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warning/0002/math0002192.html | ar5iv | text | # A note on the crystalline subrepresentation functor
## 1 Maximal Crystalline Subrepresentations
It is well known that the representations of $`G`$ with coefficients in $`_p`$-modules, the $`p`$-adic representations, have very different properties from the representations in $`𝐙_l`$-modules for $`lp`$. For example, even for a variety over $`K`$ with good reduction over $`R`$, the representation of $`G`$ on the $`p`$-adic étale cohomology is only rarely unramified.
On the other hand, $`p`$-adic Hodge theory has provided us with a fine classification of $`p`$-adic presentations together with appropriate analogies to the $`l`$-adic case. For example, the $`p`$-adic notion corresponding to an unramified $`l`$-adic representation is that of a crystalline representation. These are the representations that correspond via $`p`$-adic Hodge theory to weakly-admissible crystals (the correct $`p`$-adic analogue of local systems), whereas representations that are genuinely unramified correspond to the much smaller subcategory consisting of crystals of slope zero (see, for example, ).
We wish to continue this analogy by presenting a new class of cohomology theories associated to $`p`$-adic representations of Galois groups of local fields. The definition is very natural and elementary, and is likely to be well-known to experts. However, a specific application motivated us to commit at least a short exposition to paper:
Let $`A`$ be an abelian variety over $`K`$ and let $`𝐀`$ be its Neron model over $`R`$. Let $`A_0`$ be the special fiber of $`𝐀`$ and $`A_0^0`$ the connected component of the identity in $`A_0`$. Finally let $`\mathrm{\Gamma }=A_0(\overline{k})/A_0^0(\overline{k})`$ be the geometric points of the group of connected components of $`A_0`$. Grothendieck points out the following formula expressing the $`l`$-primary part of $`\mathrm{\Gamma }`$ in terms of Galois cohomology:
$$\mathrm{\Gamma }(l)=H^1(I,T_l(A))_{\text{tor}}$$
where $`T_l`$ refers to the $`l`$-adic Tate module and the subscript denotes the torsion subgroup. The motivating problem is that of expressing the $`p`$ part of $`\mathrm{\Gamma }`$ in an analogous ‘cohomological’ manner involving only the generic fiber.
The formula is definitely false in general if we simply substitute $`p`$ for $`l`$. An easy argument using Kummer theory shows that when $`A`$ is semi-stable over an absolutely unramified base, we actually have an injection
$$H^1(I,T_p(A))_{\text{tor}}\mathrm{\Gamma }(p)$$
which is non-surjective in general. For example, we can consider the case of an elliptic curve with split semi-stable reduction and order of discriminant $`p`$. It is an easy exercise to check that in that case, the map is surjective iff the elliptic curve has an unramified point of order $`p`$ which occurs exactly when its Tate parameter is a $`p`$-power in $`K^u`$. In short, the torsion in the Galois cohomology of $`I`$ is not big enough to capture the $`p`$-part of the component group. But notice that the Galois cohomology $`H^1(I,)`$ is just the first (right-)derived functor of the functor
$$()()^I$$
which we view as assigning to a representation its maximal unramified subrepresentation. This is an example of a ‘subrepresentation functor’ or a ‘subobject’ functor, which can occur in a wide variety of contexts whenever one has suitable subcategories of categories. On the other hand, we have already remarked that the unramified objects comprise a sub-category too small for geometric applications related to $`p`$-adic representations. This motivates us to define the crystalline subrepresentation functor
Crys
from the category of $`_p`$-representations of $`G`$ to itself. Given a $`_p`$ representation $`V`$ of $`G`$, $`\text{ Crys}(V)`$ is the maximal crystalline subrepresentation of $`V`$, where crystalline is defined in the usual way for finite-dimensional representations and in general, we say $`V`$ is crystalline if it is a direct limit of finite-dimensional crystalline subrepresentations. Equivalently, we could say $`V`$ is crystalline iff any finite dimensional subrepresentation is crystalline. This equivalence follows from the fact that the category of finite-dimensional crystalline representations is closed under sub-objects. The fact that it’s also closed under quotient objects implies that there is a well-defined notion of a ‘maximal’ crystalline subrepresentation. The functor Crys is the natural $`p`$-adic analogue of the ‘invariants under inertia’ functor on $`l`$-adic representations from the point of view of sub-representation functors. Consequently, the derived functors of Crys are natural analogues of Galois cohomology with respect to $`I`$.
To see that these notions are well-defined, we must check two things:
(1) Crys is indeed a functor: This follows from the fact that a quotient of a crystalline representation is also crystalline, so that under a map $`VW`$ of representations, the crystalline part must land in the crystalline part.
(2) Crys is left exact: The key point is that if $`UV`$ is a subrepresentation, then $`\text{ Crys}(U)=U\text{ Crys}(V)`$. The inclusion in the two directions follows from the maximality involved in the definition and the sub-object property mentioned earlier.
One could equally easily define the various ‘truncated’ functors $`\text{ Crys}_{[a,b]}`$ which associates to a representation the maximal subrepresentation with Hodge-Tate weights in the interval $`[a,b]`$. We will concentrate mostly on the functors $`\text{ Crys}_{[0,h]}`$ which we will abbreviate as $`\text{ Crys}_h`$. It will be convenient to use the term $`h`$-crystalline representations for the objects in the image of this functor. It is interesting to note that $`\text{ Crys}_0`$ is nothing but the old inertia-invariants functor, so that the sequence of functors $`\text{ Crys}_0,\text{ Crys}_1,\mathrm{}`$ and their derived functors provide natural prolongations of Galois cohomology. We see also that Crys is a bit more than just an ‘analogue’ of the inertia invariants functor. Rather, the existence of these prolongations reflect the richer structure that $`p`$-adic representations tend to have compared to their $`l`$-adic counterparts. We propose that these derived functors are natural invariants of $`p`$-adic representations (at least as natural as Galois cohomology) and should be studied seriously. One reason for thinking so stems from the application mentioned above. For this, we need to define these functors also for integral representations. Unfortunately, here the existing techniques for making the correct definitions are rather incomplete, and we can define only the truncated functors $`\text{ Crys}_i`$ for $`ip2`$. (One can actually prolong it slightly to $`i=(p1)^{}`$ in an appropriate sense, but we shall keep to the smaller truncation for simplicity of exposition.)
We also need to assume that $`K`$ is absolutely unramified so that $`K=K_0`$ and $`R=W`$. The foundational material we need is contained in the seminal paper of Fontaine and Laffaille , but the reader can find a nice summary in .
Let $`h`$ be a natural number $`p2`$. One first defines finite crystalline representations of height $`h`$, or the finite $`h`$-crystalline representations , to be the essential image of the category $`\mathrm{𝐌𝐅}_{R,\text{tor}}^h`$ (the finite-length filtered $`\varphi `$-modules of height $`h`$) under the fully-faithful functor
$$MV_{\text{crys}}^{}(M):=\text{Hom}_{\mathrm{𝐌𝐅}_R}(M,A_{\text{crys},\mathrm{}})$$
Next, one defines a finite-type $`_p`$-module $`L`$ with $`G`$-action to be $`h`$crystalline if $`L=\underset{}{\mathrm{lim}}L_i`$ where the $`L_i`$ are finite-length $`h`$-crystalline representations. The fact that $`h`$crystalline representations are closed under sub- and quotient objects follows from the corresponding property for $`\mathrm{𝐌𝐅}_{R,\text{tor}}^h`$. In particular, this implies that a finite-type $`_p`$ representation $`L`$ is $`h`$-crystalline iff $`L/p^nL`$ is $`h`$crystalline for all $`n`$ (which is the definition of ), and when $`L`$ is free, iff
$$L=V_{\text{crys}}^{}(M):=\text{Hom}_{\mathrm{𝐌𝐅}_R}(M,A_{\text{crys}})$$
for an object $`M`$ of $`\mathrm{𝐌𝐅}_R^h`$ (the finitely generated free filtered $`\varphi `$-modules of height $`h`$) ( 2.2.2).
Now for an arbitrary $`_p[G]`$-module $`V`$, we define it to be crystalline if $`V=\underset{}{\mathrm{lim}}V_i`$ where the $`V_i`$ are subrepresentations of finite-type.
We need to check that this definition is consistent with the existing one for $`_p`$-representations. Since we defined it for the infinite-dimensional case using limits from finite dimensions, we need only check it for finite-dimensional representations. So assume that $`V`$ is $`h`$crystalline in the old sense. Then $`V`$=$`\text{Hom}_{MF_K}(\mathrm{\Delta },B_{\text{crys}})`$ for some $`\mathrm{\Delta }`$ in $`\mathrm{𝐌𝐅}_K^h`$ (remarque 8.5 and 8.13 (c)). Since $`\mathrm{\Delta }`$ is $`B_{\text{crys}}`$-admissible, in particular, weakly admissible, one can find a strongly divisible lattice $`M\mathrm{\Delta }`$ which is an object of $`\mathrm{𝐌𝐅}_R^h`$. So we get $`V=L_p`$ where $`L=\text{Hom}_{\mathrm{𝐌𝐅}_R}(M,A_{\text{crys}})`$. Now, $`L`$ is $`h`$crystalline and $`V=\underset{}{\mathrm{lim}}L[1/p^n]`$ while $`L[1/p^n]L`$ (via multiplication by $`p^n`$) is $`h`$-crystalline. So $`V`$ is $`h`$-crystalline in the new sense. In the other direction, assume $`V=\underset{}{\mathrm{lim}}L_i`$ for $`h`$crystalline submodules $`L_i`$ of finite-type. Then some $`L=L_i`$ is a lattice and $`V=L_p`$. But $`L=\text{Hom}_{\mathrm{𝐌𝐅}_R}(M,A_{\text{crys}})`$ for some free $`R`$-module $`M`$ in $`\mathrm{𝐌𝐅}_R^h`$ and $`M`$ is then a strongly divisible lattice in $`\mathrm{\Delta }:=MK`$ according to the terminology of definition 7.7, and therefore, $`MK`$ is weakly admissible. Thus, by the main theorem of , $`MK`$ is $`B_{\text{crys}}`$admissible and $`V=\text{Hom}_{\mathrm{𝐌𝐅}_K}(MK,B_{\text{crys}})`$ is crystalline.
Thereby, we can define $`\text{ Crys}_h`$, the maximal $`h`$crystalline subrepresentation functor for $`hp2`$ compatibly on all $`_p[G]`$ modules. An easy consequence of the definitions is that if $`L`$ is a finitely generated free $`_p`$ representation, then $`\text{ Crys}_h(L)=\underset{}{\mathrm{lim}}\text{ Crys}_h(L/p^nL)`$. It should be emphasized that we also have Crys and all the other $`\text{ Crys}_{[a,b]}`$’s if we stick to rational representations. By the key property that
$$L_1L_2\text{ Crys}_h(L_1)=L_1\text{ Crys}_h(L_2)$$
we again have left exactness and therefore, all the right-derived functors.
A systematic study of these functors will be presented in the forthcoming Ph.D. thesis of the second author.
## 2 The $`p`$-complement to Grothendieck’s formula
In this section, we will continue to assume that $`K`$ is absolutely unramified, and furthermore, that $`p>2`$.
We will be using one more functor $`FF`$ which associates to a $`p`$-adic representation its maximal ‘finite and flat’ part. Of course, one needs to define finite flat $`p`$-adic representations in a general setting. For finite $`_p`$ representations, finite flat means the usual thing: a finite representation is finite flat if it’s isomorphic to the $`\overline{K}`$ points of a finite flat commutative group scheme over $`R`$. A finite-type $`_p`$-representation is defined to be finite flat if it is the inverse limit of finite finite flat representations (the double adjective seems unfortunately unavoidable). Finally, an arbitrary $`_p`$ representation for $`G`$ is said to be finite flat if it is the direct limit of finite flat representations of finite type.
For finite representations, the property of being finite flat is closed under passing to sub-objects and quotient objects (using Zariski closure and construction of good quotient schemes), so the same is true for any $`_p`$ representation. Thus it makes sense to speak of the maximal finite flat subrepresentation of any representation, and the associated functor $`FF`$ is left exact. Thus, we can consider its derived functors. In fact, by Fontaine-Laffaille’s description of finite flat group schemes (, section 9) $`FF`$ is nothing but $`\text{ Crys}_1`$. Notice, however, that $`FF`$ is defined over an arbitrary local field, not necessarily absolutely unramified.
We will also need the trivial observation that if
$$0M_1M_2M_30$$
is an exact sequence in $`\mathrm{𝐌𝐅}_R`$, $`M_1`$ and $`M_3`$ are in $`\mathrm{𝐌𝐅}_R^h`$, and $`M_2`$ is in $`\mathrm{𝐌𝐅}_R^h^{}`$ for some $`h^{}`$, then in fact, $`M_2`$ is in $`\mathrm{𝐌𝐅}_R^h`$. This follows by noting that the morphisms are strict so that any $`F^iM_2`$ for $`i>h`$ would have to be zero when intersected with $`M_1`$ and mapped to $`M_3`$, and hence, must be zero.
Thus, we have an obvious corresponding statement for $`h`$ and $`h^{}`$-crystalline representations.
We now return to the problem of expressing the $`p`$-part of $`\mathrm{\Gamma }`$ in an analogous manner to Grothendieck’s formula for $`lp`$
$$\mathrm{\Gamma }(l)H^1(I,T_lA)_{\text{tor}}.$$
(1)
To derive the above formula, Grothendieck shows that
$$\mathrm{\Gamma }[l^n]A[l^n]^f/(A^0[l^n])^f$$
(2)
where $`\mathrm{\Gamma }[l^n]`$ (resp. $`A[l^n]`$) denotes the kernel of multiplication by $`l^n`$ on $`\mathrm{\Gamma }`$ (resp. $`A(\overline{K})`$), and the superscript $`f`$ denotes the “finite part” (denoted the “fixed part” by Grothendieck in , section 2.2.3), i.e., the points that extend to a map from $`\text{Spec}(\overline{R})`$ to $`𝐀`$, or equivalently, the $`\overline{K}`$ points of the maximal finite flat subgroup scheme of $`𝐀[l^n]`$. Similarly, $`(A^0[l^n])^f`$ denotes the $`\overline{K}`$ points of the maximal finite flat subgroup scheme of $`𝐀^0[l^n]`$ which can also be thought of as the points of $`A[l^n]^f`$ which reduce mod $`p`$ to a point in $`A_0^0`$, the connected component of the identity in the special fiber. The key point then is that the finite part coincides with the inertia invariants of $`A[l^n]`$ (resp. $`A^0[l^n]`$) (Proposition 2.2.5) and the formula (1) follows easily.
In the case of $`l=p`$, (2) still holds (provided one assumes semi-stability), but it is no longer the case that the fixed part and inertia invariants coincide. However, we will show below that (for $`lp`$) the finite part coincides with the *maximal $`h`$-crystalline part* for any $`1hp2`$ (recall that $`p>2`$). This will allow us to derive, in a completely analogous manner to Grothendieck, the following:
###### Theorem 1
Let $`A`$ be an abelian variety over the absolutely unramified local field $`K`$ with semi-stable reduction and $`1hp2`$. Then
$$\mathrm{\Gamma }(p)R^1\text{ Crys}_h(T_pA)_{\text{tor}}.$$
Proof.
We will first show that $`\text{ Crys}_h(A[p^n])=(A[p^n])^f`$. For this, we note that the fixed part of $`A[p^n]`$ is none other than $`FF(A[p^n])`$. That is, the fixed part is finite-flat by definition, giving us one inclusion
$$(A[p^n])^fFF(A[p^n]).$$
Now let $`𝒱`$ denote the finite-flat group scheme extending $`FF(A[p^n])`$, so that if $`V`$ is the generic fiber of $`𝒱`$, we have
$$V(\overline{K})FF(A[p^n])$$
as $`G`$-modules. From the inclusion $`FF(A[p^n])A[p^n]`$, we have a map
$$VA.$$
We need to show that this map extends to a map $`𝒱𝐀`$, thereby showing that the finite part is actually “finite inside $`𝐀`$.” However, restricting to the connected component $`V^0`$ of $`V`$, we find that the image must actually land in the finite part of $`A`$. This follows because $`A[p^n]/A[p^n]^f`$ is unramified (, Proposition 5.6). By results of Raynaud , this extends to a map $`𝒱^0𝐀^f`$. Hence by Lemma 5.9.2 of , we get a unique map $`𝒱𝐀`$ extending the two previous maps, and giving us the opposite inclusion. (This is essentially the same argument as in , Lemma 6.2.)
We saw above that $`FF=\text{ Crys}_1`$, as functors. We will now show that one can replace $`\text{ Crys}_1`$ by any of the $`\text{ Crys}_h`$’s in our setting. In fact, we will see from the proof that if any general Crys functor were defined for finite representations, then that could be used as well.
We certainly have an inclusion
$$FF(A[p^n])\text{ Crys}_h(A[p^n])$$
which induces an inclusion of $`\text{ Crys}_h(A[p^n])/FF(A[p^n])`$ into the unramified $`G`$-module $`A[p^n]/FF(A[p^n])`$. Thus, $`\text{ Crys}_h(A[p^n])/FF(A[p^n])`$ is unramified as well, and hence finite-flat as a representation. Therefore $`\text{ Crys}_h(A[p^n])`$ sits in the middle of a short exact sequence whose outer terms are both crystalline of height one (actually, the last is of height 0):
$$0FF(A[p^n])\text{ Crys}_h(A[p^n])\text{ Crys}_h(A[p^n])/FF(A[p^n])0$$
By the observation made earlier, we see that $`\text{ Crys}_h(A[p^n])`$ is itself crystalline of height one, and thus equal to $`\text{ Crys}_1(A[p^n])=FF(A[p^n])=(A[p^n])^f`$. As $`A^0[p^n]/FF(A^0[p^n])`$ is contained in $`A[p^n]/FF(A[p^n])`$, it is also unramified and an entirely similar argument gives $`\text{ Crys}_h(A^0[p^n])=(A^0[p^n])^f`$ as well.
The isomorphism (2) thus becomes
$$\mathrm{\Gamma }[p^n]\text{ Crys}_h(A[p^n])/\text{ Crys}_h(A^0[p^n])$$
However,
Claim:
$$\text{ Crys}_h(A[p^n])/\text{ Crys}_h(A^0[p^n])\text{ Crys}_h(T_pA_p/p^n_p)/\text{ Crys}_h(T_pA)_p/p^n_p.$$
Proof. The equality between the ‘numerators’ is obvious, so we need to see that $`\text{ Crys}_h(T_pA)_p/p^n_p`$ is equal to $`\text{ Crys}_h(A^0[p^n])`$. But
$$\text{ Crys}_h(T_pA)\underset{}{\mathrm{lim}}(\text{ Crys}_h(A[p^n])=\underset{}{\mathrm{lim}}(A[p^n]^f)$$
Hence $`\text{ Crys}_h(T_pA)(T_pA)^f`$. Thus,
$$\text{ Crys}_h(T_pA)_p/p^n_p(T_pA)^f_p/p^n_p(A^0[p^n])^f\text{ Crys}_h(A^0[p^n]).$$
(The key point is the second isomorphism, as explained in . That is, if you take the finite part of the Tate module and then reduce mod $`p^n`$, then you end up in $`(A^0[p^n])^f`$ because the multiplication by $`p`$ map is finite and surjective only on $`𝐀^0`$.)
Applying direct limits, we get the formula
$$\mathrm{\Gamma }(p)\text{ Crys}_h(T_pA_p/_p)/\text{ Crys}_h(T_pA)_p/_p.$$
This plays a role analogous to Grothendieck’s formula (Proposition 11.2).
Following Grothendieck, we next apply $`\text{ Crys}_h`$ to the short exact sequence
$$\begin{array}{ccccccccc}0& & T_pA& & T_pA_p& & T_pA_p/_p& & 0\end{array}$$
and obtain the long exact sequence
$$\begin{array}{ccccccccc}0& & \text{ Crys}_h(T_pA)& & \text{ Crys}_h(T_pA_p)& & \text{ Crys}_h(T_pA_p/_p)& & \\ & & R^1\text{ Crys}_h(T_pA)& & R^1\text{ Crys}_h(T_pA_p)& & R^1\text{ Crys}_h(T_pA_p/_p)& & \mathrm{}\end{array}$$
Note that
$$\text{ Crys}_h(T_pA_p)=(\text{ Crys}_h(T_pA_p)T_pA)_p=\text{ Crys}_h(T_pA)_p$$
Thus the kernel of the map of
$$R^1\text{ Crys}_h(T_pA)R^1\text{ Crys}_h(T_pA_p)$$
is $`\text{ Crys}_h(T_pA_p/_p)/\text{ Crys}_h(T_pA)_p/_p`$, i.e. $`\mathrm{\Gamma }(p)`$. Since $`\mathrm{\Gamma }(p)`$ is torsion and $`R^1\text{ Crys}_h(T_pA_p)`$ torsion-free, we do indeed find that
$$\mathrm{\Gamma }(p)R^1\text{ Crys}_h(T_pA)_{\text{tor}}.$$
Remark. From the proof, it is clear that one could have just used the functor $`FF`$ for the theorem in which case one could extend the theorem to the case of $`ep2`$ by eliminating the Fontaine-Laffaille theory. However, this would have made the analogy to the $`l`$-case less natural, since a crystalline resepresentation is clearly the correct general notion which sets the formula into a broad context. In particular, the definition of $`FF`$ on $`_p`$ representations is rather artificial compared to Crys. It would of course have been nicer to replace $`\text{ Crys}_h`$ by a general Crys even for the integral representations.
Acknowledgements: We are grateful to Wiesia Niziol for pointing out the definitions in and to Ken Ribet for directing us to lemma 6.2 of . We are especially grateful to Kirti Joshi for innumerable conversations on p-adic Hodge theory.
M.K. was supported in part by NSF grant DMS-9701489
DEPARTMENT OF MATHEMATICS, UNIVERSITY OF ARIZONA,TUCSON, AZ 85721, U.S.A. E-MAIL: kim@math.arizona.edu susan@math.arizona.edu |
warning/0002/hep-th0002068.html | ar5iv | text | # Extended Conformal Symmetry
## 1 Introduction
Scale invariance provides an successful class of physical (string) models in 2-dim. One element of the success in 2-dim is the infinite character of the conformal Lie algebra in that dimension. However, as we show here, the standard formulation of scale invariance introduced by Weyl is incomplete. By studying scale covariance \- the grading of fields by conformal weight implicit in Weyl or conformal geometries - we find a noncentral, conformal extension of any Lie group containing dilatations. We find that the infinite Heisenberg and Virasoro algebras are natural adjuncts of scale covariance in any dimension. When the extended groups are gauged the resulting spaces possess a natural commutator product and a scale-invariant inner product.
Initially motivated by the assignment of weights to classical fields, the use of scale covariance is transformed to a necessity by the well-known result that in quantum field theory, scale-invariant theories spontaneously develop scales. Thus, we build a conformally covariant theory based on conformal or other scaling symmetries. We start with the observation that it is not only the units of length that are arbitrary. In addition, at the outset we can assign an arbitrary conformal weight to the coordinates – then the scale invariance of the action fixes the weights of other fields. When this freedom is described as a symmetry we are led to introduce a set of operators that change the assigned weight. The full collection of such operators forms a noncentral, discrete extension of the conformal group which satisfies any of the standard definitions of conformal transformations (i.e., preserve angles, preserve ratios of infinitesimal displacements, preserve light cones , ).
Extended scaling groups and their gauge theories have a range of novel properties. First, they fill some gaps in the previous formulation. In the past, the covariance of fields has been managed by inserting the conformal weight by hand into gauge transformations and the covariant derivative. But in extended conformal gauge theories, the geometric structure automatically recognizes the weights of fields. This provides a check on our development, indicating the correctness of the extension.
What is more interesting than completing Weyl’s picture, however, is the potential usefulness of the new structures for building new field theory models. Extended scaling groups are Lie groups with infinitely many connected components. Tangent vectors to the group manifold are both homothetic tensors and elements of an infinite dimensional vector space, which obey a commutator algebra containing the infinite Heisenberg algebra, the Virasoro algebra, homothetic algebras and Kac-Moody algebras as sub-algebras. These vectors also have an invariant scalar product which is positive definite on the conformally self-dual subspace.
We make use of a recent advance in conformal gauge theory in $`n`$-dim to display some of these properties of the extended symmetry. In the past, conformal gaugings (-) treated the special conformal transformations as additional symmetries of the final gauged geometry. But biconformal gauging requires the special conformal transformations together with the translations to span the base manifold, giving a $`2n`$-dim space with $`n`$-dim Weyl symmetry group . Because the biconformal gauging has a dimensionless volume form, it is possible to write a scale-invariant action linear in the biconformal curvatures. This action leads to two involutions of the $`2n`$-dim space which show it to be spanned by a Ricci-flat, $`n`$-dim Riemannian spacetime together with a flat $`n`$-dim Riemannian spacetime . Thus, higher-dim conformal symmetry is consistent with vacuum general relativity in a straightforward way, without requiring a quadratic action or compensating fields. Recently, these results have been extended to include certain matter sources (,). Below we illustrate the use of extended conformal symmetry by performing a biconformal gauging of the extended conformal group.
The layout of the paper is as follows. In the next section, we derive the form of extended scaling groups from the usual representation for dimensionful fields. Then, in Sec.(3), we develop a vector representation on which the extended group acts effectively. The vector representation is isomorphic to the space of tangent vectors to the extended group manifold. We show that these vectors satisfy a commutator algebra containing the infinite Heisenberg and Virasoro algebras and define a class of inner products. Finally, in Sec.(4), we indicate a few of the consequences of the structures of Sec.(3) for field theories based on the extended conformal group. Sec.(5) contains a brief summary.
## 2 Extended Scale Invariance
In this section and the next, we establish our main results. Here we define the extended conformal group, find its center and find their quotient group. The quotient group acts transitively and effectively. In the next section, we find a vector representation of the group and show that the vectors close under commutation and have an invariant, indefinite scalar product.
To build the group, we first note that nothing physical is changed if we assign an arbitrary conformal weight to the spacetime coordinates, adjusting the weights of other fields accordingly. For simplicity, we take the possible weight assignments to lie in the set $`Z`$ of integers. This much is necessary to account for our ability to step the conformal weight up or down arbitrarily many times by integrating and differentiating fields.
Now we start with the standard definition: The extended conformal group, $`𝒞_E,`$ is the group of transformations that preserves angles. We understand this definition to include transformations of conformal weight, so, in addition to the usual conformal transformations, $`𝒞_E`$ includes an infinite discrete part which formalizes the conformal grading. The development is as follows.
Suppose $`𝒢`$ is a Lie symmetry group including the $`1`$-dim dilatational subgroup, $`e^{\lambda D},`$ with generator $`D.`$ Then the generators $`\{G_A\}=\{D,G_\alpha \}`$ of the Lie algebra, $`_g,`$ of $`𝒢`$ satisfy commutation relations of the form
$`[D,G_\alpha ]`$ $`=`$ $`c_{0\alpha }^\beta G_\beta +c_{0\alpha }^0D`$ (1)
$`[G_\alpha ,G_\beta ]`$ $`=`$ $`c_{\alpha \beta }^\gamma G_\gamma +c_{\alpha \beta }^0D`$ (2)
Examples for $`𝒢`$ include the homogeneous and inhomogeneous homothetic (or Weyl) groups and the conformal group. We wish to study tensor field representations $`\varphi \mathrm{\Phi }`$ of $`𝒢`$ which in addition to their transformation properties under $`𝒢`$
$$\varphi ^{}(x^{})=e^{\lambda ^AG_A}\varphi (x)$$
(3)
are assigned geometric units from the set $`L=\{(length)^m|mZ\}.`$ This graded field representation therefore takes the product form
$$\varphi (x)(length^k)=\varphi (x)l^k\mathrm{\Phi }L$$
(4)
We seek an extension, $`𝒢_E,`$ of the original group $`𝒢,`$ that acts on the representation $`\mathrm{\Phi }L.`$ The extension should include a subgroup, $`𝒥`$, of operators $`J_\mathrm{\Sigma }:LL.`$ Thus, $`𝒥`$ is a subgroup of the automorphism group of $`L;`$ in addition we ask for $`𝒥`$ to be closed under the action of $`𝒢`$. Since we expect, for example, Lorentz transformations or translations to commute with changes of assigned conformal weight, this amounts to closure under dilatations, $`e^{\lambda D}𝒥e^{\lambda D}𝒥`$. The largest such group is easily seen (see Appendix A) to be the set $`\{J_k|J_kl^m=l^{m+k}\}`$ satisfying $`e^{\lambda D}J_ke^{\lambda D}=e^{\lambda k}J_k.`$ It follows that
$`[J_k,J_m]`$ $`=`$ $`0`$ (5)
$`[J_m,D]`$ $`=`$ $`mJ_m`$ (6)
$`[J_m,G_\alpha ]`$ $`=`$ $`0`$ (7)
$`J_kJ_m`$ $`=`$ $`J_{k+m}`$ (8)
We form a group containing both $`D`$ and $`\{J_k\}`$ by writing
$$h(\alpha ,k,\lambda )=e^\alpha J_ke^{\lambda D}$$
(9)
Then the group product is
$`h(\alpha ,k,\lambda )h(\beta ,m,\gamma )`$ $`=`$ $`e^{\alpha +\beta \lambda m}J_{k+m}e^{(\lambda +\gamma )D}`$ (10)
$`=`$ $`h(\alpha +\beta \lambda m,k+m,\lambda +\gamma )`$ (11)
The identity element is $`h(0,0,0)=J_0,`$ and the inverse to $`h(\alpha ,k,\lambda )`$ is $`h(\alpha +\lambda k,k,\lambda ).`$ Notice that the positive real factor $`e^\alpha `$ is necessary in order to accomodate the factor $`e^{\lambda m}`$ that arises in commuting $`e^{\lambda D}`$ and $`J_m`$ into standard form after taking a product.
Next, we extend the dilatation factor, $`e^{\lambda D},`$ to include the rest of the group $`𝒢`$. As noted above, the remaining generators $`G_\alpha `$ may be expected to commute with $`J_k`$. Therefore a general element $`g𝒢_E`$ takes the form
$$g(\alpha ,k,\lambda ,\lambda ^\alpha )=e^\alpha J_ke^{\lambda D+\lambda ^\alpha G_\alpha }$$
(12)
where $`\lambda D+\lambda ^\alpha G_\alpha `$ is a general element of $`_g`$. We easily check that $`𝒢_E`$ is a Lie group. In particular, the inverse of $`g`$ is given by
$$g^1=e^\alpha e^{\lambda D\lambda ^\alpha G_\alpha }J_k=e^{\alpha \lambda k}J_ke^{\lambda D\lambda ^\alpha G_\alpha }𝒢_E$$
(13)
The Lie algebra, $`_{g_E},`$ of $`𝒢_E`$ is $`_g`$ extended by the identity to include the necessary factor $`e^\alpha :_{g_E}=_g\mathrm{𝟏}.`$ The $`J_k`$ operators which change conformal weights give the group an infinite number of distinct connected components. Clearly, each of these components is a manifold of dimension ($`dim𝒢+1)`$ which is in $`11`$ correspondence with $`𝒢R^+`$. Thus, as a manifold, $`𝒢_E`$ is homeomorphic to the direct product, $`𝒢R^+𝐙.`$
In order to construct a gauge field theory based on $`𝒢_E,`$ we study the adjoint action of $`𝒢_E`$ on the manifold $`=𝒢_E,`$ seeking the maximal effective subgroup. This will be $`𝒢_E`$ modulo its center, where the center of $`𝒢_E`$ is the set
$$K=\{g|gpg^1=p,p𝒢_E\}=\{e^\alpha \}$$
(14)
Notice that the elements $`J_k`$ do not lie in the center because for a general $`p=J_me^{\alpha +\lambda D+\lambda ^\alpha G_\alpha }`$ we have
$$J_kpJ_k=J_k(J_me^{\alpha +\lambda D+\lambda ^\alpha G_\alpha })J_k=e^{\lambda k}pp$$
(15)
However, the projective subgroup $`𝒢_E/K`$ is isomorphic to the direct product $`𝒢𝒥`$ because $`J_k`$ and $`g_c=`$ $`e^{\lambda D+\lambda ^\alpha G_\alpha }`$ now commute in $`𝒢_E/K`$:
$$J_kg_c=J_ke^{\lambda D+\lambda ^\alpha G_\alpha }=e^{k\lambda }g_cJ_kg_cJ_k$$
(16)
The quotient $`𝒢_E/K`$ is the maximal effective subgroup. The quotient introduces central charges into both the Lie algebra of $`𝒢_E/K,`$ and into the commutator algebra of the extended representation and tangent space. The central charges for $`𝒢_E/K`$ depend only on the original Lie algebra of $`𝒢.`$ Those of the tangent space are discussed below.
We note also that $`𝒢_E/K`$ is transitive, since any two elements $`p=J_kg_{c1}`$ and $`q=J_mg_{c2}`$ are connected by the element $`r=J_{mk}g_{c2}g_{c1}^1:`$
$$rp=J_{mk}g_{c2}g_{c1}^1J_kg_{c1}J_mg_{c2}=q$$
(17)
We now find a representation for $`𝒢_E,`$ which turns out to be isomorphic to the tangent space, $`T𝒢_E`$ of $`𝒢_E.`$
## 3 The extended representation space
In this section, we continue to study properties of $`𝒢_E,`$ finding a faithful vector representation for $`𝒢_E`$ and developing its properties. Normally, the Lie algebra of a Lie group provides an adequate vector representation for a Lie group. However, when we wish to represent discrete symmetries this representation is not faithful and we must look at infinitesimal transformations in each connected component, rather than just a neighborhood of the identity. To illustrate this point, we first consider the trivial example of $`O(3),`$ including parity. Then we apply the technique to $`𝒢_E.`$
For $`O(3),`$ a general group element may be written in the form $`g=P_\alpha e^{\frac{1}{2}\lambda ^iM_i}`$ where $`\alpha \{+,\},`$ $`P_+=\mathrm{𝟏},P_{}=\mathrm{𝟏}`$ and the $`M_i`$ generate rotations. The $`P_\alpha `$ have product $`P_\alpha P_\beta =P_{\alpha \times \beta }.`$ An element of the Lie algebra is simply $`v=v^iM_i,`$ which is insufficient to represent the action of parity. Instead, we expand about each connected component to find six generators, $`P_+M_i`$ and $`P_{}M_i,`$ and a representation of the form
$$v=v_+^i(P_+M_i)+v_{}^i(P_{}M_i)$$
(18)
The enlarged vector space is sufficient to span such common indefinite parity combinations as
$$\stackrel{}{y}=\stackrel{}{w}+\stackrel{}{u}\times \stackrel{}{v}=(\stackrel{}{u}\times \stackrel{}{v})^k(P_+M_k)+w^k(P_{}M_k).$$
(19)
Notice that the basis vectors form a commutator algebra under the usual Leibnitz rule for the commutator of a product. For example, the commutator of two vectors in the $`P_{}`$ sector appropriately lie in the $`P_+`$ sector since
$$[v,w]=v^iw^j[P_{}M_i,P_{}M_j]=(\stackrel{}{v}\times \stackrel{}{w})^kP_+M_k$$
(20)
Similarly, to find a faithful representation of $`𝒢_E`$ we consider the expansion of $`𝒢_E`$ about each $`J_k.`$ We find
$$e^\alpha J_ke^{\lambda D+\lambda ^\alpha G_\alpha }J_k(1+\alpha 1+\lambda D+\lambda ^\alpha G_\alpha )$$
(21)
so that the representation has the basis
$$A=(J_k\mathrm{𝟏},J_kD,J_kG_\alpha )(J_k,L_k,G_\alpha ^k)$$
(22)
where we have defined
$`L_k`$ $``$ $`J_kD`$ (23)
$`G_\alpha ^k`$ $``$ $`J_kG_\alpha `$ (24)
It is straightforward to show that vectors of the form
$$v=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\left(v^kJ_k+v^{0k}L_k+v^{\alpha k}G_\alpha ^k\right)V$$
(25)
give a faithful representation of $`𝒢_E.`$ In fact, this representation is isomorphic to the tangent space to $`𝒢_E.`$ The components of $`v`$ have conformal weights $`(k,k,k+w_\alpha )`$ where $`w_\alpha `$ is the weight of $`G_\alpha .`$
Since $`[J_m,J_n]=0`$ and because of the form of the product $`J_kJ_m`$ given in eq.(8), this space also admits a commutator product between vectors, again defined using $`[AB,C]=A[B,C]+[A,C]B`$. This algebra is generally not a Lie algebra because the Jacobi identities fail. Instead, we have Jacobi relations, given in Appendix B. All of the non-Jacobi terms are proportional to $`c_{\alpha \beta }^0`$ or $`c_{0\alpha }^0,`$ so for the case of homogeneous or inhomogeneous homothetic algebras, the usual Jacobi identities do hold. Recalling that the projective quotient also introduces central charges into $`V`$, the algebra of $`V`$ becomes
$`[G_\alpha ^k,G_\beta ^m]`$ $`=`$ $`c_{\alpha \beta }^\gamma G_\gamma ^{k+m}+c_{\alpha \beta }^0L_{k+m}+c_{\alpha \beta }\delta _{k+m}^0`$ (26)
$`[L_k,G_\alpha ^m]`$ $`=`$ $`(c_{0\alpha }^\beta m\delta _\alpha ^\beta )G_\beta ^{k+m}+c_{0\alpha }^0L_{k+m}`$ (27)
$`[L_{k,}L_m]`$ $`=`$ $`(km)L_{k+m}+ak(k^21)\delta _{k+m}^0`$ (28)
$`[J_m,L_k]`$ $`=`$ $`mJ_{k+m}+b_k\delta _{k+m}^0`$ (29)
$`[J_k,J_m]`$ $`=`$ $`ck\delta _{k+m}^0`$ (30)
The algebra has several readily identifiable properties.
1. The subalgebra generated by $`\{G_\alpha ^0,L_0,J_0=\mathrm{𝟏}\}`$ is the original extended Lie algebra of $`𝒢_E,`$ while $`\{G_\alpha ^0,L_0\}`$ is the Lie algebra of $`𝒢`$.
2. The subalgebra generated by $`\{J_k\}`$ is the infinite Heisenberg algebra.
3. The subalgebra generated by $`\{L_k\}`$ is the Virasoro algebra.
4. For any proper Lie subalgebra of $`_g,`$ with basis $`\{H_\beta \}\{G_\alpha \},`$ the set $`\{H_\beta ^k\}`$ is a basis for the associated Kac-Moody algebra. For example, when $`𝒢`$ is the conformal group, eq.(26) includes the Poincaré-Kac-Moody algebra. Notice that the commutator algebra does not contain the Kac-Moody algebra of $`𝒢`$ because of the nontrivial commutator for the $`L_k.`$ This commutator is nontrivial precisely because the extension is noncentral, i.e., the dilatation generator $`D`$ measures the weight of $`J_k`$ the same way it measures the weight of $`G_A.`$
We may also define a class of indefinite, weight-$`m`$ scalar products for $`V`$ whenever the original group has a nontrivial Killing metric, $`K_{AB}=c_{EA}^Fc_{FB}^E.`$ Using the adjoint representation the Lie algebra of $`𝒢`$ we have
$$K_{AB}=tr(G_AG_B)$$
(31)
so that if we define
$$Tr(J_kG_A)=\delta _k^0tr(G_A)$$
(32)
we have
$`v,w_m`$ $``$ $`Tr(vJ_mw)`$ (33)
$`=`$ $`{\displaystyle \underset{k}{}}v^{A,k}w^{B,mk}K_{AB}=w,v_m`$ (34)
where we have used $`tr(G_A)=0.`$ The result has weight $`m`$ because we require the weight of $`v=v^mJ_m`$ to be zero. Of greatest interest is the scale-invariant case, $`m=0,`$
$$v,wv,w_0=Tr(vw)=\underset{k}{}v^{Ak}w^{B,k}K_{AB}$$
(35)
We can find a subspace on which he $`0`$-weight inner product defines a norm. Let the conformal dual of a vector be defined as
$$\overline{v}v^{kA}J_kG_A$$
(36)
and define $`v`$ to be self-dual if there exists a gauge in which $`v=\overline{v}.`$ Then a gauge-invariant norm is given on the space of self-dual vectors by
$$v^{2}v,v=\underset{k}{}v^{kA}v^{kB}K_{AB}$$
(37)
In general, $`v^{2}`$ shares the signature of $`K_{AB}.`$ For self-dual vectors of the form $`v_v=v^kJ_k`$ (vertical on the bundle defined below) the norm is positive definite.
As an example of the use of the extended conformal group, we now consider its biconformal gauging.
## 4 Gauging the extended conformal group
The connected component of the extended conformal group only differs from the conformal group through the presence of the positive real factor, $`e^\alpha ,`$ and this part is factored out to produce an effective group action. Therefore, there is no difference between the local structure of extended biconformal gauge theory and the biconformal gauging described in (-). For this reason, we give only a brief summary of the biconformal space, then move directly to some properties of graded tensor fields.
We first consider the action of $`𝒢_E/K`$ on $`=𝒢_E/K.`$ As just noted, the local structure of $``$ is described by the original Lie algebra of $`𝒢.`$ This means that even though $``$ has multiple connected components, the connection is still a $`𝒢`$-valued $`1`$-form on $`.`$ Each of the connected components therefore shares the same curvature, and $``$ is homeomorphic to a direct product, $`J_0.`$ The $`𝒢`$-valued connection is sufficient to provide a unique $`𝒢`$ mapping along any given curve between any two points on the same connected component, while the discrete operators $`J_k`$ map uniquely from component to component. Thus, the direct product allows us to define a $`𝒢_E/K`$ action on $`.`$ Moreover, the tangent bundle to $``$ includes a copy of the tangent space associated with each connected component, giving an infinite-dimensional irreducible vector representation of $`𝒢_E/K`$. The direct product structure of the base manifold, $`J_0,`$ means the manifold itself may be treated as a trivial bundle with projection $`\pi :J_k\mathrm{𝟏}.`$ The tangent space $`T𝒢_E`$ may be divided into horizontal and vertical vector spaces using this projection.
Our notation follows that of refs and , and is based on $`O(n,2)`$. The fibre bundle is given by the quotient $`𝒞/𝒲`$, with the connection $`1`$-form
$$\omega =\frac{1}{2}\omega _b^aM_a^b+\omega _0^0D$$
(38)
where the $`M_b^a`$ generate Lorentz transformations. The biconformal base manifold is spanned by the $`2n`$ $`1`$-forms $`(\omega _0^a,\omega _a^0)`$. Because these basis forms have opposite scaling weights, biconformal geometry has a scale invariant volume form, and allows us to write a scale invariant action linear in the curvature tensors without the use of compensating fields. The resulting field equations, subject to a constraint of minimal torsion, lead to a foliation of the $`2n`$-dim space by $`n`$-dim Riemannian spacetimes satisfying the vacuum Einstein equation. The theory therefore makes close contact with general relativity. The structure equations, curvatures and gravitational field equations are reported in detail in .
We now consider tensor fields on such a biconformal geometry. The effect of such matter fields on the results of are under separate investigation , so we will limit our focus to a few basic properties of graded matter fields in flat space. For tangent vectors we have
$$v=v^mJ_m+v_m^aP_a^m+v_a^mK_m^a$$
(39)
where the generators $`P_a`$ and $`K^a`$ are for translations and co-translations, respectively. The covariant derivative is given by
$`𝐃v`$ $`=`$ $`𝐝v+[\omega ,v]`$ (40)
$`=`$ $`(𝐝v^m+m\omega _0^0v^m)J_m`$ (43)
$`+(𝐝v_m^av_m^c\omega _c^a+(m+1)\omega _0^0v_m^a)P_a^m`$
$`+(𝐝v_a^m+\omega _d^cv_c^m+(m1)\omega _0^0v_a^m)K_m^a`$
which shows that the scale-covariant derivative, with appropriate weights, emerges correctly. This result confirms our claim at the start, that the extended conformal group completes Weyl’s description of scale invariance. We can now differentiate arbitrary weight fields correctly without inserting the weights by hand.
In addition to providing a consistent formalism, the extended group has interesting field theoretic properties. Consider the dynamics of a general weight vector, $`v=v^mJ_m.`$ We easily write a massive, scale-invariant action for $`v,`$ using the invariant inner product:
$$\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(𝐃v,𝐃v+\mathrm{m}^2v,v)𝚽$$
(44)
where $`𝐃v,𝐃v=K^{AB}D_Av^nD_Bv^n`$ and $`𝚽`$ is the dimensionless biconformal volume element. Then weight of the mass is taken as $`1,`$ so that $`\mathrm{m}^2v,v=\mathrm{m}^2v^{n+2}v^n.`$ Neglecting gravitational effects, the field equations are
$$K^{AB}D_AD_Bv^n=\mathrm{m}^2v^{n+2}$$
(45)
where indices $`A,B,\mathrm{}`$ run over the full basis, $`\omega ^A=(\omega _0^a,\omega _a^0)`$ and $`K^{AB}`$ is the projection of the conformal Killing metric to the base space, $`K^{ab}=K_{ab}=0,K_b^a=K_b^a=\delta _b^a`$. Substituting the form of $`K^{AB}`$ in the expressions above gives
$$K^{AB}D_AD_Bv^n=(D^aD_a+D_aD^a)v^m$$
(46)
While eq.(45) appears to be a straightforward classical wave equation, there are two important differences. First, we note that the presence of the mass term couples component fields of different conformal weight. This means that if the theory is not to break conformal invariance, it must be massless. If we keep the mass we necessarily produce mixing between different weight fields. Of course, such mixing is also produced by generic potentials, $`U(v).`$
To see the second difference, recall the commutator algebra satisfied by vertical tangent vectors. For simplicity, let the spacetime be flat and let $`\pi _m`$ be the canonically conjugate momentum to $`v^m.`$ Then
$$\pi _m=\frac{}{\left(\frac{v^m}{x^0}\right)}=\frac{v^m}{y_0}=^0v^m$$
(47)
Notice that, for scale-invariant actions, canonical conjugacy and conformal conjugacy always coincide. For $`v^m`$ and $`\pi _m,`$ the commutator algebra with central charges gives
$$[v,\pi ]=amv^m^0v^m$$
(48)
Therefore, we have a nonvanishing commutation relation between a field and its conjugate momentum. It is shown in that the term on the right is proportional to one component of the Weyl vector. The commutator arises because of the central charge, $`a.`$ Recall that the central charge, $`a,`$ is a necessary consequence of the exponential factor in the original group and the demand that the group act effectively.
The commutator in eq.(48) is quadratic in the fields, and therefore of the same order as the source terms for the gravitational sector of the geometry. For this reason, we cannot explore the consequences of eq.(48) further here – the only solution available is the $`m=0`$ case of a scalar field (see for a description of this solution, and for the mathematical details). More general classes of solution are under active investigation.
What can be said at this point is that in order to maintain the commutator algebra of the tangent space while using the vectors for field theory, we must use something akin to the techniques of quantum field theory. Indeed, $`v`$ is already in the form of the usual Hiesenberg operator expansion for a quantum field. Thus, in addition to successfully formalizing certain details of scale invariance, the structures arising from the extended conformal group have unexpected properties which might shed some insight onto our understanding of quantum systems.
Before concluding, we note one further possible form of the scalar field action, in which we take the norm instead of the inner product:
$$\frac{1}{2}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}𝐃v^2𝚽$$
(49)
Here we define $`𝐃v^2K^{AB}D_A\overline{v},D_Bv.`$ In this case, we can add a mass term $`\mathrm{m}^2v^2𝚽`$ only for fields $`v^k`$ of definite weight $`k=1,`$ but for this one case there is no mixing of conformal weights.
## 5 Conclusion
We have shown that our freedom to assign an arbitrary conformal weight to spacetime coordinates leads to a noncentral, discrete extension of the conformal group. This discrete extension applies to any Lie group containing dilatations. We examined the properties of the resulting extended scaling groups and their gauge theories.
We showed the presence of central charges in the extended algebra. A priori, the extended symmetries necessarily have non-trivial center. The maximal effective and transitive subgroup is therefore projective. When the resulting central charges are introduced, the discrete part of the group leads to the infinite Heisenberg algebra and the dilatational part leads to the Virasoro algebra. These algebras govern tangent vectors when extended groups are gauged, leading to nonvanishing commutators between canonically conjugate fields.
Also, extended scaling groups fill some gaps in the previous formulation of scale invariance. By formalizing the use of conformal weights as part of the initial symmetry, the covariant derivative and gauge properties are automatically appropriate to the weights of fields.
Acknowledgement The author thanks André Wehner for his careful reading of this manuscript.
Appendix A: Scale-covariant weight maps
Here we find the maximal scale-covariant subgroup of the automorphism group of the set of integer-valued conformal weights.
Let $`\mathrm{\Phi }`$ be the automorphism group of the integers, $`\mathrm{\Phi }=\{\phi |\phi :ZZ,bijective\}`$. The set $`L=\{(length)^m|mZ\}`$ provides a representation $`\mathrm{\Phi }_L`$ of $`\mathrm{\Phi }`$ by defining $`J_\phi :LL`$ as
$$J_\phi l^k=l^{\phi (k)}$$
(50)
We seek the maximal scale-covariant subgroup $`\mathrm{\Phi }_D`$ of $`\mathrm{\Phi }_L,`$ in the sense that the action of dilations, $`e^{\lambda D},`$ should be well-defined on $`\mathrm{\Phi }_D.`$ That is, $`\mathrm{\Phi }_D`$ must be closed under the action of dilatation:
$$e^{\lambda D}J_\phi e^{\lambda D}=\underset{\phi ^{}}{}f_{\phi \phi ^{}}J_\phi ^{}\mathrm{\Phi }_D$$
(51)
for all $`J_\phi `$ in $`\mathrm{\Phi }_D.`$ Consider the action of eq.(51) on $`L.`$ We compute, for all $`k,`$
$`e^{\lambda D}J_\phi e^{\lambda D}l^k`$ $`=`$ $`{\displaystyle \underset{\phi ^{}}{}}f_{\phi \phi ^{}}J_\phi ^{}l^k`$ (52)
$`e^{\lambda \phi (k)\lambda k}l^{\phi (k)}`$ $`=`$ $`{\displaystyle \underset{\phi ^{}}{}}f_{\phi \phi ^{}}l^{\phi ^{}(k)}`$ (53)
which is satisfied iff both
$`f_{\phi \phi ^{}}`$ $`=`$ $`e^\alpha \delta _{\phi \phi ^{}}`$ (54)
$`\lambda \phi (k)\lambda k`$ $`=`$ $`\alpha (\phi )`$ (55)
with $`\alpha (\phi )`$ independent of $`k.`$ Thus, for all elements of the subgroup, $`\phi =\phi ^{}`$ and $`\phi (k)=k+\alpha (\phi )/\lambda Z.`$ Therefore, the action of $`\phi `$ is characterized by a single integer $`m=`$ $`\alpha (\phi )/\lambda Z`$ and is given explicityly by $`\phi (k)=km.`$ Labelling $`J_\phi `$ as $`J_m,`$ we have
$$e^{\lambda D}J_me^{\lambda D}=e^{\lambda m}J_m$$
(56)
or infinitesimally, $`[J_m,D]=mJ_m`$.
We conclude that the set $`\mathrm{\Phi }_D=`$ $`\{J_m|J_ml^k=l^{km},mZ\}`$ with $`J_m`$ satisfying
$`[J_m,J_n]`$ $`=`$ $`0`$ (57)
$`[J_m,D]`$ $`=`$ $`mJ_m`$ (58)
is the maximal covariant subset. The set is easily seen to form a subgroup since $`J_0`$ is the identity and $`J_k`$ has inverse $`J_k.`$ Also note that the set $`\{D,J_m\}`$ is the basis for an infinite Lie algebra.
Appendix B: Nonvanishing Jacobi relations
The non-vanishing Jacobi relations for the tangent commutator algebra involve only the structure constants $`c_{\alpha \beta }^0`$ and $`c_{0\alpha }^0.`$ For the homothetic algebras, both of these vanish, while for the conformal group, $`c_{0\alpha }^0=0.`$ For a general scaling algebra, the Jacobi identities are replaced by the following Jacobi relations:
$`[G_\alpha ^k,[G_\beta ^m,G_\gamma ^n]]_{etcyc}`$ $`=`$ $`(nc_{\alpha \beta }^0\delta _\gamma ^\rho +kc_{\beta \gamma }^0\delta _\alpha ^\rho `$ (60)
$`+mc_{\gamma \alpha }^0\delta _\beta ^\rho )G_\rho ^{k+m+n}`$
$`[G_\alpha ^k,[G_\beta ^m,L^n]]_{etcyc}`$ $`=`$ $`nc_{\alpha \beta }^0L^{k+m+n}`$ (62)
$`+(mc_{0\alpha }^0\delta _\beta ^\rho kc_{0\beta }^0\delta _\alpha ^\rho )G_\rho ^{k+m+n}`$
$`[G_\alpha ^k,[G_\beta ^m,J_n]]_{etcyc}`$ $`=`$ $`nc_{\alpha \beta }^0J_{k+m+n}`$ (63)
$`[G_\alpha ^k,[L^m,L^n]]_{etcyc}`$ $`=`$ $`(mn)c_{0\alpha }^0L^{k+m+n}`$ (64)
$`[G_\alpha ^k,[L^m,J_n]]_{etcyc}`$ $`=`$ $`nc_{0\alpha }^0J_{k+m+n}`$ (65)
For the homothetic algebras, the usual Jacobi identities hold, so that the extended homothetic algebra is an infinite dimensional Lie algebra. For the conformal group (with a similar simplification for any Lie algebra with a definite weight basis, so that $`c_{0\alpha }^0=0)`$ the non-vanishing Jacobi relations are
$`[P_a^k,[K_b^m,G_\gamma ^n]]_{etcyc}`$ $`=`$ $`2(n\eta _{ab}G_\gamma ^{k+m+n}`$ (67)
$`+kc_{b\gamma }^0P_a^{k+m+n}+mc_{\gamma a}^0K_b^{k+m+n})`$
$`[P_a^k,[K_b^m,L^n]]_{etcyc}`$ $`=`$ $`2n\eta _{ab}L^{k+m+n}`$ (68)
$`[P_a^k,[K_b^m,J_n]]_{etcyc}`$ $`=`$ $`2n\eta _{ab}J_{k+m+n}`$ (69) |
warning/0002/astro-ph0002479.html | ar5iv | text | # Resonant scattering of X-rays by the warm intergalactic medium
## 1 Introduction
At high redshift the majority of baryons are contained in a photoionized IGM with temperature $`10^4\mathrm{K}`$ which is responsible for the strong Ly$`\alpha `$ absorption in the spectra of high redshift objects (see Rauch 1998 and Weinberg 1999 for reviews). At low redshift, however, there are little observational constraints on the thermal state and baryon content of the IGM (e.g. Barcons, Fabian & Rees 1991). Numerical simulations predict that a considerable fraction of all baryons is still contained in a warm IGM which traces the filamentary and sheet-like distribution of the dark matter (e.g. Ostriker & Cen, 1996, Cen & Ostriker 1999). This phase of the IGM should have densities of a few up to a few tens times the mean baryonic density. The temperatures should range between a “ photoionization temperature” at low densities ($`10^310^4\mathrm{K}`$) and the virial temperature of sheets and filaments ( $`10^6\mathrm{K}`$). These temperatures might be enhanced due to the energy input by star formation and AGN and it has been argued that such an energy input is necessary to explain the X-ray luminosity temperature relation of galaxy clusters (Kaiser 1991; Ponman, Cannon & Navarro 1999; Pen 1999). This “warm” phase of the IGM is of special interest as it will contain a record of the energy and metals expelled from galaxies. The typical surface brightness of filaments due to local thermal emission of the diffuse gas is, however, a factor 100 or more smaller than that of the X-ray background. This local thermal emission will be detectable with XMM only for especially strong filaments (Pierre, Bryan & Gastaud 1999) but see Scharf et al. (1999) for a claimed detection of such a filament with ROSAT. It was also suggested that the warm IGM produces measurable absorption in the resonant transitions of heavy elements such as oxygen or iron if there is a bright quasar behind the gas (Shapiro and Bahcall 1980, Aldcroft et al., 1994, Hellsten, Gnedin, Miralda–Escude 1998, Perna and Loeb, 1998, Markevitch, 1999). This seems a promising method to study the warm IGM, especially as long as high energy resolution is only possible for bright sources (i.e. using gratings). It will, however, only give information along the line-of-sight to point sources. We explore here the possibility to investigate the warm IGM by its emission due to resonant scattering of X-ray background photons by He and H–like ions of heavy elements. While for the high density and temperatures in galaxy clusters the local thermal emission clearly dominates over resonant scattering the opposite is true for the more moderate densities and temperatures expected in the sheet-like and filamentary structures of the warm phase of the IGM.
In section 2 we discuss the importance of resonant scattering relative to the local thermal emission due to collisional excitation and ionization for the temperatures and densities prevalent in the warm IGM. Section 3 discusses the detectability of filamentary structures with upcoming satellite missions. In section 4 we discuss the illumination of filamentary structures by individual bright sources and section 5 contains our conclusions.
A Hubble constant of $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and an Einstein-de-Sitter Universe was assumed throughout the paper. For the baryon density we take $`\rho _{\mathrm{bar}}=3.610^{31}\mathrm{g}\mathrm{cm}^3`$ ($`\mathrm{\Omega }_{\mathrm{bar}}h^2=0.02`$, Burles et al., 1999). The abundances of heavy elements were assumed to be a constant fraction of solar abundances as given by Feldman (1992). For the local XRB spectrum we use the simple approximation $`I(E)=I_0E^{1.3}e^{\frac{E}{40\mathrm{keV}}}`$ for energies $`E1\mathrm{keV}`$ and a power law with a photon index of 2 below 1 keV, where $`I_0=8\mathrm{photons}\mathrm{cm}^2\mathrm{sec}^1\mathrm{keV}^1\mathrm{sr}^1`$ (e.g. Barcons & Fabian 1992, Miyaji et al. 1998).
## 2 Photoionization balance and emissivity of the warm gas
### 2.1 Typical densities and temperatures
We discuss here the emission from typical sheet-like and filamentary structures in the warm intergalactic medium taking into account the resonant scattering of X-ray background photons. The characteristic sheets and filaments seen in numerical simulations are the result of the non-linear collapse of density perturbation imprinted onto the matter distribution in the early Universe. Sheets form from perturbations which collapse along one axis (Zeldovich 1970) while filaments have collapsed along two axis. This results in typical overdensities of a few in sheets and a few tens in filaments. The typical Thomson optical depth scales linearly with the length scale of the density perturbation and will be about $`(0.1,0.5)h^1\times 10^4(R/8h^1\mathrm{Mpc})`$ for sheets and filaments, respectively. Note that $`R=8h^1\mathrm{Mpc}`$ is approximately the scale on which the present-day density field has gone non-linear. The space density of structures larger than this decreases exponentially. Typical temperatures in the warm IGM are somewhat uncertain but will be strongly correlated with density. In the absence of energy input from star formation and AGN the temperatures at low densities are set by the balance of photoheating and adiabatic cooling and lie in the range 3000 -10000 K (e.g. Hui & Gnedin 1997). At higher densities the gas will be shock-heated. In virialized regions the temperatures will be set by the virial temperature of sheets and filaments and should be about $`10^6\mathrm{K}`$. The energy input from star formation and AGN will raise this temperatures especially at low densities. Ponman, Cannon & Navarro (1999) e.g. suggest a minimum “entropy” of $`100h^{1/3}\mathrm{keV}\mathrm{cm}^2`$ to explain the X-ray luminosity temperature relation of galaxy clusters at the faint end.
In the discussion below we consider two examples: emission from regions with an overdensity of (5,30), corresponding to a electron density of ($`10^6cm^3`$,$`610^6cm^3`$), a temperature of ($`2\times 10^5K`$,$`10^6K`$), and a metallicity of (10%,30%) solar which may resemble a typical sheet and filament, respectively. For a size of 8 Mpc a sheet and filament will have a Thomson optical depth of $`(0.2,1)\times 10^4`$ respectively. The emission spectra are calculated below for zero redshift.
### 2.2 Thermal emission of the gas in collisional equilibrium
For gas with given density and temperature in pure collisional equilibrium (coronal approximation) the following processes contribute to the X–ray emission: continuum emission (free-free and bound-free), recombination lines and emission lines excited by electron collisions. In the following we call these emission mechanisms collectively “thermal” emission. We used the code MEKA \[Mewe, Gronenschild & van den Oord 1985, Mewe, Lemen & van den Oord 1985, Kaastra 1992\] as implemented in the software package XSPEC v10 \[Arnaud 1996\] to calculate this thermal emission. For our adopted parameters the X–ray emissivity of the gas in filaments and sheets is extremely low (due to both the low temperature and the low density of the gas). Detecting these structures requires very high sensitivity. The thermal emission from a typical filament in the 0.5–1 keV energy band (in units of $`\mathrm{keV}\mathrm{sec}^1\mathrm{cm}^2\mathrm{keV}^1\mathrm{sr}^1`$ to facilitate comparison with the XRB surface brightness) is shown in Fig. 1 by the dotted curve. For comparison the intensity of the XRB is shown by the solid line. For our canonical sheet the emissivity is several orders of magnitude lower and well below the limits of the plot.
The intrinsic width of the lines should be dominated by the Hubble flow, peculiar and turbulent velocities rather than thermal broadening. In the considered energy range of 0.5–1 keV the velocity dispersion of 500-1000 $`\mathrm{km}\mathrm{s}^1`$ implies a width of $`1`$–2 eV. The spectra shown in all figures are convolved with a Gaussian with a FWHM of 2 eV. We note here that this kind of resolution may be achieved by projected X-ray missions like Constellation-X or XEUS.
### 2.3 Influence of the XRB on the thermal emission
The IGM is exposed to XRB photons. These photons change the ionization balance of the IGM, producing ions at higher ionization stage than expected for pure collisional ionization at a given temperature. The importance of photoionization depends on the temperature and density of the gas. The total ionization rate of a given ion is
$`n_e<\sigma v>_{\mathrm{ion}}+4\pi {\displaystyle I(E)\sigma _{\mathrm{ph}}(E)𝑑E},`$ (1)
where $`n_e`$ is the electron density, and $`<\sigma v>_{\mathrm{ion}}`$ is the collisional ionization rate (in $`\mathrm{cm}^3\mathrm{sec}^1`$) which is a function of temperature. The second term in eq.(1) accounts for photoionization, where $`I(E)`$ is the background intensity ($`\mathrm{photons}\mathrm{cm}^2\mathrm{sec}^1\mathrm{keV}^1\mathrm{sr}^1`$) and $`\sigma _{\mathrm{ph}}(E)`$ is the photoionization cross section. We used the approximations for photoionization cross sections given by Verner & Yakovlev (1995) and Verner et al. (1996). For simplicity we neglected ejection of multiple electrons which may follow innershell ionization. For the other processes affecting the ionization balance (i.e. collisional ionization and photo and dielectronic recombinations) we used the values adopted in the MEKA code. The characteristic time for photoionization of oxygen ions by XRB photons, $`t=\left(4\pi I(E)\sigma _{\mathrm{ph}}(E)𝑑E\right)^1`$ is $`3.5\times 10^9\mathrm{yr}`$ for OVII and $`10^{10}\mathrm{yr}`$ for OVIII. This is somewhat shorter than the Hubble time and ionization equilibrium is approximately established<sup>1</sup><sup>1</sup>1For the high abundances of OVII and OVIII in which we are interested here the recombinations time scales are of the same order.. At the densities and temperatures typical for the warm IGM the oxygen is mainly in the form of He and H–like ions, as shown in Fig.3. The symbols (dots for OVIII and circles for OVII ions respectively) show the areas on the temperature/overdensity plot where the fraction of He and H–like ions of oxygen is larger than 30%. OVII and OVIII more or less trace the density temperature relation of the warm IGM. A fraction larger than 30% is expected for practically the whole range of densities and temperatures prevalent in the warm IGM (see also Hellsten et al. 1998).
The change of the ionization balance due to photoionization affects the bound-free radiation, recombination lines and strength of the lines excited by electron collisions. The free-free emission does not change compared to the gas at the same temperature in collisional equilibrium. The corresponding thermal emission spectra (i.e. the sum of the free-free, bound-free continuums, recombination and collisionally excited lines) are shown by the grey lines in Fig.1,2 for our canonical filament and sheet, respectively. The change of the ionization state due to photoionization strongly enhances the X–ray emissivity of the gas. Photoabsorbed XRB photons are effectively converted into recombination radiation (in the form of bound-free radiation and recombination lines). The effect is especially strong when the temperature of the gas is low, e.g. around $`10^5K`$. In this case the gas does hardly emit any X–rays if photoionization is neglected.
### 2.4 Resonant scattering vs local thermal emission
Resonant scattering of XRB photons is even more important than the enhancement of the thermal emission of the gas due to photoionization. The emissivity of resonantly scattered radiation ($`\mathrm{cm}^3\mathrm{sec}^1`$) can be written as,
$`ϵ_{\mathrm{rsc}}(E_\mathrm{l})=4\pi I(E_\mathrm{l})n_i\sigma _0,`$ (2)
where $`E_\mathrm{l}`$ is the line energy and $`n_i`$ is number density of a given ion. Hereby $`\sigma _0=\frac{\pi e^2}{m_ec}f_{\mathrm{ul}}F`$, where $`f_{\mathrm{ul}}`$ is the oscillator strength, $`e`$ is the electron charge, $`m_e`$ is the electron mass, $`c`$ is the speed of light, and $`F=4.14\times 10^{18}`$ is the conversion factor from Hz to keV. The number density of a given ion is
$`n_i=n__\mathrm{H}Z_{}Zf__\mathrm{X},`$ (3)
where $`n__\mathrm{H}`$ is the hydrogen density in the warm IGM, $`Z_{}`$ is the solar abundance of a given element relative to hydrogen, $`Z`$ is the abundance of the element relative to solar abundance and $`f__\mathrm{X}`$ is the fraction of atoms in a given ionization state.
The resonantly scattered flux in a given line is proportional to the number density of the ion. It is convenient to express its intensity in terms of the intensity of the Thomson scattered continuum, which is proportional to the density of the electrons. The emissivity of the Thomson scattered continuum is given by,
$`ϵ_{\mathrm{tsc}}(E_\mathrm{l})=4\pi I(E_\mathrm{l})n_e\sigma __\mathrm{T},`$ (4)
where $`n_e`$ is the electron density and $`\sigma __\mathrm{T}`$ is the Thomson cross section. The total scattered spectrum (i.e. the sum of resonantly and Thomson scattered radiation) is shown in Fig.1,2 by the thick solid curves<sup>2</sup><sup>2</sup>2Here we use the list of the strong resonant lines compiled by Verner ($`http://www.pa.uky.edu/verner/atom.html`$)..
The ratio of resonantly and Thomson scattered emissivity (i.e. the equivalent width) is,
$`E_{_{\mathrm{EW}}}{\displaystyle \frac{\pi e^2}{m_ec}}f_{\mathrm{ul}}{\displaystyle \frac{Z_{}Zf__\mathrm{X}}{\sigma __\mathrm{T}}}`$ (5)
This expression is of course only valid for an optically thin medium. The equivalent widths for He and H–like ions are given in Table 1 assuming unity for the fraction of the ionization state and solar abundances. Along an isoelectronic sequence (e.g. for He–like ions of heavy elements) the oscillator strength for a given transition is approximately constant. Not surprisingly the resonant transitions of H and He-like oxygen (OVII and OVIII) are particularly strong due to the high oxygen abundance.<sup>3</sup><sup>3</sup>3For neutral gas an analogous relation exists between the Thomson scattered continuum and intensity of the fluorescent lines (e.g. Vainshtein, Sunyaev 1980). Resonant scattering and photoionization have a comparable effective cross sections, but an additional branching ratio (radiative decay vs autoionization – Auger effect) enters the expression for the equivalent width in the case of fluorescent lines. This branching ratio is small for light elements (e.g. $``$0.005 for the oxygen 1s-2p transition) and reaches $``$0.3 for iron. Because of this the equivalent width of the fluorescent lines from neutral gas is always much lower (except for the iron line at 6.4 keV) than in the case of resonant scattering for H- and He-like ions.
We now compare the emissivity of the gas due to scattering and that due to thermal emission (taking into account photoionization by XRB). The ratio depends strongly on the temperature and density of the gas as shown in Fig.4. Both scattered and thermal emission were integrated over the energy range 0.5–1 keV. At low densities and low temperatures the ratio is about 3–4. For these parameters photoionization strongly dominates over collisional ionization. As a result the ratio of the scattered and thermal emission (the latter being dominated by recombination radiation) is proportional to the factor $`\frac{{\scriptscriptstyle n_iI(E_\mathrm{l})\sigma _0}}{{\scriptscriptstyle n_iI(E)\sigma _{\mathrm{ph}}(E)𝑑E}}`$, where the summation in the denominator is over all ions and in the numerator over all ions and lines. For strong individual lines this ratio is about 3–4. For our canonical filament the 0.5–1 keV emissivity of the gas due to resonant scattering exceeds the thermal emission of the photoionized gas by a factor of about 2. In Fig.5 we compare the emissivity (including scattered radiation) of the gas photoionized by the XRB to the pure thermal emission of the gas if it were in collisional ionization equilibrium. The ratio of the emissivities (integrated over the 0.5–1 keV energy band) is shown as a contour plot. For our canonical filament this ratio is about 5–10 while for the sheet it is more than $`10^5`$. We conclude that estimates for the detectability of the warm IGM which neglect either scattering of the soft X-ray background and/or photoionization of the IGM by the soft X-ray background are overly pessimistic.
Finally, we compare the emissivity due to resonant scattering in the OVII line which is likely to be the strongest line at the relevant temperature and densities, to the total integrated emissivity in the energy range 0.5–1 keV (Figure 6). For a typical filament about 30 percent of the total flux is emitted in the resonant line of OVII .
### 2.5 The contribution of resonant scattering to the diffuse X-ray background
The intensity of a background scattered by a medium with uniform density is given by a formula similar to the Gunn-Peterson relation (Gunn and Peterson, 1965, Shapiro and Bahcall 1980, Aldcroft et al., 1994). In an Einstein-de-Sitter Universe this takes the form,
$`{\displaystyle \frac{I_{\mathrm{rsc}}(E_0)}{I_{_{\mathrm{XRB}}}(E_0)}}={\displaystyle \frac{c}{H_0}}ZZ_{}f__\mathrm{X}n_0{\displaystyle \frac{\sigma _0}{E_0}}(1+z)^{0.5}.`$ (6)
Here $`I_{\mathrm{rsc}}(E_0)`$ is the intensity of the scattered background at the observed energy $`E_0`$ ($`E_0<E_l`$). Equation (6) assumes that temperature, abundance and ionization state of the gas do not evolve with redshift and that the XRB is due to distant sources. For the $`1s^2`$$`1s2p^1P`$ transition of He–like oxygen the above gives approximately,
$`{\displaystyle \frac{I_{\mathrm{rsc}}(E_0)}{I_{_{\mathrm{XRB}}}(E_0)}}0.04\left({\displaystyle \frac{Z}{0.3}}\right)\left({\displaystyle \frac{f_{_{\mathrm{OVII}}}}{0.3}}\right).`$ (7)
The exact contribution of resonantly scattered photons to the XRB is difficult to assess and will depend on the detailed density, temperature and metal distribution of the IGM. From equation (7) we can, however, infer that the contribution should be at the percent level for energies below the strong resonance lines of oxygen.
From equation (7) we can also see that the warm IGM is optically thin in the strong OVII resonance line up to an overdensity of at least 30 even if the metallicity is high and the OVII fraction is large. For other lines the optical depth will be generally smaller.
## 3 Detectability of filamentary structures in the warm IGM
### 3.1 Detecting resonance scattering from diffuse gas in emission
Note that neither the number nor the energy of photons change during resonant scattering. Resonant scattering would not change the XRB flux if the XRB were completely homogeneous and isotropic<sup>4</sup><sup>4</sup>4 To a smaller extent this is also true for the photoabsorbed photons of the XRB which are reemitted as recombination photons.. However, the XRB is emitted by discrete sources and resonance scattering converts photons emitted by compact sources into a diffuse background.
The resonantly scattered emission will be detectable in images in which discrete sources making up a significant fraction of the total background have been removed. This is demonstrated in Fig.7 which shows the spectrum of our canonical filament (taking into account the XRB) and the case were 90 percent of the background has been removed. ROSAT has e.g. resolved 70-80% of the XRB in the energy range 0.5–2 keV (Hasinger et al., 1998). At the lower end of this range the situation is somewhat unclear mainly due to the emission of our own Galaxy.
The spectral resolution of X-ray instruments is dramatically increasing. XMM and Chandra will have the first instruments with high spectral resolution (using gratings) but only for bright compact sources. With projected missions like Constellation-X and XEUS imaging with a spectral resolution of about 2 eV will become possible. The large ratio of resonantly scattered emission to the Thomson scattered continuum and the local thermal emission makes the search for the spectral feature of resonantly scattered radiation very worthwhile. Detecting the warm IGM by photons scattered in the OVII or OVIII resonant transitions has another big advantage. Filaments and sheets contain a large number of faint galaxy clusters and galaxy groups. The emission from the dense hot gas in these clusters and groups generally dominates the thermal emission from filaments. However, in this dense and hot gas oxygen is generally completely stripped from electrons and no resonant scattering will occur in these regions. The resonantly scattered radiation will thus be a good tracer of the diffuse gas in filaments. The only contamination should be due to resonant oxygen scattering by the warm gas which may be contained in the galaxies within the filaments.
The highest signal to noise ratio can be achived if OVII or OVIII absorption lines are observed in the spectrum of a very bright compact source located behind the filament as first suggested by Shapiro and Bahcall (1980). The longest exposures (deep surveys) however are usually collected for fields without strong X–ray sources. For such fields the detection of OVII or OVIII resonant lines in emission in the residual background is favorable compared to the detection of these lines in absorption using the combined spectrum of all resolved background sources in the field. The number of line photons in the residual background is approximately equal to the number of line photons absorbed from the spectra of all resolved sources. The intensity of the residual background is, however, lower if most of the background is resolved. This results in a higher signal-to-noise ratio for the residual background. It means that in order for such a method to be useful the particle and other “non cosmological” detector backgrounds have to be low compared to the combined intensity of the resolved sources in the field.
### 3.2 Detecting filaments with upcoming satellite missions
In an Einstein de Sitter Universe our canonical filament of comoving size ($`16h^1\mathrm{Mpc}\times 3h^1\mathrm{Mpc}`$ has angular extent $`(183^{}\times 34^{},50^{}\times 9^{},31^{}\times 6^{})`$ at $`z=(0.1,0.5,1)`$, respectively. In the case of a filament, completely filling the FOV the expected number of counts in the oxygen line is proportional to the product of the effective area and the FOV (solid angle) of the telescope. This value is largest for XMM and is comparable for the Wide field imager of XEUS. Estimates of the effective area, FOV and energy resolution are summarized in Table 2. For a nearby ($`z0.1`$) filament with Thomson optical depth of $`10^4`$ XMM (EPIC pn detector) will detect about 200 counts in the OVII line in a $`10^5`$ s exposure. For Chandra the expected number of counts is about an order of magnitude lower.
The total particle background in a 60 eV band pass (the resolution of XMM at 0.5 keV) will be at the level of $``$50–60 counts for a $`10^5`$ s exposure. The residual unresolved extragalactic background (assuming that 90% will be resolved) will produce $``$ 1000 counts. The biggest problem is probably Galactic emission. Galactic emission is usually modeled as a combination of two thermal (i.e. thermal emission of optically thin plasma) components with temperatures $``$0.15 and 0.05 keV respectively (e.g. Miyaji et al., 1998). A redshift of 0.1 should, however, be sufficient to distinguish oxygen lines produced by a filament from oxygen lines produced in the Galaxy. Galactic emission will produce several thousand counts in a 60 eV wide band pass. For a $`10^5`$ second XMM exposure the signal-to-noise ratio will thus be of order a few. Similar signal-to-noise ratios can be achieved for the OVIII line at 0.65 keV. Such signal-to-noise ratios should be sufficient for a meaningful cross-correlation with the redshift distribution of galaxies in the same field. Longer exposure times, filaments with larger Thomson optical depth, mapping of larger areas of sky or illumination by a nearby bright X-ray source as described in the next section will be necessary to increase the signal-to-noise ratio.
The anticipated FOV of Constellation X and of the Narrow field imager of XEUS is rather small (see Table 2). The total number of photons detected in the OVII and OVIII lines for a filament filling the whole FOV will therefore be an order of magnitude lower for these missions than for XMM for a similar exposure time(except for the Wide field imager of XEUS, which will detect a similar number of photons). The signal-to-noise ratio (assuming pure statistical errors due to residual unresolved extragalatic and Galactic background within the band pass set by the energy resolution of the telescope) will nevertheless be similar due to the better energy resolution of XEUS and Constellation X. For an efficient study of the diffuse emission of filaments a larger FOV would be very important. If it were possible to increase the FOV to e.g. $`5^{}\times 5^{}`$ (below 1 keV) for high energy resolution detectors like those of Constellation X and XEUS then these missions should be able to detect resonantly scattered emission of filaments along any line of sight during a 50ks exposure.
## 4 Illumination by a nearby bright source
Close to a bright source the emission due to resonant scattering is obviously enhanced both relative to the local emission and relative to the background. Bright sources may have soft X-ray luminosities of $`10^{46}\mathrm{erg}\mathrm{sec}^1`$ or larger and exceed the soft X-ray background by a factor $`(r/r_0)^2`$ out to a radius $`r_030(L/10^{46}\mathrm{erg}\mathrm{sec}^1)^{0.5}\mathrm{Mpc}`$ and the signal-to-noise ratio will be enhanced by the same factor. Particular interesting is the illumination of strong filaments close to rich galaxy clusters by either the X-ray emission from the core of the cluster or by a bright AGN contained in the cluster. This is illustrated in Fig.8, where we show the intrinsic and scattered emission for the same parameters as in Fig. 1, but with a factor 100 larger intensity of the illuminating continuum. For simplicity we assumed the same shape for the quasar spectrum as for the XRB. The increased X-ray flux also affects the ionization state of the gas. The relevant quantity is the ionization parameter which is proportional to the ionizing flux. The regions in Fig.3 will e.g. shift linearly to the right with increasing flux (at least approximately for the moderate densities and temperatures where collisional ionization is not important). This shifts the dominant ionization state from OVII to OVIII. The strongest line in Fig.8 is indeed that of OVIII at 0.65 keV. The local thermal emission is also strongly enhanced (compare the dotted to the grey and thick solid curves in Fig.8). Very close to bright sources oxygen may be even completely stripped due to photoionization. The fraction of H–like ions is then inversely proportional to the blazar flux while the number of scattered photons (e.g. in the H-like 1s-2p transition) is proportional to the blazar flux and the density of H-like ions. The resonantly scattered flux will then not depend on the illuminating flux anymore and will instead be proportional to the recombination rate. The number of resonantly scattered photons will still be larger than the number of recombination line photons by a factor of few. On the other hand the Thomson scattered flux is simply proportional to the intensity of illuminating radiation. The relative importance of the Thomson scattered continuum is thus strongly enhanced close to a strong ionizing source.
The number of bright AGN which are available for an illumination of filaments is much larger than naively inferred from the AGN luminosity function. A considerable fraction of bright AGN are blazars with strongly beamed radiation. If these typically have an opening angle of 5 degrees there will be about 2000 objects which point away of us for any observed blazar.
Note that in the case of illumination by individual bright sources the assumption of an isotropic phase function for the resonant scattering which we made so far is not valid. The phase function will depend on the particular transition. For the $`1s^2`$$`1s2p^1P`$ transitions of He–like ions the phase function will e.g. be the familiar dipole (see e.g. Chandrasekhar 1960) and the scattered radiation will be polarized.
The illumination of filaments by an AGN may also constrain the lifetime of these objects. If the active phase of AGN lasts for about $`10^7`$years then the observed size of the illuminated region will be of the order of 3 Mpc (depending on the angle of the beam to the line of sight). Of course the observed flux in the scattered lines depends on a number of factors (density, temperature, blazar flux, heavy elements abundance etc.). Various structures in the surface brightness which may appear as a result of scattering of beamed radiation are considered by Gilfanov, Sunyaev & Churazov (1987) and Wise & Sarazin 1990.
The time, needed to establish photoionization equilibrium is inversely proportional to the photoionization rate and is equal to $`10^9\left(\frac{D}{20\mathrm{Mpc}}\right)^2\left(\frac{L}{10^{47}\mathrm{erg}\mathrm{sec}^1}\right)^1\mathrm{yr}`$. Thus if the blazar was active for a shorter period of time or the beam direction varies (e.g. due to precession of the jet) then the influence of the photoionization by the blazar emission will be lower accordingly. The analysis of gas illuminated by a nearby sources is therefore not straightforward. It will, however, be easier to interpret for beamed sources than for isotropic sources.
## 5 Conclusions
We have demonstrated that for the temperatures and densities prevalent in the the filamentary and sheet-like structures of the warm intergalactic medium illumination by XRB photons significantly increases the X–ray emissivity of the medium. The strongest contribution to the X–ray emissivity is due to resonant scattering of XRB photons. This resonantly scattered emission can be detected in images from which a significant fraction of the XRB due to compact sources has been removed.
Resonant lines of He and H-like oxygen are most promising for detecting filamentary structures in the warm IGM. The fraction of oxygen in the ionization states OVII or OVIII is larger than 30 percent practically for the whole range of densities and temperatures expected in the warm IGM.
Estimates for the detectability of filaments in the warm IGM which take only local thermal emission into account may have been overly pessimistic. XMM should be able to detect the radiation resonantly scattered by diffuse gas in a filament with Thomson optical depth of $`10^4`$ and metallicity 0.3 solar at $`z0.1`$ with signal-to-noise of a few.
For filaments close to massive cluster the signal will be enhanced due to illumination from the cluster core and/or AGN contained in the cluster. The most promising sources emitting high X-ray intensities are beamed blazars which should have a high space density. For AGN and especially blazars constraints on the duration of the active phase may also be obtainable.
## 6 Acknowledgements
We thank the referee K. Yamashita for helpful comments. MH gratefully acknowledges the hospitality of the Institute for Theoretical Physics Santa Barbara where this research was completed. This research was supported in part by the National Science Foundation under Grant No. Phy94-07194. |
warning/0002/physics0002035.html | ar5iv | text | # The Emergence of the Planck Scale
## 1 Introduction
About a century ago Max Planck had pointed out that the quantity $`\left(\frac{\mathrm{}G}{c^3}\right)^{\frac{1}{2}}10^{33}cms`$ is a fundamental length. This so called Planck length ties up Quantum Mechanics, Gravitation and Special Relativity and leads to the Planck mass $`10^5gms`$. It is but natural that the Planck length has played a crucial role in Quantum Gravity as also in String Theory which includes a description of Gravitation, unlike Quantum Theory or Quantum Field Theory.
It turns out to be the scale at which we have no longer the smooth space time of Classical Theory and Quantum Theory, but rather we have the space time foam of Wheeler. This is inextricably linked with gravitational collapse which has been described by Wheeler as ”The greatest crisis of Physics”. As he puts it, ”These are small scale fluctuations telling one that something like gravitational collapse is taking place everywhere in space and all the time; that gravitational collapse is in effect perpetually being done and undone …. at the Planck scale of distances.” In this space time foam, worm holes and non local effects abound.
On the other hand there is also a stochastic fluctuational picture of space time that deals with phenomena at the Compton wavelength scale and leads to meaningful physics and cosmology including a unified description of gravitation and electromagnetism consistent with observation. In this picture, space time has been considered to be a random heap of elementary particles. If we consider a typical elementary particle to be a pion with Compton wavelength $`l`$, then the above picture leads to a dispersion length in the Gaussian distribution $`\sqrt{N}l,N10^{80}`$ being the number of elementary particles in the universe, this being the correct dimension of the universe itself.
We will now show a parallel between the Planck length considerations and the Compton wavelength considerations referred to above, which will then show us how the Planck length considerations emerge.
## 2 The Emergence of the Planck Scale
We first show the parallels between the Compton wavelength picture and the Planck length picture. We note that in the former scenario, particles are fluctuationally created at the Compton wavelength from a background pre space time Zero Point Field (ZPF) of the kind considered in stochastic electrodynamics\[9, r10\]. The energy content in terms of the magnetic field of such a particle is given by (Cf.ref.)
$$\mathrm{\Delta }B\frac{(\mathrm{}c)^{1/2}}{L^2}$$
(1)
where $`L`$ is the dimension under consideration, which in this case is of the order of the particle’s Compton wavelength. We note that in (1) if $`\mathrm{}c`$ or equivalently $`137e^2`$ is replaced by its gravitational counterpart, namely $`137Gm^2`$ then we get, as in the fluctuation of the metric,
$$\mathrm{\Delta }g\frac{L_P}{L}$$
(2)
where $`L_P`$ is the Planck length and $`L`$ as in (1) is of the order of the dimension under consideration.
The space time foam referred to above arises at the Planck scale because the right hand side in (2) becomes unity, indicating perpetual collapse and creation.
From this point of view, as Wheeler points out our space time is an approximation, an average swathe at the Planck scale of several probable spaces and topologies which form the super space (Cf.ref.). There is an immediately parallel in terms of the Compton wavelength considerations also: As pointed out by Nottale, Abbot-wise, El Naschie, the author and others the Quantum behaviour below a critical length is fractal and as pointed out by the author, our space time is the thick brush stroke of thickness of the order of the Compton wavelength of a jagged, fractal coastline like underpinning.
In the light of the above considerations the fluctuational creation of particles considered by Hayakawa and the author have a parallel in the non local worm hole related appearance of particles and fields at the Planck scale.
We will now quantify the above parallels and show the actual emergence of the Planck scale particles from the Compton wavelength considerations.
We first observe that in an actual random heap of particles, the smaller particles (in our case those having smaller Compton wavelengths and therefore higher mass) tend to settle down together due to gravity. In a fluctuationally created random heap of particles, there is no gravity, but as this space time heap is not only non differentiable, but is also not required to be even a continuum the random motion would have a similar effect: Of the $`N^{}=\sqrt{N}`$ particles which are less dispersed, $`\sqrt{N^{}}`$ particles would similarly fluctuationally, that is non locally be together. This fluctuationally bound group would have a mass $`\sqrt{N^{}}m10^5gms`$ or the Planck mass, since $`m`$ is the mass of the pion. (Cf.ref.”Ramification” for another interesting perspective).
One way of looking at this is that in the above scenario, space time no longer has the rigid features of Classical and Quantum Physics - on the average it is a measure of dispersion of a random distribution of particles which themselves have a stochastic underpinning. So the length scale or dispersion would be less, the less dispersed the random collection of particles is - this leads to the Planck scale from the Compton scale. However it must be borne in mind that a Planck mass has a life time $`10^{42}`$ seconds, and can hardly be detected.
The Planck scale corresponds to the extreme classical limit of Quantum Mechanics, as can be immediately seen from the fact that the Planck mass $`m_P10^5gms`$ corresponds to a Schwarzchild Black Hole of radius $`L_P10^{33}cms`$, the Planck length. At this stage the spinorial Quantum Mechanical feature as brought out by the Kerr-Newman type Black Hole and the Compton wavelength (Cf.detailed discussion in refs.) disappears. Infact at the Planck scale we have
$$\frac{Gm_P}{c^2}=\mathrm{}/m_Pc$$
(3)
In (3), the left side gives the Schwarzchild radius while the right side gives the Quantum Mechanical Compton wavelength. Another way of writing (3) is,
$$\frac{Gm_P^2}{e^2}1,$$
(4)
Equation (4) expresses the well known fact that at this scale the entire energy is gravitational, rather than electromagnetic, in contrast to equation (1) for a typical elementary particle mass, vi.,
$$Gm^2\frac{1}{\sqrt{N}}e^210^{40}e^2$$
Interestingly from the background ZPF, Planck particles can be produced at the Planck scales given by (3), exactly as in the case of pions, as seen earlier. They have been considered to be what may be called a Zero Point Scale. But these shortlived Planck particles can at best describe a space time foam.
We will now throw further light on the fact that at the Planck scale it is gravitation alone that manifests itself. Indeed Rosen has pointed out that one could use a Schrodinger equation with a gravitational interaction to deduce a mini universe, namely the Planck particle. The Schrodinger equation for a self gravitating particle has also been considered, from a different point of view. We merely quote the main results.
The energy of such a particle is given by
$$\frac{Gm^2}{L}\frac{2m^5G^2}{\mathrm{}^2}$$
(5)
where
$$L=\frac{\mathrm{}^2}{2m^3G}$$
(6)
(5) and (6) bring out the characteristic of the Planck particles and also the difference with elementary particles, as we will now see.
We first observe that for a Planck mass, (5) gives, self consistently,
$$\text{Energy}=m_Pc^2,$$
while (6) gives,
$$L=10^{33}cms,$$
as required.
However, the situation for pions is different: They are parts of the universe and do not constitute a mini universe. Indeed, if, as above there are $`N`$ pions in the universe, then the total gravitational energy is given by, from (5),
$$\frac{NGm^2}{L}$$
where now $`L`$ stands for the radius of the universe $`10^{28}cm`$. As this equals $`mc^2`$, we get back as can easily be verified, the pion mass!
Indeed given the pion mass, one can verify from (6) that $`L=10^{28}cms`$ which is the radius of the universe, $`R`$. Remembering that $`R\frac{c}{H}`$, (6) infact gives back the supposedly mysterious and adhoc Weinberg formula, relating the Hubble constant to the pion mass.
This provides a justification for taking a pion as a typical particle of the universe, and not a Planck particle, besides re-emphasizing the basic unified picture of gravitation and electromagnetism. It must be mentioned that just as the Planck particle constitutes a mini universe or Black Hole, so also the $`N10^{80}`$ pion filled universe can itself considered to be a Black hole!
To proceed, let us now use the fact that our minimum space time intervals are $`(l_P,\tau _P)`$, the Planck scale, instead of $`(l,\tau )`$ of the pion, as above.
With this new limit, it can be easily verified that the total mass in the volume $`l^3`$ is given by
$$\rho _P\times l^3=M$$
(7)
where $`\rho _P`$ is the Planck density and $`M`$ is the mass of the universe.
Moreover the number of Planck masses in the above volume $`l^3`$ can easily be seen to be $`10^{60}`$. However, it must be remembered that in the physical time period $`\tau `$, there are $`10^{20}`$ (that is $`\frac{\tau }{\tau _P})`$ Planck life times. In other words the number of Planck particles in the physical interval $`(l,\tau )`$ is $`N10^{80}`$, the total particle number, as if all these were the seeds of the fixed number of $`N`$ particles in the universe. This is symptomatic of the fact that instead of the elementary particle Compton wavelength scale of the physical universe we are using the Planck scale (cf. also considerations before equation (3)).
That is from the typical physical interval $`(l,\tau )`$ we recover the entire mass and also the entire number of particles in the universe, as in the Big Bang theory. This also provides the explanation for the above puzzling relations like (7).
That is the Big Bang theory is a characterization of the new Compton wavelength model in the classical limit at Planck scales, but then, in this latter case we cannot deduce from theory the relations like the Dirac coincidences or the Weinberg formula.
In the spirit of, one can now see the semi-classical and Quantum Mechanical divide between Planck particles and elementary particles in the following way. We will see that Planck particles have a life time given by the Hawking Radiation Law of Black Hole Thermodynamics, whereas elementary particles are characterised by Quantum Mechanical life times.
It is well known that the life time due to the Hawking Radiation Law is given by
$$t=\frac{G^2m^3}{\mathrm{}c^4}$$
(8)
which for the Planck particles gives the usual Planck time.
However this formulation is not valid for elementary particles. In this case, we consider the gravitational energy $`\mathrm{\Delta }E`$ of a pion as given by an equation like (5) and use instead the Quantum Mechanical relation
$$\mathrm{\Delta }E.\mathrm{\Delta }t\mathrm{}$$
(9)
to get
$$Gm_\pi ^2(\mathrm{}/m_\pi c)\mathrm{\Delta }t\mathrm{}$$
(10)
which is correct if in (9) $`\mathrm{\Delta }t\frac{1}{H}`$, the age of the universe! (cf.also ref.)). In this case equation (10) gives the well known and supposedly mysterious and empirical formula of Weinberg referred to earlier, viz.,
$$m_\pi ^3\frac{H\mathrm{}^2}{Gc}$$
(11)
One way of looking at this is that it is the emergence of Quantum Mechanical effects and electromagnetism at the Compton wavelength scales from classical gravitational considerations at the Planck scale as seen above, which gives stability to the universe as expressed by (9) and (10).
All this has been justified from stochastic considerations.
Another way of looking at all this is the following: The gravitational constant $`G`$ is taken to be a universal constant in most conventional theories. However in the above formulation it turns out that,
$$G=\frac{G_0}{\sqrt{N}}\frac{1}{T}$$
(12)
where $`N`$ is the number of elementary particles in the universe and $`T`$ is the age of the universe. This time varying gravitational constant can be shown to lead to consistent results including an explanation for the all important precision of the perihelion of the Planet Mercury . The equation (12) also shows a Machian or holistic character. In any case for a single particle universe, $`N=1`$ the $`G`$ above leads to the Planck length or Planck mass, while for $`N10^{80}`$ the same equation leads to the pion Compton wavelength and the usual Physics and Cosmology. Infact if the pion Compton time scales $`(l,\tau )`$ tends to zero or the Planck scale we recover the big bang scenario and the usual space time of Classical and Quantum Physics or the Prigogine Cosmology. In these cases we cannot explain the large number ”coincidences” and Weinberg’s mysterious formula (11), whereas at the elementary particle Compton scale these features can be deduced as consequences of the theory. |
warning/0002/quant-ph0002067.html | ar5iv | text | # Rules for Integrals over Products of Distributions from Coordinate Independence of Path Integrals
## I Introduction
While quantum mechanical path integrals in curvilinear coordinates can be defined uniquely and independently of the choice of coordinates within the time-sliced formalism , a perturbative definition on a continuous time axis poses problems. To exhibit the difficulties, consider the associated partition function calculated for periodic paths on the imaginary-time axis $`\tau `$:
$$Z=𝒟q(\tau )\sqrt{g}e^{𝒜[q]},$$
(1)
where $`𝒜[q]`$ is the euclidean action with the general form
$$𝒜[q]=𝑑\tau \left[\frac{1}{2}g_{\mu \nu }(q(\tau ))\dot{q}^\mu (\tau )\dot{q}^\nu (\tau )+V(q(\tau ))\right].$$
(2)
The dots denote $`\tau `$-derivatives, $`g_{\mu \nu }(q)`$ is a metric, and $`g=detg`$ its determinant. The path integral may formally be defined perturbatively as follows: The metric $`g_{\mu \nu }(q)`$ is expanded around some point $`q_0^\mu `$ in powers of $`\delta q^\mu q^\mu q_0^\mu `$. The same thing is done with the potential $`V(q)`$. After this, the action $`𝒜[q]`$ is separated into a free part $`𝒜_0[q_0;\delta q]\frac{1}{2}g_{\mu \nu }(q_0))\dot{q}^\mu \dot{q}^\nu +\frac{1}{2}\omega ^2\delta q^\mu \delta q^\nu `$, and an interacting part $`𝒜_{\mathrm{int}}[q_0;\delta q]𝒜[q]𝒜_0[q_0;\delta q]`$.
A first problem is encountered in the square root in the functional integration measure in (1). Taking it into the exponent and expanding it in powers of $`\delta q`$ we define an effective action$`𝒜_\sqrt{g}=\frac{1}{2}\delta (0)\mathrm{log}[g(q_0+\delta q)/g(q_0)]`$, which contains the $`\delta `$-function at the origin $`\delta (0)`$. It represents formally the inverse infinitesimal lattice spacing on the time axis, and is equal to the infinite number $`\delta (0)𝑑p/(2\pi )`$.
The second problem arises in the expansion of $`Z`$ in powers of the interaction, Performing all Wick contractions, $`Z`$ is expressed as a sum of loop diagrams. There are interaction terms involving $`\dot{q}^2q^n`$ which lead to Feynman integrals over products of distributions. The diagrams contain three types of lines representing the correlation functions
$`\mathrm{\Delta }(\tau \tau ^{})`$ $``$ $`q(\tau )q(\tau ^{})=\text{ }\text{}\text{ },`$ (3)
$`_\tau \mathrm{\Delta }(\tau \tau ^{})`$ $``$ $`\dot{q}(\tau )q(\tau ^{})=\text{ }\text{}\text{ },`$ (4)
$`_\tau _\tau ^{}\mathrm{\Delta }(\tau \tau ^{})`$ $``$ $`\dot{q}(\tau )\dot{q}(\tau ^{})=\text{ }\text{}\text{ }.`$ (5)
The right-hand sides define the line symbols to be used in Feynman diagrams for the interaction terms.
Explictly, the first correlation function reads
$$\mathrm{\Delta }(\tau ,\tau ^{})=\frac{1}{2\omega }e^{\omega |\tau \tau ^{}|}.$$
(6)
The second correlation function (4) has a discontinuity
$$_\tau \mathrm{\Delta }(\tau ,\tau ^{})=\frac{1}{2}ϵ(\tau \tau ^{})e^{\omega |\tau \tau ^{}|},$$
(7)
where
$$ϵ(\tau \tau ^{})1+2_{\mathrm{}}^\tau 𝑑\tau ^{\prime \prime }\delta (\tau ^{\prime \prime }\tau ^{})$$
(8)
is a distribution which vanishes at the origin and is equal to $`\pm 1`$ for positive and negative arguments, respectively. The third correlation function (5) contains a $`\delta `$-function:
$$_\tau _\tau ^{}\mathrm{\Delta }(\tau ,\tau ^{})=\delta (\tau \tau ^{})\frac{\omega }{2}e^{\omega |\tau \tau ^{}|},$$
(9)
The temporal integrals over products of such distributions are undefined .
In this note we define them uniquely by setting up rules between these and integrals over products of nonsingular correlation functions $`\mathrm{\Delta }(\tau \tau ^{})`$, plus integrals over pure products of $`\delta `$-functions. These will be defined uniquely by the requirement of coordinate invariance of the path integral (1).
The internal consistency of these definitions is ensured by previous work of the present authors. In Ref. , we have shown that Feynman integrals in momentum space can be uniquely defined as $`ϵ0`$ -limits of $`1ϵ`$-dimensional integrals via an analytic continuation à la ’t Hooft and M. Veltman . This definition makes path integrals coordinate independent. In Ref. we have given rules for calculating the same results directly from the Feynman integrals in the $`1ϵ`$ -dimensional time space.
The present approach has the important advantage making superfluous the somewhat tedious analytic continuation to $`1ϵ`$ dimensions. In fact, it does not require specifying any regularization scheme. In addition, it gives a foundation of a new and general mathematics of extending the theory of distributions from a linear space to products.
## II Model System
The announced derivation of the identities will be based on the requirement of coordinate independence of the exactly solvable path integral of a point particle of unit mass in a harmonic potential $`\omega ^2x^2/2`$, over a large imaginary-time interval $`\beta `$,
$$Z_\omega =𝒟x(\tau )e^{𝒜_\omega [x]}=e^{\mathrm{Tr}\mathrm{log}(^2+\omega ^2)}=e^{\beta \omega /2}.$$
(10)
The action is
$$𝒜_\omega =\frac{1}{2}𝑑\tau \left[\dot{x}^2(\tau )+\omega ^2x^2(\tau )\right].$$
(11)
A coordinate transformation turns (10) into a path integral of the type (1) with a singular perturbation expansion. From our work in Refs. we know that all terms in this expansion vanish in dimensional regularization. Here we shall require the vanishing to find the desired identities for integrals over products of distributions.
For simplicity we assume the coordinate transformation to preserve the symmetry $`xx`$ of the initial oscillator, such its power series expansion starts out like $`x(\tau )=f(q(\tau ))=qgq^3/3+g^2aq^5/5\mathrm{}`$, where $`g`$ is a smallness parameter, and $`a`$ an extra parameter. We shall see that the identities are independent of $`a`$, such that $`a`$ will merely serve to check the calculations. The transformation changes the partition function (10) into
$$Z=𝒟q(\tau )e^{𝒜_J[q]}e^{𝒜[q]},$$
(12)
where is $`𝒜[q]`$ is the transformed action, whereas $`𝒜_J[q]`$ an effective action coming from the Jacobian of the coordinate transformation:
$$𝒜_J[q]=\delta (0)𝑑\tau \mathrm{log}\frac{\delta f(q(\tau ))}{\delta q(\tau )}.$$
(13)
The transformed action is decomposed into a free part
$$𝒜_\omega [q]=\frac{1}{2}𝑑\tau [\dot{q}^2(\tau )+\omega ^2q^2(\tau )],$$
(14)
and an interacting part, which reads to second order in $`g`$:
$`𝒜_{\mathrm{int}}[q]={\displaystyle \frac{1}{2}}{\displaystyle }d\tau \{g[2\dot{q}^2(\tau )q^2(\tau )+{\displaystyle \frac{2\omega ^2}{3}}q^4(\tau )]`$ (15)
$`+g^2[(1+2a)\dot{q}^2(\tau )q^4(\tau )+\omega ^2({\displaystyle \frac{1}{9}}+{\displaystyle \frac{2a}{5}})q^6(\tau )]\}.`$ (16)
To the same order in $`g`$, the Jacobian action (13) is
$`𝒜_J[q]=\delta (0){\displaystyle 𝑑\tau \left[gq^2(\tau )+g^2\left(a\frac{1}{2}\right)q^4(\tau )\right]}.`$ (17)
For $`g=0`$, the transformed partition function (12) coincides with (10). When expanding $`Z`$ of Eq. (12) in powers of $`g`$, we obtain sums of Feynman diagrams contributing to each order $`g^n`$, which must vanish to ensure coordinate invariance. By considering only connected Feynman diagrams, we are dealing directly with the ground state energy.
## III Free Energy Density
The graphical expansion for the ground state energy will be carried here only up to three loops. At any order $`g^n`$, there exist different types Feynman diagrams with $`L=n+1,n,`$ and $`n1`$ number of loops coming from the interaction terms (16) and (17), respectively. The diagrams are composed of the three types of lines in (3)–(5), and new interaction vertices for each power of $`g`$. The diagrams coming from the Jacobian action (17) are easily recognized by an accompanying power of $`\delta (0)`$.
At first order in $`g`$, there exists only three diagrams, two originated from the interaction (16), one from the Jacobian action (17):
$$g\text{ }\text{}\text{ }g\omega ^2\text{ }\text{}\text{ }+g\delta (0)\text{ }\text{}\text{ }.$$
(18)
At order $`g^2`$, we distinguish several contributions. First there are two three-loop local diagrams coming from the interaction (16), and one two-loop local diagram from the Jacobian action (17):
$`g^2[\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}3}({\displaystyle \frac{1}{2}}+a)\text{ }\text{}\text{ }+\mathrm{\hspace{0.17em}15}\omega ^2({\displaystyle \frac{1}{18}}+{\displaystyle \frac{a}{5}})\text{ }\text{}\text{ }`$ (19)
(20)
$`3(a{\displaystyle \frac{1}{2}})\delta (0)\text{ }\text{}\text{ }].`$ (21)
We call a diagram local if it involves no temporal time integral. The Jacobian action (17) contributes further the nonlocal diagrams:
$`{\displaystyle \frac{g^2}{2!}}\left\{2\delta ^2(0)\text{ }\text{}\text{ }4\delta (0)[\text{ }\text{}\text{ }+\text{ }\text{}\text{ }+2\omega ^4\text{ }\text{}\text{ }]\right\}.`$ (22)
(23)
The remaining diagrams come from the interaction (16) only. They are either of the three-bubble type, or of the watermelon type, each with all possible combinations of the three line types (3)–(5): The sum of all three-bubbles diagrams is
$`{\displaystyle \frac{g^2}{2!}}[4\text{ }\text{}\text{ }+\mathrm{\hspace{0.17em}\hspace{0.17em}2}\text{ }\text{}\text{ }+\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}2}\text{ }\text{}`$ (24)
$`+8\omega ^2\text{ }\text{}\text{ }+8\omega ^2\text{ }\text{}\text{ }+8\omega ^4\text{ }\text{}\text{ }],`$ (25)
while the watermelon-like diagrams contribute
$`{\displaystyle \frac{g^2}{2!}}\mathrm{\hspace{0.17em}4}[\text{ }\text{}\text{ }+4\text{ }\text{}\text{ }+\text{ }\text{}\text{ }+4\omega ^2\text{ }\text{}\text{ }+{\displaystyle \frac{2}{3}}\omega ^4\text{ }\text{}\text{ }].`$ (26)
Since the equal-time expectation value $`\dot{q}(\tau )q(\tau )`$ vanishes according to Eq. (7), there are a number of trivially vanishing diagrams, which have been omitted.
In our previous papers , all integrals were calculated individually in $`D=1\epsilon `$ dimensions, taking the limit $`\epsilon 0`$ at the end. The results for the integrals ensured that the sum of all Feynman diagrams contributing to each order $`g^n`$ vanishes.
## IV Rules for Integrals over Distributions
As a first step in calculating the Feynman integrals we express singular time derivatives $`\dot{\mathrm{\Delta }}(\tau )`$, $`\ddot{\mathrm{\Delta }}(\tau )`$ in terms of regular correlation functions $`\mathrm{\Delta }(\tau )`$, plus integrals over powers of $`\delta `$-functions. The tools for this will be partial integrations and the inhomogeneous field equation satisfied by the correlation function
$$\ddot{\mathrm{\Delta }}(\tau )=\overline{}𝑑k\frac{k^2}{k^2+\omega ^2}e^{ik\tau }=\delta (\tau )+\omega ^2\mathrm{\Delta }(\tau ),$$
(27)
Most simply, we have for integrals over products of two correlation functions the relation
$$𝑑\tau \left[\dot{\mathrm{\Delta }}^2(\tau )+\omega ^2\mathrm{\Delta }^2(\tau )\right]=\mathrm{\Delta }(0).$$
(28)
To prove this, we integrate the first term partially,
$$𝑑\tau \dot{\mathrm{\Delta }}^2(\tau )=𝑑\tau \mathrm{\Delta }(\tau )\ddot{\mathrm{\Delta }}(\tau ),$$
(29)
with no boundary term due to the exponential vanishing at infinity of all functions involved. Using now the field equation (27) and the property of $`\delta `$-function that
$$𝑑\tau f(\tau )\delta (\tau )=f(0),$$
(30)
for any smooth test function $`f(\tau )`$, we obtain (28).
We now turn to singular integrals involving $`\ddot{\mathrm{\Delta }}^2(\tau )`$. Using the same tools, we obtain, in the same way, the relation
$`{\displaystyle 𝑑\tau \left[\ddot{\mathrm{\Delta }}^2(\tau )+2\omega ^2\dot{\mathrm{\Delta }}^2(\tau )+\omega ^4\mathrm{\Delta }^2(\tau )\right]}={\displaystyle 𝑑\tau \delta ^2(\tau )}.`$ (31)
The last integral is undefined. Before fixing its value in the next section, we shall derive relations for integrals over singular products of four correlation functions. First for $`\ddot{\mathrm{\Delta }}(\tau )\mathrm{\Delta }^3(\tau )`$. Using again the field equation (27), we find
$$𝑑\tau \ddot{\mathrm{\Delta }}(\tau )\mathrm{\Delta }^3(\tau )=\mathrm{\Delta }^3(0)\omega ^2𝑑\tau \mathrm{\Delta }^4(\tau ).$$
(32)
By a partial integration, the left-hand side becomes
$$𝑑\tau \ddot{\mathrm{\Delta }}(\tau )\mathrm{\Delta }^3(\tau )=3𝑑\tau \dot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }^2(\tau ),$$
(33)
leading to
$$𝑑\tau \dot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }^2(\tau )=\frac{1}{3}\mathrm{\Delta }^3(0)\frac{1}{3}\omega ^2𝑑\tau \mathrm{\Delta }^4(\tau ).$$
(34)
Invoking once more the field equation (27), we obtain the integral
$$𝑑\tau \ddot{\mathrm{\Delta }}(\tau )\dot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }(\tau )=\omega ^2𝑑\tau \dot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }^2(\tau ),$$
(35)
where we have used $`\dot{\mathrm{\Delta }}(0)=0`$. Due to Eq. (34), this takes the form
$$𝑑\tau \ddot{\mathrm{\Delta }}(\tau )\dot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }(\tau )=\frac{1}{3}\omega ^2\mathrm{\Delta }^3(0)\frac{1}{3}\omega ^4𝑑\tau \mathrm{\Delta }^4(\tau ).$$
(36)
A further partial integration reduces the integral
$`{\displaystyle 𝑑\tau \dot{\mathrm{\Delta }}^4(\tau )}=3{\displaystyle 𝑑\tau \mathrm{\Delta }(\tau )\dot{\mathrm{\Delta }}^2(\tau )\ddot{\mathrm{\Delta }}(\tau )}`$ (37)
to (36), such that we arrive at the relation
$$𝑑\tau \dot{\mathrm{\Delta }}^4(\tau )=\omega ^2\mathrm{\Delta }^3(0)+\omega ^4𝑑\tau \mathrm{\Delta }^4(\tau ).$$
(38)
We now consider an integral over $`\ddot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }^2`$. Applying again the field equation (27), we find the relation
$`{\displaystyle 𝑑\tau \ddot{\mathrm{\Delta }}^2(\tau )\mathrm{\Delta }^2(\tau )}`$ $`=`$ $`{\displaystyle 𝑑\tau \mathrm{\Delta }^2(\tau )\delta ^2(\tau )}`$ (39)
$``$ $`2\omega ^2\mathrm{\Delta }^3(0)+\omega ^4{\displaystyle 𝑑\tau \mathrm{\Delta }^4(\tau )},`$ (40)
The relations Eqs. (31) and (40) have reduced all integrals over singular products of correlation functions to regular integrals plus two undefined integrals containing $`\delta ^2(\tau )`$. We are now going to show, that the reparametrization invariance of path integrals requires the following rules for integrals over products of two $`\delta `$-functions in Eqs. (31) and (40):
$$𝑑\tau \delta ^2(\tau )=\delta (0),$$
(41)
and further
$$𝑑\tau f(\tau )\delta ^2(\tau )=f(0)\delta (0),$$
(42)
for any smooth test function $`f(\tau )`$.
## V Imposing Reparametrization Invariance
To first order in $`g`$, the sum of Feynman diagrams (18) must vanish:
$$\text{ }\text{}\text{ }+\omega ^2\text{ }\text{}\text{ }\delta (0)\text{ }\text{}\text{ }=0.$$
(43)
The analytic form of this relation is
$$\left[\ddot{\mathrm{\Delta }}(0)+\omega ^2\mathrm{\Delta }(0)\delta (0)\right]\mathrm{\Delta }(0)=0,$$
(44)
and the vanishing is a direct consequence of the field equation (27) for the correlation function at origin.
At order $`g^2`$, the same equation reduces the sum of all local diagrams in (21) to a finite result plus a term proportional to $`\delta (0)`$:
$`[3({\displaystyle \frac{1}{2}}+a)\ddot{\mathrm{\Delta }}(0)+15({\displaystyle \frac{1}{18}}+{\displaystyle \frac{a}{5}})\omega ^2\mathrm{\Delta }(0)`$ (45)
$`3(a{\displaystyle \frac{1}{2}})\delta (0)]\mathrm{\Delta }^2(0)=[3\delta (0){\displaystyle \frac{2}{3}}\omega ^2\mathrm{\Delta }(0)]\mathrm{\Delta }^2(0).`$ (46)
Representing right-hand side diagrammatically, we obtain the identity
$$\Sigma (\text{21})=3\delta (0)\text{ }\text{}\text{ }\frac{2}{3}\omega ^2\text{ }\text{}\text{ },$$
(47)
where $`\Sigma (\text{21})`$ denotes the sum of all diagrams in Eq. (21). Using the identity (28) together with the field equation (27), we reduce the sum (23) of all one and two-loop bubbles diagrams to terms involving $`\delta (0)`$ and $`\delta ^2(0)`$:
$`{\displaystyle \frac{1}{2!}}\{2\delta ^2(0){\displaystyle }d\tau \mathrm{\Delta }^2(\tau )`$ (48)
$`4\delta (0){\displaystyle }d\tau [\mathrm{\Delta }(0)\dot{\mathrm{\Delta }}^2(\tau )\ddot{\mathrm{\Delta }}(0)\mathrm{\Delta }^2(\tau )+2\omega ^2\mathrm{\Delta }(0)\mathrm{\Delta }^2(\tau )]\}`$ (49)
$`=2\mathrm{\Delta }^2(0)\delta (0)+\delta ^2(0){\displaystyle 𝑑\tau \mathrm{\Delta }^2(\tau )}.`$ (50)
Hence we find the diagrammatic identity
$$\frac{1}{2!}\Sigma (\text{23})=2\delta (0)\text{ }\text{}\text{ }+\delta ^2(0)\text{ }\text{}\text{ }.$$
(51)
Now, the terms accompanying $`\delta ^2(0)`$ turn out to be canceled by similar terms coming from the sum of all three-loop bubbles diagrams in (25). In fact, the identities (28) and (31) lead to
$`{\displaystyle \frac{1}{2!}}{\displaystyle }d\tau [4\mathrm{\Delta }(0)\ddot{\mathrm{\Delta }}(0)\dot{\mathrm{\Delta }}^2(\tau )+2\mathrm{\Delta }^2(0)\ddot{\mathrm{\Delta }}^2(\tau )`$ (52)
$`+2\ddot{\mathrm{\Delta }}^2(0)\mathrm{\Delta }^2(\tau )+8\omega ^2\mathrm{\Delta }^2(0)\dot{\mathrm{\Delta }}^2(\tau )`$ (53)
$`8\omega ^2\mathrm{\Delta }(0)\ddot{\mathrm{\Delta }}(0)\mathrm{\Delta }^2(\tau )+8\omega ^4\mathrm{\Delta }^2(0)\mathrm{\Delta }^2(\tau )]`$ (54)
$`=\left[{\displaystyle 𝑑\tau \delta ^2(\tau )}+2\delta (0)\right]\mathrm{\Delta }^2(0)\delta ^2(0){\displaystyle 𝑑\tau \mathrm{\Delta }^2(\tau )}.`$ (55)
Thus, we find the diagrammatic identity for all bubbles diagrams
$$\frac{1}{2!}\Sigma (\text{23})\frac{1}{2!}\mathrm{\Sigma }(\text{25})=𝑑\tau \delta ^2(\tau )\text{ }\text{}\text{ }.$$
(56)
Finally, the relations (34), (40), (36) and (38) reduce the sum (26) of all watermelon-like diagrams to a finite contribution plus the integral involving $`\delta ^2(\tau )`$:
$`{\displaystyle \frac{4}{2!}}{\displaystyle }d\tau [\mathrm{\Delta }^2(\tau )\ddot{\mathrm{\Delta }}^2(\tau )+4\mathrm{\Delta }(\tau )\dot{\mathrm{\Delta }}^2(\tau )\ddot{\mathrm{\Delta }}(\tau )`$ (57)
$`+\dot{\mathrm{\Delta }}^4(\tau )+4\omega ^2\mathrm{\Delta }^2(\tau )\dot{\mathrm{\Delta }}^2(\tau )+{\displaystyle \frac{2}{3}}\omega ^4\mathrm{\Delta }^4(\tau )]`$ (58)
$`=2{\displaystyle 𝑑\tau \mathrm{\Delta }^2(\tau )\delta ^2(\tau )}+{\displaystyle \frac{2}{3}}\omega ^2\mathrm{\Delta }^3(0).`$ (59)
Combining these with all local diagrams (47), we easily verify that all finite contributions cancel each other leading to the diagrammatic identity
$`\Sigma (\text{21}){\displaystyle \frac{4}{2!}}\mathrm{\Sigma }(\text{26})`$ (60)
$`=\left[3\delta (0)2\mathrm{\Delta }^2(0){\displaystyle 𝑑\tau \mathrm{\Delta }^2(\tau )\delta ^2(\tau )}\right]\text{ }\text{}\text{ }.`$ (61)
If the singular terms in Eqs. (56) and (61) are to sum up to zero, as required by the coordinate invariance of perturbatively defined path integrals, we must have the integration rules for the square distribution (41) and (42), which determined completely the right-hand sides of relations (31) and (40).
The procedure can easily be continued to higher-loop diagrams to obtain integrals over any desired products of singular correlation functions, and over products of $`\delta `$-functions.
At no place do we have to specify the value of $`\delta (0)`$ and the regularization scheme. There is a perfect cancellation of all powers of $`\delta (0)`$ arising from the expansion of the Jacobian action, and this is the reason why the so-called Veltman rule of setting $`\delta (0)=0`$ can be used everywhere without problems.
## VI Summary
In this note we have set up simple rules for relating singular to regular Feynman integrals which avoid the explicit calculation of dimensionally regularized integrals over products of distributions. These rules follow directly from the invariance of perturbatively defined path integral under coordinate transformations. Our procedure is independent of of regularization prescriptions, using only the fact that regularized integrals can be integrated by parts. The results are, of course, perfectly compatible with those derived before in Refs. by dimensional regularization.
Just as in the time-sliced definition of path integrals in curved spcae in Ref. , there is absolutely no need for extra compensating potential terms found necessary in the treatments in Refs. . |
warning/0002/astro-ph0002028.html | ar5iv | text | # ChaMP and the High Redshift Quasars in X-rays
## 1. Introduction
NASA’s Chandra X-ray Observatory was launched on July 23, 1999. The Chandra Multiwavelength Project (ChaMP) will combine radio to X-ray observations of serendipitous Chandra sources, with emphasis on optical identification. The ChaMP is superior to previous X-ray surveys because of (1) unprecedented X-ray positional accuracy ($`1^{\prime \prime }`$), (2) X-ray flux limits 20 times deeper than current wide area surveys (down to $`f(0.53.5keV)2\times 10^{16}`$ erg s<sup>-1</sup> cm<sup>-2</sup>), (3) larger sky coverage ($``$ 8 deg<sup>2</sup>) per year than current deep surveys.
## 2. Prediction of Redshift Distribution of Quasars in ChaMP Fields
The X-ray Luminosity Function at $`z=0`$ is described as
$$\mathrm{\Phi }(L_X)=\mathrm{\Phi }_1^{}L_{44}^{\gamma 1}forL<L^{}(0)$$
$$\mathrm{\Phi }(L_X)=\mathrm{\Phi }_2^{}L_{44}^{\gamma 2}forL>L^{}(0)$$
where L<sub>44</sub> is the X-ray luminosity in $`10^{44}`$ erg s<sup>-1</sup>. The redshift evolution of the luminosity function is characterized by
$$L_X(z)=L_X(0)(1+z)^k$$
Continuity of the luminosity function at the break luminosity requires that
$$\mathrm{\Phi }_1^{}=\mathrm{\Phi }_2^{}L_{44}^{(\gamma 1\gamma 2)}$$
The total number $`N`$ of quasars in the sample is obtained by integrating the luminosity function over luminosity and volume, i.e.,
$$N=\mathrm{\Phi }(L_X,z)\mathrm{\Omega }(L_X,z)𝑑V(z)𝑑L_X$$
Here $`\mathrm{\Omega }(L_X,z)`$ is the solid angle covered by the survey as a function of redshift and luminosity. The parameters of the X-ray luminosity function determined by Boyle et al. (1993) are as follows: $`\gamma 1=1.7\pm 0.2`$, $`\gamma 2=3.4\pm 0.1`$, $`\mathrm{log}L^{}(0)=43.84`$, $`\mathrm{\Phi }_1^{}=5.7\times 10^7Mpc^3(10^{44}ergs^1)^{\gamma 11}`$. Following Comastri et al. (1995), we have used k=2.6 and increased the the normalization $`\mathrm{\Phi }_1^{}`$ by 20%.
The X-ray logN-logS Curve: Using the above luminosity function we derived the number density of quasars as a function of observed flux. The luminosity function was integrated over the luminosity range $`10^{42}<L_X<10^{48}`$ erg s<sup>-1</sup> and the redshift range $`0<z<4`$. H$`{}_{0}{}^{}=50`$ and q$`{}_{0}{}^{}=0`$ were assumed throughout. The predicted logN-logS curve is shown in figure 1.
Since the unabsorbed sources dominate at the faint end in the soft X-ray range, and since they are likely to be observed at high redshift, in the present analysis we will concentrate on unabsorbed sources only. The absorbed sources would contribute an additional $``$ 60% (Comastri et al. 1995), making the total number consistent with the extrapolation of the empirical determination of logN-logS (Hasinger et al. 1993). The flux of unabsorbed quasars is given by $`fE^\alpha `$ and in the soft X-ray band, $`\alpha `$ is typically 1.3.
The ChaMP Sky Coverage: The ChaMP Cycle 1 consists of 85 extragalactic fields, —b—$`>20^o`$. From all the Chandra cycle 1 fields we have excluded (1) deep fields of PI survey observations, (2) fields with extended sources & planetary targets, (3) ACIS sub-arrays and continuous clocking modes. See figure 2 for ChaMP sky coverage as a function of flux limit.
Cumulative Number Distribution in ChaMP: Integrating the predicted logN-logS over the ChaMP sky coverage, we obtained the cumulative number distribution of quasars in the ChaMP fields (figure 3). The total number in soft band is expected to be over 1500 for unabsorbed sources and over 2500 total.
Predicted Redshift Distribution: The histogram (figure 4) shows the predicted number distribution of quasars in ChaMP fields. Over 200 quasars will be detected in the redshift range $`3<z<4`$ and over 400 quasars in $`2<z<3`$.
## 3. Comparison with Previous X-ray Surveys
We will be able to determine the X-ray luminosity function and its redshift evolution with unprecedented accuracy.
## References
Boyle, B., Griffiths, R., Shanks, T., Stewart, G., & Georgantopoulus, I. 1993, MNRAS, 260, 49
Comstri, A., Setti, G., Zamorani, G., & Hasinger, G. 1995, A&A, 296, 1
Gioia et al. 1990, ApJS, 72, 567
Hasinger, G. et al. 1993, A&A, 275, 1
It’s my pleasure (SM) to thank A. Comastri for useful discussions. This work is supported in parts by NASA grant NAG5-3249 (LTSA). |
warning/0002/physics0002050.html | ar5iv | text | # Scaling in dynamical Turing pattern formation: density of defects frozen into permanent patterns
## Motivation and summary of results
Long time ago Turing pointed out that simple reaction-diffusion (RD) systems of equations can account for formation of biological patterns. The mainstream of research, as reviewed in Ref., is devoted to RD models in continuous space. The continuum RD-patterns are smooth and nonpermanent. On the other hand, it is an empirical fact that even the nearest neighbor cells can differ sharply in their biological functions and their sets of expressed genes. Moreover many biological patterns are permanent. Even the most primitive viruses, like the much studied bacteriophage $`\lambda `$ , possess genetic switches that discriminate between different developmental pathways and make a once chosen pathway permanent. It is reasonable to assume that cells of higher organisms can also lock their distinctive sets of expressed genes.
The Turing patterns on figures in the review are contaminated with defects. If we insist on pattern permanence, we must accept that patterns are permanent together with their defects. Sometimes, like for the animal coat patterns, permanent defects can provide an animal with its own characteristic life-long but not inheritable ”fingerprints”. In other cases, like formation of vital organ structures, a single defect can be fatal. In this situation it is important to understand better the origin of defects.
In this paper we use a simple toy model which in principle should give a homogeneous Turing pattern. Defects are particularly manifest on such a simple background. The model has two genes $`A`$ and $`B`$. The genes are strong mutual repressors. The strong intracellular mutual inhibition is the factor responsible for pattern permanence. Both genes are activated simultaneously in a given cell when a level of their common activator $`a`$ exceeds its critical value $`a_c`$. Pattern formation in RD models of Ref. was simulated with fixed model parameters. In this paper we turn on the activator level $`a`$ at a finite rate to find a scaling relation between density of defects and the rate. Strong mutual intracellular inhibition stabilizes the pattern together with its defects. We obtain permanent domains of $`A`$-phase and domains of $`B`$-phase divided by sharp cell-size boundaries.
We also show that an inhomogeneous activation of the genes can result in a perfect defect-free homogeneous pattern. At first the activator $`a`$ exceeds $`a_c`$ in a small seed area where, say, the gene $`A`$ is chosen. Then $`a`$ slowly spreads around gradually activating more and more cells. The initial choice of $`A`$ is imposed via intercellular coupling on all newly activated cells. The inhomogeneous activation can be sufficiently characterized by a velocity $`v`$ with which the critical $`a=a_c`$ surface spreads. Thanks to the strong mutual intracellular inhibition there is a nonzero threshold velocity $`v_c`$, such that for $`v<v_c`$ the formation of defects is completely suppressed. In this way the very mutual inhibition which is responsible for stability of defects can be harnessed to get rid of them.
The genetic network that we use in our toy model is functionally equivalent to the genetic toggle switch which was syntetized by the authors of the recent paper . In that paper the network is studied experimentally in a single ”cell”. It would be interesting to generalize the experiment to a ”multicellular” structure.
## The toy model
For the sake of definiteness we take a genetic network with two genes $`A`$ and $`B`$. $`A`$ and $`B`$ are mutual repressors. The network is symmetric under exchange $`AB`$. Expression of both genes is initiated by a common activator $`a`$. Let $`A(t,\stackrel{}{x})`$ and $`B(t,\stackrel{}{x})`$ denote time-dependent protein concentrations in the cell at the position $`\stackrel{}{x}`$. $`\stackrel{}{x}`$ belongs to a discrete square lattice with a lattice constant of $`1`$. Evolution of the protein concentrations is described by the stochastic differential equations
$`\dot{A}(t,\stackrel{}{x})=RS_A(t,\stackrel{}{x})A(t,\stackrel{}{x}),`$ (1)
$`\dot{B}(t,\stackrel{}{x})=RS_B(t,\stackrel{}{x})B(t,\stackrel{}{x}).`$ (2)
The last terms in these equations are responsible for the protein degradation. $`R`$ is a transcription rate. $`S_{A,B}(t,\stackrel{}{x})\{0,1\}`$ are dichotomic stochastic processes. They switch on $`(01)`$ and off $`(10)`$ transcription of a given gene. For simplicity the processes are assumed to have the same constant switch-off rate $`r^{\mathrm{off}}`$. The switch-on rates depend on concentrations
$`r_A^{\mathrm{on}}(t,\stackrel{}{x})=a(t)F\left[WB(t,\stackrel{}{x})+V{\displaystyle \underset{\mathrm{n}.\mathrm{n}.\stackrel{}{y}}{}}A(t,\stackrel{}{y})\right],`$ (3)
$`r_B^{\mathrm{on}}(t,\stackrel{}{x})=a(t)F\left[WA(t,\stackrel{}{x})+V{\displaystyle \underset{\mathrm{n}.\mathrm{n}.\stackrel{}{y}}{}}B(t,\stackrel{}{y})\right].`$ (4)
$`W,V`$ are positive coupling constants, $`a(t)`$ is a concentration of the activator. $`F[z]`$ is a smooth step-like sigmoidal function; the function $`F[z]=10^3\mathrm{exp}(z2.2)/[1+\mathrm{exp}(z2.2)]`$ was used in our numerical simulations. In this model the genes A and B are mutual repressors $`(W>0)`$. There is a ”ferromagnetic” coupling between nearest-neighbor cells $`(V>0)`$; expression of $`A`$ in a given cell enhances expression of $`A`$ in its nearest neighbors.
The model is motivated by a genetic switch between two mutual repressors like the one studied in the phage $`\lambda `$ and in the E. coli switch . The mutual repressors have a common promoter site on DNA. A necessary condition for expression of any of them is a binding of an activator molecule to their promoter site . The concentrations $`A`$ and $`B`$ influence its affinity to the promoter site. The gene expression is intermittent because of binding and unbinding of activator molecules. The nearest-neighbor coupling is possible thanks to signalling through intercellular membrane channels.
In an adiabatic limit, when switching of $`S_{A,B}`$ is much faster than protein expression and degradation, the processes $`S_{A,B}`$ can be replaced by their time averages,
$`\dot{A}(t,\stackrel{}{x})={\displaystyle \frac{Ra(t)F\left[WB(t,\stackrel{}{x})+V_{\mathrm{n}.\mathrm{n}.\stackrel{}{y}}A(t,\stackrel{}{y})\right]}{r^{\mathrm{off}}+a(t)F\left[WB(t,\stackrel{}{x})+V_{\mathrm{n}.\mathrm{n}.\stackrel{}{y}}A(t,\stackrel{}{y})\right]}}A(t,\stackrel{}{x}),`$ (5)
$`\dot{B}(t,\stackrel{}{x})={\displaystyle \frac{Ra(t)F\left[WA(t,\stackrel{}{x})+V_{\mathrm{n}.\mathrm{n}.\stackrel{}{y}}B(t,\stackrel{}{y})\right]}{r^{\mathrm{off}}+a(t)F\left[WA(t,\stackrel{}{x})+V_{\mathrm{n}.\mathrm{n}.\stackrel{}{y}}B(t,\stackrel{}{y})\right]}}B(t,\stackrel{}{x}).`$ (6)
Here we temporarily neglect any noise terms.
## Attractor structure
In a subspace of uniform configurations $`A(t),B(t)`$ these equations simplify to the dynamical system
$`\dot{A}={\displaystyle \frac{RaF\left[WB+2dVA\right]}{r^{\mathrm{off}}+aF\left[WB+2dVA\right]}}A,`$ (7)
$`\dot{B}={\displaystyle \frac{RaF\left[WA+2dVB\right]}{r^{\mathrm{off}}+aF\left[WA+2dVB\right]}}B,`$ (8)
where $`2d`$ is the number of nearest neighbors in $`d`$ dimensions.
The RHS’s of these equations define a velocity field on the $`AB`$ plane, which is not a gradient field. The velocity field has attractor structure which depends on the activator level $`a`$. There are two critical activator levels $`a_{c_1}<a_{c_2}`$. For $`a<a_{c_1}`$ there is one attractor at $`[A,B]=[\gamma (a),\gamma (a)]`$ with an increasing function $`\gamma (a)`$. In the range $`a_{c_1}<a<a_{c_2}`$ there are three attractors: the old $`[\gamma (a),\gamma (a)]`$ plus a new symmetric pair of $`[\alpha (a),\beta (a)]`$ and $`[\beta (a),\alpha (a)]`$ with $`\alpha (a)>\beta (a)`$. For $`a_{c_2}<a`$ there remain only the two broken symmetry attractors $`[\alpha (a),\beta (a)]`$ and $`[\beta (a),\alpha (a)]`$. The functions $`\alpha (a),\beta (a)`$ and $`\gamma (a)`$ are plotted in Fig.1.
If we start in the $`[A,B]=[0,0]`$ state and slowly increase $`a`$-level, the system will stay in the $`\gamma \gamma `$-phase until we reach $`a=a_{c_2}`$. At $`a=a_{c_2}^+`$ the system will roll into $`\alpha \beta `$ or $`\beta \alpha `$-phase. On the other hand, if we start from $`a_{c_2}<a`$ with the system in, say, $`\alpha \beta `$-phase, then we will have to decrease $`a`$ down to $`a=a_{c_1}`$, where $`\alpha \beta `$ becomes unstable towards the symmetric $`\gamma \gamma `$-phase. The discontinuous jumps of the concentrations are illustrated in Fig.1. This hysteresis loop is characteristic for first order phase transitions. In the adiabatic limit, where fluctuations are small, there are no short cuts via bubble nucleation. When $`a_{c_1}`$ ($`a_{c_2}`$) is approached from above (below), the correlation length of small fluctuations around this uniform state diverges like in a continuous phase transition. The critical regime is narrow in the adiabatic limit so we can rely on the mean field approximation.
## A finite rate Turing transition
Let us think again about starting from $`[A,B]=[0,0]`$ and continuously increasing $`a(t)`$ above $`a_{c_2}`$. At $`a_{c_2}^+`$ the $`\gamma \gamma `$ state becomes unstable and the system has to choose between the $`\alpha \beta `$ and $`\beta \alpha `$ attractors. If $`a(t)`$ is increased at a finite rate, then there are finite correlated domains which make the choice independently. Despite divergence of the correlation length at $`a_{c_2}^{}`$, the critical slowing down results in a certain finite correlation length $`\widehat{\xi }`$ ”frozen” into the fluctuations. This scale defines density of defects in the Turing pattern. This effect is well known in cosmology and condensed matter physics as Kibble-Zurek scenario . In those contexts the defects disappear rapidly as a result of phase ordering kinetics. We will see that in our gene network model the defect pattern is permanent. This effect results from a combination of the histeresis loop and the discreteness of the cell lattice.
To be more quantitative we substitute $`A(t,\stackrel{}{x})=\gamma (a(t))+\delta A(t,\stackrel{}{x})`$ and $`B(t,\stackrel{}{x})=\gamma (a(t))+\delta B(t,\stackrel{}{x})`$ into Eqs.(5) and linearize them in $`\delta A,\delta B`$. The linearized equations can be diagonalized by $`\varphi =\delta A\delta B`$ and $`\psi =\delta A+\delta B`$. After Fourier transformation in space
$$\varphi (t,\stackrel{}{x})=d^dk\stackrel{~}{\varphi }(t,\stackrel{}{k})e^{i\stackrel{}{k}\stackrel{}{x}}$$
(9)
they become
$`\dot{\varphi }(t,\stackrel{}{k})=Rs_\varphi (t,\stackrel{}{k})+{\displaystyle \frac{r^{\mathrm{off}}Ra(t)F_a^{}}{[r^{\mathrm{off}}+a(t)F_a]^2}}\left[W\varphi (t,\stackrel{}{k})+Ve_\stackrel{}{k}\varphi (t,\stackrel{}{k})\right]\varphi (t,\stackrel{}{k}),`$ (10)
$`\dot{\psi }(t,\stackrel{}{k})=Rs_\psi (t,\stackrel{}{k})+{\displaystyle \frac{r^{\mathrm{off}}Ra(t)F_a^{}}{[r^{\mathrm{off}}+a(t)F_a]^2}}\left[W\psi (t,\stackrel{}{k})+Ve_\stackrel{}{k}\psi (t,\stackrel{}{k})\right]\psi (t,\stackrel{}{k}),`$ (11)
where $`e_\stackrel{}{k}=2_{i=1}^d\mathrm{cos}k_i`$ in $`d`$ dimensions and we skipped the tildas over Fourier transforms. $`F^{}[z]=dF[z]/dz`$ and we used the shorthands $`F_a^{(^{})}=F^{(^{})}[(W+2dV)\gamma (a(t))]`$. $`Rs_{\varphi ,\psi }`$ are noises which result from fluctuations in $`RS_{A,B}`$. In the adiabatic limit they can be approximated by white noises (both in space and in time) with small magnitude.
The next step is to linearize $`a(t)`$ around its critical value $`a(t)=a_{c_2}+t/\tau `$, where $`\tau `$ is the transition rate. This linearization gives
$$\frac{r^{\mathrm{off}}Ra(t)F_a^{}}{[r^{\mathrm{off}}+a(t)F_a]^2}=c_0+c_1\frac{t}{\tau }+O[(t/\tau )^2].$$
(12)
Approximating $`e_\stackrel{}{k}=2d\stackrel{}{k}^2`$ in Eqs.(10,11) and keeping only leading terms in $`t/\tau `$ and in $`k^2`$ we get
$`\dot{\varphi }(t,\stackrel{}{k})=Rs_\varphi (t,\stackrel{}{k})+\left[({\displaystyle \frac{c_1}{c_0}}){\displaystyle \frac{t}{\tau }}(c_0V)\stackrel{}{k}^2\right]\varphi (t,\stackrel{}{k}),`$ (13)
$`\dot{\psi }(t,\stackrel{}{k})=Rs_\psi (t,\stackrel{}{k})[2c_0W+c_0V\stackrel{}{k}^2]\psi (t,\stackrel{}{k}).`$ (14)
Here we used the identity $`c_0[W+2dV]=1`$, which has to be satisfied because, by definition, $`\varphi (t,\stackrel{}{0})`$ is a zero mode at $`a_{c_2}`$. The $`\psi `$ modes are stable for any $`\stackrel{}{k}`$. The $`\varphi `$-modes in the neighborhood of $`\stackrel{}{k}=\stackrel{}{0}`$ become unstable for $`t>0`$ (or $`a_{c_2}<a`$). Eq.(13) is a standard linearized Landau model with the symmetry breaking parameter $`(c_1/c_0)(t/\tau )`$ changing sign at $`t=0`$. The length scale $`\widehat{\xi }`$ frozen into fluctuations at $`t>0`$ can be estimated following the classic argument given by Zurek . For $`t<0`$ the model (13) has an instantaneous relaxation time $`c_0\tau /c_1|t|`$ and an instantaneous correlation length $`c_0\sqrt{V\tau /c_1|t|}`$. They both diverge at $`t=0^{}`$. The fluctuations can no longer follow the increasing $`a(t)`$ when their relaxation time becomes equal to the time still remaining to the transition at $`a=a_{c_2}`$, $`c_0\tau /c_1|t||t|`$. At this instant the correlation length is
$$\widehat{\xi }\left(\frac{V^{1/2}c_0^{3/4}}{c_1^{1/4}}\right)\tau ^{1/4}.$$
(15)
This scale determines the typical size of the $`\alpha \beta `$\- and $`\beta \alpha `$-domains. The scaling relation $`\widehat{\xi }\tau ^{1/4}`$ was verified by numerical simulations illustrated at figures 2 and 3. The domain structures generated in the simulations turned out to be permanent.
The domain structures are permanent because already at $`a_{c_2}`$ the width of the domain wall interpolating between $`\alpha \beta `$ and $`\beta \alpha `$ is less then the cell size (lattice spacing). The nearest neighbor cells across the wall express different genes. The width (the healing length) is determined by the longest length scale of fluctuations around the $`\alpha \beta `$\- or $`\beta \alpha `$-state. These correlation lengths are plotted in Fig.4. For $`aa_{c_2}`$ they are substantially less than $`1`$. In the adiabatic limit, where the noises are weak, the domain wall cannot evolve because it would have to overcome a prohibitive potential barrier. On a cellular level the barrier originates from the mutual inhibition between $`A`$ and $`B`$ in a single cell. Roughly speaking, much above $`a_{c_1}`$ each cell is locked in its gene expression state and insensitive to its nearest neighbors’ states.
## Inhomogeneous activation
The intracellular mutual inhibition stabilizes the Turing pattern but it also stabilizes the defects frozen into the pattern. With the $`\widehat{\xi }\tau ^{1/4}`$ scaling the number of defects is rather weakly dependent on $`\tau `$. There may be not enough time during morphogenesis to get rid of the defects by simply increasing $`\tau `$. However, it is possible to generate a defect-free pattern by spatially inhomogeneous switching of the activator level $`a`$. For example, its concentration can exceed $`a_{c_2}`$ at one point at first, where the cells happen to pick (or are forced to pick), say, $`\alpha \beta `$-phase, and then the activator can gradually spread around so that the initial seed of $`\alpha \beta `$-cells gradually imposes their choice on the whole system. For continuous transitions this effect was described in Ref.().
The effect of defect suppression in inhomogeneous activation can be most easily studied in a one dimensional version of the model (1). Suppose that a smooth activator front is moving across the one dimensional chain of cells with a velocity $`v`$, $`a(t,x)a_{c_2}+(vtx)/v\tau `$ close to $`x=vt`$ where $`a=a_{c_2}`$. For definiteness we impose two asymptotic conditions: for $`vtx`$ ( where $`a<a_{c_2}`$) the cells are in the $`\gamma \gamma `$-state, and for $`xvt`$ (where $`a>a_{c_2}`$) they are in $`\alpha \beta `$-phase. We can expect that as the $`a`$-front moves to the right it is followed by the $`\alpha \beta `$ front gradually entering the area formerly occupied by the $`\gamma \gamma `$-phase. If the concentration front is fast enough to move in step with the activator front, then the $`\alpha \beta `$-phase will gradually fill the whole system. If, on the other hand, the concentration front is slower than the activator front then the front of the $`\alpha \beta `$-phase will lag behind the $`a=a_{c_2}`$ front. The gap between the two fronts will grow with time. The gap will be filled with the unstable $`\gamma \gamma `$-phase ($`a>a_{c_2}`$ behind the $`a`$-front). When the gap becomes wide enough, then $`\gamma \gamma `$-state will be able to decay towards the $`\beta \alpha `$-state. A domain of $`\beta \alpha `$-phase will eventually be nucleated behind the $`a`$-front. Now the $`\beta \alpha `$-domain will grow behind the $`a`$-front until its front lags sufficiently behind so that a new domain of $`\alpha \beta `$-phase will be nucleated. In this way the activator front will leave behind a landscape of alternating $`\alpha \beta `$\- and $`\beta \alpha `$-domains qualitatively the same as for homogeneous activation.
The success of the inhomogeneous activation depends on the relation between the velocity $`v`$ of the $`a`$-front and that of the concentration front. As illustrated in Fig.4 fluctuations around the $`\alpha \beta `$-state have two families of modes each with a different correlation length. For any $`a`$ each $`\stackrel{}{k}`$-mode within each family has a different diffusion velocity: a ratio of its wavelength to its relaxation time. The lowest of these diffusion velocities, $`v_c(a)`$, is the maximal velocity at which the $`\alpha \beta `$-phase can spread into the area occupied by the $`\gamma \gamma `$-phase. $`v_c(a_{c_2})v_{c_2}>0`$ because at $`a=a_{c_2}`$ the $`\alpha \beta `$-state is stable (the hysteresis loop again!). $`v_c(a)`$ increases with an increasing $`a`$. If $`v<v_{c_2}`$ the $`\alpha \beta `$-front moves in step with the $`a`$-front; its tail spreads into the $`vt<x`$ area imposing an $`\alpha \beta `$-bias on the fluctuations around $`\gamma \gamma `$-state. The $`\alpha \beta `$-phase spreads without nucleation of any $`\beta \alpha `$-domains. For $`v<v_{c_2}`$ a defect-free uniform Turing pattern forms behind the activator front. Results from numerical simulations of the inhomogeneous activation are presented in Fig.5.
## More complicated patterns
Finally, it is time to comment on more complicated models which are expected to give more complicated patterns than the (in principle) uniform pattern discussed so far. Let us pick a zebra pattern for example. For the uniform pattern the first mode to become unstable in Eq.(13) is the $`\stackrel{}{k}=\stackrel{}{0}`$ mode. The final pattern has an admixture of $`\stackrel{}{k}`$’s in a range $`\widehat{\xi }^1`$ around $`\stackrel{}{k}=\stackrel{}{0}`$. In distinction, for the zebra pattern the first unstable modes are those on the circle $`|\stackrel{}{k}|=2\pi /L`$, where $`L`$ is the spacing between zebra stripes. The final pattern has an admixture of $`\stackrel{}{k}`$’s in a ring of thickness $`\widehat{\xi }^1`$ around the circle $`|\stackrel{}{k}|=2\pi /L`$, compare results for Swift-Hohenberg equation in Ref.. This admixture results in defects frozen into zebra pattern. The inhomogeneous activation can be applied in the zebra case too. In addition it can be used to arrange the stripes. An activator spreading from an initial point would result (at least close to the initial point) in concentric black and white rings. A front of activator moving through the system would comb the stripes perpendicular to the front.
Acknowledgements. I would like to thank M.Sadzikowski and W.Zurek for useful comments on the manuscript. |
warning/0002/math0002243.html | ar5iv | text | # Seiberg-Witten invariants of non-simple type and Einstein metrics
## 1 Introduction
A smooth Riemannian manifold $`(M,g)`$ is said to be *Einstein* if its Ricci curvature tensor $`r`$ is a multiple of the metric i.e.
$$r=\lambda g.$$
Not every smooth compact oriented $`4`$-manifold admits such a metric. A well known obstruction is given by the following result due to N. Hitchin and J. Thorpe (see ). *If $`M`$ is a compact oriented $`4`$-manifold and $`e(M)<\frac{3}{2}|\sigma (M)|`$ then $`M`$ does not admit an Einstein metric, where $`e`$ and $`\sigma `$ respectively denote the Euler characteristic and the signature.* The Gauss-Bonnet-like formula
$$2e(M)\pm 3\sigma (M)=\frac{1}{4\pi ^2}_M\left(\frac{s^2}{24}\frac{|r_0|^2}{2}+2|W_\pm |^2\right)𝑑\mu ,$$
implies Hitchin-Thorpe’s inequality because Einstein metrics are characterized by the vanishing of $`r_0`$, and this is the only negative term in the above integrand. Here $`s,r_0,W_+,W_{}`$ respectively denote the scalar, trace-free Ricci, self-dual Weyl, and anti-self-dual Weyl curvature tensors of a Riemannian metric.
As C. LeBrun showed in this result can be improved using careful estimates on the $`L^2`$-norm of the scalar curvature tensor $`s`$ and the $`L^2`$-norm of the self-dual part of the Weyl tensor $`W_+`$ arising from the Seiberg-Witten equations if, for example, the smooth $`4`$-manifold $`M`$, admits a symplectic form. To obtain these estimates C. LeBrun used that such an $`M`$ admits irreducible solutions to the Seiberg-Witten equations for every metric $`g`$ rather than actually using the fact that $`M`$ has non-trivial Seiberg-Witten invariant. Our main result is
###### Theorem A.
Let $`(M,𝔠)`$ be a smooth compact Kähler surface with a $`\mathrm{Spin}^c`$-structure $`𝔠`$. There is a canonical $`\mathrm{Spin}^c`$ structure in the connected sum manifold $`M\mathrm{\#}(S^1\times S^3)`$ which we will denote by $`𝔠_{0,1}`$. Moreover $`d(𝔠_{0,1})=d(𝔠)+1`$. If $`𝔠`$ is a non-trivial SW-class for $`M`$ then $`𝔠_{0,1}`$ is a $`B`$-class for the connected sum $`M\mathrm{\#}(S^1\times S^3)`$.
The equality $`d(𝔠_{0,1})=d(𝔠)+1`$ implies that $`𝔠_{0,1}`$ is not induced by an almost complex structure, and the statement *$`𝔠_{0,1}`$ is a $`B`$-class* implies that there exist irreducible solutions to the Seiberg-Witten equations for every Riemannian metric. The technique that we have used to produce these $`\mathrm{Spin}^c`$-structures does not rely on the well-known gluing-argument (compare with ).
The main application of our result is
###### Theorem B.
For each admissible pair $`(m,n)`$ there exist an infinite number of non-homeomorphic compact oriented $`4`$-manifolds which have Euler characteristic $`m`$, signature $`n`$, with free fundamental group and which do not admit an Einstein metric.
Similar examples but with very complicated fundamental group have been obtained by A. Sambusetti using connected sums with real or complex hyperbolic $`4`$-manifolds.
I would like to thank Prof. C. LeBrun for all the useful comments, and the time he spent reading previous versions of this manuscript.
## 2 SW-Moduli Space
###### Definition 1.
Let $`(M,𝔠)`$ be a smooth compact oriented $`4`$-manifold with a $`\mathrm{Spin}^c`$-structure $`𝔠`$. Let $`L_𝔠=det(𝔠)`$ be the determinant line bundle associated to $`𝔠`$. Fix a Riemannian metric $`g`$ on $`M`$. The configuration space $`𝒞(𝔠)`$ consist of pairs $`(A,\varphi )`$ , where $`A`$ is an $`U(1)`$-connection on $`L_𝔠`$ and $`\varphi 𝒞^{\mathrm{}}(S^+(𝔠))`$ is a self-dual spinor. We say that $`(A,\varphi )`$ satisfy the Seiberg-Witten equations (SW-equations) if and only if
$`D_A\varphi `$ $`=0`$
$`F_A^+`$ $`=q(\varphi ),`$
where $`q(\varphi )=\varphi \varphi ^{}\frac{|\varphi |^2}{2}\text{Id}`$.
###### Remark.
$`D_A`$ is the associated Dirac operator of the $`\mathrm{Spin}^c`$-bundle, and $`F_A^+`$ is the self-dual part of the curvature associated to the connection $`A`$, thought of as an endomorphism of the self-dual spinors.
###### Definition 2.
We say that an element $`(A,\varphi )`$ is irreducible if $`\varphi 0`$, otherwise it is reducible. We denote by $`𝒞^{}(𝔠)`$ the open subset of irreducible configurations, by $`𝒢(𝔠)=\{\sigma :MS^1\}`$ the gauge group, and by $`^{}(𝔠)=𝒞^{}(𝔠)/𝒢(𝔠)`$ the open subset of irreducible equivalence classes.
The naive definition of the Seiberg-Witten moduli space would be:
$$_g(𝔠)=\{(A,\varphi )𝒞(𝔠)|D_A\varphi =0,F_A^+=q(\varphi )\}/𝒢(𝔠),$$
but in order to use the usual analytical tools, one has to extend the $`𝒞^{\mathrm{}}`$ objects to appropriate Sobolev spaces. From now on we extend the configuration space $`𝒜(𝔠)`$ and the gauge group $`𝒢(𝔠)`$ by requiring $`A`$ and $`\varphi `$ to be in $`L_2^2`$ and $`\sigma `$ to be in $`L_3^2`$. The SW-equations and the gauge actions make sense in this context also and we define:
###### Definition 3.
The Seiberg-Witten moduli space is:
$$_g(𝔠)=\{(A,\varphi )𝒞(𝔠)|D_A\varphi =0,F_A^+=q(\varphi )\}/𝒢(𝔠),$$
where $`𝒜(𝔠)`$ and $`𝒢(𝔠)`$ are the extended configuration space and gauge group. The formal dimension (computed using the Atiyah-Singer index theorem) of this moduli space is
$$d(𝔠)=\frac{c_1^2(𝔠)(2e(M)+3\sigma (M))}{4}.$$
In general there is no reason to expect that the moduli space form a smooth manifold. The best we can hope for is that *generically* it does. The next Theorem guarantees that this is the case. For the proof see .
###### Theorem 1.
Suppose that $`b_2^+>0`$. Fix a metric $`g`$ on $`M`$. Then for a generic $`𝒞^{\mathrm{}}`$ self-dual $`2`$-form $`h`$ on $`M`$ the following holds. For any $`\mathrm{Spin}^c`$-structure $`𝔠`$ on $`M`$ the moduli space $`_g(𝔠,h)(𝔠)`$ of gauge equivalence classes of pairs $`[A,\varphi ]`$ which are solutions to the perturbed SW-equations
$`D_A\varphi `$ $`=0`$
$`F_A^+q(\varphi )`$ $`=ih`$
form a smooth compact submanifold of $`^{}(𝔠)`$ of dimension $`d(𝔠)`$.
Also in it is shown that if $`b_2^+>1`$ then the bordism class of $`_g(𝔠,h)`$ is an invariant of the smooth structure of $`M`$ and the $`\mathrm{Spin}^c`$-structure $`𝔠`$ on $`M`$. We will denote by $`(𝔠)`$ this bordism class.
###### Proposition 2.
Consider a fixed $`U(1)`$-connection $`A`$ on $`L_𝔠`$. Let $`[A_i,\varphi _i]`$ be solutions to the SW-equations, and let $`(A_i,\varphi _i)`$ be the unique representatives such that $`A_iA`$ is co-closed (gauge fixing condition, see ), for $`i=1,2`$. If $`\varphi _1=\varphi _2`$ then $`A_1=A_2`$.
###### Proof.
The first thing to notice is that $`A_2=A_1+\theta `$, where $`\theta `$ is a co-closed $`1`$-form. Since $`(A_1,\varphi _1)`$ and $`(A_2,\varphi _2)`$ are solutions to the SW-equations we have
$`F_{A_1}^+`$ $`=q(\varphi _1)`$
$`=q(\varphi _2)`$
$`=F_{A_2}^+.`$
Therefore
$`F_{A_2}^+F_{A_1}^+=0`$ $`(d\theta )^+=0`$
$`d\theta =d\theta `$
$`dd\theta =dd\theta =0`$
$`dd\theta =0`$
$`\delta d\theta =0.`$
This last statement and the fact $`\delta \theta =0`$ implies that
$`\mathrm{\Delta }\theta `$ $`=d\delta \theta +\delta d\theta `$
$`=\delta d\theta `$
$`=0.`$
Since $`(A_i,\varphi _i)`$ $`i=1,2`$ are solutions to the Seiberg-Witten equations we have
$`0`$ $`=D_{A_2}\varphi _2`$
$`=D_{A_1+\theta }\varphi _1`$
$`=D_{A_1}\varphi _1+\theta \varphi _1`$
$`=\theta \varphi _1,`$
*multiplying* by $`\theta `$ both sides of the equality we get that $`|\theta |^2\varphi _1=0`$. Taking the point-wise norm we will have $`|\theta |^2|\varphi _1|=0`$. If we denote by $`Z_{|\theta |^2}`$ and $`Z_{|\varphi _1|}`$ the set of points where $`|\theta |^2`$ and $`|\varphi _1|`$ vanish respectively, and we denote by $`Z_{|\theta |^2}^c`$ and $`Z_{|\varphi _1|}^c`$ their corresponding complements, we will have that $`Z_{|\varphi _1|}^cZ_{|\theta |^2}`$, therefore if $`[A_1,\varphi _1]`$ is not a reducible solution then $`Z_{|\varphi _1|}^c`$ is a non-empty open set. By a result of N. Aronszajn (see ) we will have that $`\theta =0`$, since it vanishes in an open set. ∎
Since $`𝒞(𝔠)`$ is an affine space it is contractible. Also the space of reducible configurations $`𝒜(𝔠)\times \{0\}`$ is contractible and has infinite codimension in $`𝒞(𝔠)`$. Since $`𝒞^{}(𝔠)`$ is open in $`𝒞(𝔠)`$ and it is the complement of $`𝒜(𝔠)\times \{0\}`$ then it is contractible. $`^{}(𝔠)=𝒞^{}(𝔠)/𝒢(𝔠)`$ is the classifying space of $`𝒢(𝔠)=Map(M,S^1)`$ since $`𝒢(𝔠)`$ acts freely on $`𝒞^{}(𝔠)`$.
Moreover,
$$Map(M,S^1)Map(M,S^1)_o\times \pi _0(Map(M,S^1)),$$
where $`Map(M,S^1)_o`$ denotes homotopically constant maps. $`Map(M,S^1)_o`$ can be identified with $`S^1`$, therefore $`Map(M,S^1)S^1\times H^1(M;)`$, so the classifying space for $`Map(M,S^1)`$ is weakly homotopically equivalent to $`^{\mathrm{}}\times \frac{H^1(M;)}{H^1(M^{})}`$, and
$$H^{}(^{}(𝔠);)[U]\mathrm{\Omega }^{}H^1(M;),$$
(1)
where $`U`$ is a generator for $`H^{}(^{\mathrm{}};)`$.
###### Definition 4.
The Seiberg-Witten invariant $`SW(𝔠)`$ for the $`\mathrm{Spin}^c`$-structure $`𝔠`$ is defined as follows
$$SW(𝔠)=\{\begin{array}{cc}U^{d(𝔠)/2},(𝔠|_{^{}(𝔠)}\hfill & \text{if }d(𝔠)\text{ is even}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
It is easy to see that this invariant is a cobordism invariant of the moduli space $`(𝔠)`$, therefore it does not depend on the metric we used to define the Dirac operator, it does define an invariant of the smooth manifold $`M`$.
From this definition it is easy to see that we are loosing information about the moduli space. For example if the moduli space is odd dimensional this invariant is zero, even though the moduli itself may not represent a trivial bordism class in $`^{}(𝔠)`$.
###### Definition 5.
Let $`(M,𝔠)`$ be a smooth compact oriented $`4`$-manifold with a $`\mathrm{Spin}^c`$-structure $`𝔠`$. We will say that $`𝔠`$ is a *$`B`$-class* if for some (then for any) Riemannian metric $`g`$ on $`M`$, the moduli space $`_g(𝔠)`$ of irreducible solutions to the SW-equations is a smooth manifold of dimension $`d(𝔠)0`$ that represents a non-trivial bordism class in $`^{}(𝔠)`$, i.e. there exists $`\eta H^{}(^{}(𝔠);)`$ of degree $`d(𝔠)`$ such that
$$\eta ,(𝔠)|_{^{}(𝔠)}0.$$
## 3 SW-Equations and Conformal Structures
It is easy to see that conformal changes on the metric can be lifted to a fixed $`\mathrm{Spin}^c`$-structure, and one can study the associated change in the Dirac operator. A basic important fact is that *the Dirac operator remains essentially invariant under all conformal changes of the metric*.
We now make this statement precise. Let $`(M,𝔠)`$ be a fixed smooth compact oriented $`n`$-manifold with a fixed $`\mathrm{Spin}^c`$-structure $`𝔠`$ and a fixed Hermitian structure $`h`$ on the determinant line bundle $`L_𝔠`$. Fix a Riemannian metric $`g`$ on $`M`$ and consider the conformally related metric $`g_f=e^{2f}g`$, where $`f`$ is a smooth function on $`M`$. To each $`g`$-orthonormal tangent frame $`\{e_i\}_{i=1\mathrm{}n}`$ we can associate the $`g_f`$-orthonormal frame $`\{e_i^{}\}_{i=1\mathrm{}n}`$, where $`e_i^{}=\psi _f(e_i)=e^fe_i`$ for each $`i`$. This map induces a bundle isometry between the bundles $`S(𝔠)`$ and $`S^{}(𝔠)`$. Let $`\mathrm{\Psi }_f=e^{\frac{n1}{2}f}\psi _f`$. The resulting map is a bundle isomorphism which is conformal on each fiber. The proof of the following proposition is similar to the one found in , pages $`132134`$, since we are not changing the $`U(1)`$-connection on $`L_𝔠`$.
###### Proposition 3.
Let $`D_A`$ and $`D_A^{}`$ be the Dirac operators (induced by the $`U(1)`$-connection $`A`$) defined over the conformally related Riemannian manifolds $`(M,g)`$ and $`(M,g_f)`$ respectively. Then
$$\mathrm{\Psi }_fD_A=D_A^{}\mathrm{\Psi }_f$$
###### Corollary 4.
There is bijection between $`\mathrm{ker}D_A`$ and $`\mathrm{ker}D_A^{}`$.
Let $`(M,𝔠)`$ be a fixed smooth compact oriented $`4`$-manifold with a fixed $`\mathrm{Spin}^c`$-structure $`𝔠`$. We want to relate the moduli spaces $`(𝔠)`$ and $`^{}(𝔠)`$ for two Riemannian metrics $`g`$ and $`g_f`$ (respectively) in the same conformal class. It is well known (see ) that both moduli spaces represent the same bordism class (in $`^{}(𝔠)`$), but when one of the metrics is Kähler, both moduli spaces are diffeomorphic (see proposition 6) .
###### Proposition 5.
Let $`(M,𝔠)`$ be a fixed smooth compact oriented $`4`$-manifold with a fixed $`\mathrm{Spin}^c`$-structure $`𝔠`$. Let $`g`$ be a fixed Riemannian metric on $`M`$ and consider the conformal metric $`g_f=e^{2f}g`$. Solutions to the Seiberg-Witten equation for the metric $`g_f`$ are in one-to-one correspondence with solutions of the following pair of equations:
$$\begin{array}{cc}\hfill D_A\varphi & =0\hfill \\ \hfill F_A^+& =e^fq(\varphi ).\hfill \end{array}$$
($`SW_f`$)
The one-to-one correspondence is given by the map $`(A,\varphi )(A,\mathrm{\Psi }_f\varphi )`$.
###### Proof.
This is a consequence of Proposition 3, the expression for $`q`$ (see Definition 1) and that $`^{}|_^2=|_^2`$, where $``$ and $`^{}`$ are the Hodge operators of $`g`$ and $`g_f`$, respectively. ∎
###### Proposition 6.
Let $`(M,g)`$ be a Kähler surface with Kähler metric $`g`$. Then for any smooth function $`f:M`$
* If the degree of $`K_M`$ is negative the only solutions to ($`SW_f`$) are reducible, i.e. $`_{e^{2f}g}(𝔠)=\mathrm{}`$.
* Let $`𝔠`$ be the $`\mathrm{Spin}^c`$-structure determined by the complex structure. If the degree of $`K_M`$ is positive then $`\mathrm{\#}_{e^{2f}g}(𝔠)=1`$.
###### Proof.
The proof of this proposition can be carry out following the steps in the proof of Proposition 7.3.1 in pg. 119, replacing the expression for $`q`$ with $`e^fq`$. ∎
###### Remark.
Note that $`\mathrm{\#}_{e^{2f}g}(𝔠)=1`$ is stronger than $`SW_{e^{2f}g}(𝔠)=1`$, which we already knew (see ).
## 4 SW-Moduli Space of a Manifold with a Cylindrical End
The last result shows that if $`(M,g)`$ is a Kähler surface with $`\mathrm{deg}(K_M)>0`$ the Seiberg-Witten moduli space for any metric $`g_f=e^{2f}g`$ in the same conformal class of $`g`$ consists of a single point. In this Section we extend this result to a manifold with finitely many cylindrical ends.
###### Definition 6.
We will say that $`(M_{\mathrm{}},g_{\mathrm{}})`$ is a *manifold with a cylindrical end modeled on $`^+\times S^3`$*, if $`M_{\mathrm{}}`$ is diffeomorphic to $`M\{p\}`$ where $`M`$ is a closed manifold, and $`F:U_p\{p\}^+\times S^3`$ where $`F(x)=(\mathrm{log}(|x|^1),x/|x|)`$ is a diffeomorphism such that $`(g_{\mathrm{}})|_{U_p\{p\}}`$ is the $`F`$-pull-back of the standard product metric $`dt^2+g_{S^3}`$ on $`^+\times S^3`$ and $`U_p`$ is a neighborhood of $`p`$.
If $`(M,g)`$ is a Riemannian manifold such that $`g`$ is flat in a $`\delta `$-neighborhood of $`p`$, where $`\delta <\mathrm{inj}(M,g)`$, there is a canonical way to produce a manifold with a cylindrical end using the conformal class of $`g`$. Here $`\mathrm{inj}(M,g)`$ denotes the injectivity radius of $`(M,g)`$. Choose a function $`\lambda _l:(0,1][1,\mathrm{})`$ which satisfies
$$\lambda _l(r)=\{\begin{array}{cc}1\hfill & \text{if }0re^l\delta ^3\hfill \\ \delta ^2/r\hfill & \text{if }e^l\delta ^2r\delta ^2\hfill \\ 1\hfill & \text{if }r\delta .\hfill \end{array}$$
(2)
Consider the sequence of functions $`\{f_l\}`$, where $`e^{f_l(x)}=\lambda _l(|x|)`$ and the sequence of metrics $`g_l=e^{2f_l}g`$. This sequence of metrics converges in the $`𝒞𝒪`$-topology on $`M\{p\}`$ to a metric $`g_{\mathrm{}}`$. The pair $`(M\{p\},g_{\mathrm{}})`$ is a manifold with a cylindrical end. We will denote by $`\mathrm{\Psi }_l`$ the associated conformal isomorphism defined above proposition 3.
The SW-equations make perfectly good sense on a manifold with a cylindrical end, but in order to use the usual analytical tools, one has to extend the $`𝒞^{\mathrm{}}`$ objects to appropriate weighted Sobolev spaces (see ). From now on every time we work on a manifold with finitely many cylindrical ends we extend the configuration space $`𝒜(𝔠)`$ and the gauge group $`𝒢(𝔠)`$ by requiring $`A`$ and $`\varphi `$ to be in $`L_{2,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$ and $`\sigma `$ to be in $`L_{3,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$. The $`L_{q,ϵ}^p(M_{\mathrm{}},g_{\mathrm{}})`$ norm is defined as
$$h_{p,q,ϵ}=e^{\stackrel{~}{ϵ}t}h_{p,q},$$
where $`\stackrel{~}{ϵ}`$ is a smooth non-decreasing function with bounded derivatives, $`\stackrel{~}{ϵ}:M[0,ϵ]`$, such that $`\stackrel{~}{ϵ}(x)0`$ for $`xB_\delta (p)`$ and $`\stackrel{~}{ϵ}(x)ϵ>0`$ for $`xB_{\delta ^2}(p)`$.
Here we choose the weight $`ϵ<1`$ because we want to produce solutions on the manifold with cylindrical end from solutions on the manifold $`(M,g)`$ via the conformal process ($`g_lg_{\mathrm{}}`$) using proposition 5.
###### Proposition 7.
Let $`(M,g)`$ be any Riemannian $`4`$-manifold, where $`g`$ is flat in the neighborhood of some point $`pM`$ . If $`(A,\varphi )`$ is a solution of ($`SW_f`$) on $`(M,g)`$ (where $`f=f_{\mathrm{}}`$) then $`(A,\mathrm{\Psi }_{\mathrm{}}\varphi )`$ is a solution of the SW-equations on $`(M_{\mathrm{}},g_{\mathrm{}})`$, such that $`(A,\mathrm{\Psi }_{\mathrm{}}\varphi )L_{1,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$.
###### Proof.
The fact that $`(A,\mathrm{\Psi }_{\mathrm{}}\varphi )`$ satisfies the SW-equations follows from proposition 5. We just need to show that $`(A,\mathrm{\Psi }_{\mathrm{}}\varphi )L_{1,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$. In order to do this, we will use the metric $`g`$ as the background metric.
$`\mathrm{\Psi }_{\mathrm{}}\varphi _{2,1,ϵ}^2`$ $`=e^{\stackrel{~}{ϵ}t}\mathrm{\Psi }_{\mathrm{}}\varphi _{2,1}^2`$
$`{\displaystyle _{MB_\delta (p)}}(|\varphi |^2+|\varphi |^2)𝑑\mu +`$
$`{\displaystyle _{^+\times S^3}}(|e^{\stackrel{~}{ϵ}t}\mathrm{\Psi }_{\mathrm{}}\varphi |_{\mathrm{}}^2+|e^{\stackrel{~}{ϵ}t}_t\mathrm{\Psi }_{\mathrm{}}\varphi |_{\mathrm{}}^2)𝑑t𝑑\mu _{S^3}`$
$`={\displaystyle _{MB_\delta (p)}}(|\varphi |^2+|\varphi |^2)𝑑\mu +`$
$`{\displaystyle _{B_\delta (p)\{p\}}}(|r^{ϵ+3/2}\varphi |^2+|r^{ϵ+1+3/2}_r\varphi |^2){\displaystyle \frac{1}{r}}𝑑r𝑑\mu _{S^3}`$
$`={\displaystyle _{MB_\delta (p)}}(|\varphi |^2+|\varphi |^2)𝑑\mu +`$
$`{\displaystyle _{B_\delta (p)\{p\}}}r^{2ϵ1}(|\varphi |^2+|r_r\varphi |^2)r^3𝑑r𝑑\mu _{S^3}`$
$`C\varphi _{2,1}^2.`$
To prove that $`AL_{1,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$ we need to recall that
$$_{g_f}|_^p=e^{(n2p)f}_g|_^p$$
where $`g_f=e^{2f}g`$. The computation is very similar to the one above. ∎
Our next task is to show that there is no loss of generality in assuming that a Kähler metric $`g`$ is flat in a neighborhood of some point.
###### Proposition 8.
Let $`(M^{2n},g)`$ be a Kähler $`2n`$-manifold with Kähler metric $`g`$ and induced Kähler form $`\omega `$ . There is no local obstruction to finding a Kähler metric on $`M`$, flat in a neighborhood of a point (a finite collection of points) without changing the Kähler class of $`\omega `$.
###### Proof.
Let $`pM`$. The existence of such metric is equivalent to finding a neighborhood $`U`$ of $`p`$, and a Kähler form $`\omega ^{}`$ in the same Kähler class of $`\omega `$, such that $`\omega ^{}|_U=\omega _0=_{i=1}^ndz^id\overline{z^i}`$. It is well known that there exist an $`ϵ`$-neighborhood $`U_p`$ of $`p`$ and a function $`f:U_p`$ such that $`\omega |_{U_p}=i\overline{}(z\overline{z}+f(z))>0`$, where $`|f(z)|o(|z|^4)`$ and $`|z|`$ denotes the distance (using the Kähler metric $`g`$) on $`U_p`$ to $`p`$. Let $`𝒦^{\mathrm{}}(f)`$ be the space of smooth functions on $`M`$ that satisfy
$$𝒦^{\mathrm{}}(f)=\{h_{s,t}𝒞^{\mathrm{}}(M)|h(z)=f(z)\text{ if }|z|<s,h(z)=0\text{ if }t<|z|\}$$
where $`0<s<tϵ`$, depend on $`h`$. Observe that if $`f`$ is zero we do not have anything to prove, otherwise $`0𝒦^{\mathrm{}}(f)`$, but $`0𝒦^{3+\alpha }(f)`$, where $`𝒦^{3+\alpha }(f)`$ denotes the completion of $`𝒦^{\mathrm{}}(f)`$ in the $`𝒞^{3+\alpha }`$ topology. To see this consider the one-parameter family of functions $`h_k(z)=\rho (k|z|)f(z)`$, where $`\rho `$ is a smooth bump function such that
$$\rho (r)=\{\begin{array}{cc}1\hfill & \text{if }0<r<1/2\hfill \\ 0\hfill & \text{if }1/2<r<1.\hfill \end{array}$$
All these functions are in $`𝒦^{\mathrm{}}(f)`$ and satisfy
$`|h_k(z)|`$ $`o(|z|^4)`$
$`|h_k(z)|`$ $`o(|z|^3)`$
$`|^2h_k(z)|`$ $`o(|z|^2)`$
$`|^3h_k(z)|`$ $`o(|z|)`$
$`|^4h_k(z)|`$ $`o(1).`$
It is not difficult to see that $`h_k0`$ in the $`𝒞^{3+\alpha }`$ topology. It is important to recall that the set $`𝒫(\omega )`$ of smooth functions $`h`$ such that $`\omega _h=\omega +i\overline{}h>0`$, is open in the $`𝒞^{\mathrm{}}`$ topology. This two facts allow us to find $`h_{s,t}𝒦^{\mathrm{}}(f)𝒫(\omega )`$, $`𝒞^{3+\alpha }`$ close to $`0`$, such that
$`\omega _{h_{s,t}}`$ $`=\omega +i\overline{}h_{s,t}>0`$
$`=\omega _0+i\overline{}(f+h_{s,t}),`$
therefore we have
$$\omega _{h_{s,t}}|_{B_s(p)}=\omega _0,$$
where $`B_s(p)=\{zU_p||z|<s\}`$. ∎
###### Corollary 9.
For any compact oriented Kähler surface $`(M,g)`$ with canonical line bundle $`K_M`$ of positive degree, where $`g`$ is flat in a neighborhood of some point, the *induced* manifold with a cylindrical end $`(M_{\mathrm{}},g_{\mathrm{}})`$ admits solutions to the SW-equations.
In order to prove that the Seiberg-Witten moduli space of a manifold with a cylindrical end consists of only one point if $`\mathrm{deg}(K_M)>0`$, we will need the following technical result.
###### Proposition 10.
Let $`(M_{\mathrm{}},g_{\mathrm{}})`$ be a $`4`$-manifold with a cylindrical end. If
$$(A_{\mathrm{}},\varphi _{\mathrm{}})𝒞^{\mathrm{}}L_{k,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})$$
is a solution of the SW-equations on the manifold with cylindrical end $`(M_{\mathrm{}},g_{\mathrm{}})`$, then $`(A_{\mathrm{}},\mathrm{\Psi }_{\mathrm{}}^1\varphi _{\mathrm{}})`$ extends to a smooth solution of ($`SW_f`$) on $`(M,g)`$, replacing the strictly positive function $`e^f`$ by the non-negative function
$$\lambda _{\mathrm{}}(x)=\{\begin{array}{cc}|x|/\delta ^2\hfill & \text{if }|x|<\delta ^2\hfill \\ 1\hfill & \text{if }|x|>\delta \hfill \end{array}$$
###### Proof.
It is easy to see that $`(A_{\mathrm{}},\mathrm{\Psi }_{\mathrm{}}^1\varphi _{\mathrm{}})L^2(M,g)`$, as it is to see that $`(A_{\mathrm{}},\mathrm{\Psi }_{\mathrm{}}\varphi _{\mathrm{}})`$ is a solution of ($`SW_f`$) with function $`\lambda _{\mathrm{}}`$ replacing $`e^f`$. The first equation in ($`SW_f`$) tell us that $`\mathrm{\Psi }_{\mathrm{}}^1\varphi _{\mathrm{}}`$ is a holomorphic section on $`M\{p\}`$. Using Hartog’s Theorem we can extend this to a holomorphic section on $`M`$. All the analysis done in proving proposition 6 can be carry out if we replace the strictly positive function $`e^f`$ in ($`SW_f`$) by a non-negative function $`\lambda _{\mathrm{}}`$ whose zero set has measure zero. ∎
###### Corollary 11.
Let $`(M,g)`$ be a compact oriented Kähler surface with canonical line bundle $`K_M`$ of positive degree, where $`g`$ is flat in a neighborhood of some point. Then there exists a solution $`(A_{\mathrm{}},\varphi _{\mathrm{}})𝒞^{\mathrm{}}L_{k,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$ of the SW-equations on $`(M_{\mathrm{}},g_{\mathrm{}})`$. This solution is unique up to gauge equivalence.
###### Proof.
Since all the analysis done in proving proposition 6 can be carry out if we replace the strictly positive function $`e^f`$ in ($`SW_f`$) by a non-negative function $`\lambda _{\mathrm{}}`$ whose zero set has measure zero, existence is a consequence of corollary 9 and uniqueness is obtained using proposition 10 and proposition 6. ∎
## 5 Holonomy, Connected Sums with $`S^1\times S^3`$ and SW-Invariants
Consider the diffeomorphism
$$F:^4\{0\}\times S^3,F(x)=(\mathrm{log}|x|,\frac{x}{|x|}).$$
It is easy to see that the pull-back of the standard product metric $`g`$ on $`\times S^3`$ under this diffeomorphism is given by
$$F^{}g(\xi ,\eta )=\frac{1}{|x|^2}\xi ,\eta $$
for $`|x|1`$. Fix $`\delta >0`$ and choose a function $`\lambda _l:(0,1][1,\mathrm{})`$ as in (2) and consider the metric
$$g_l(\xi ,\eta )=\lambda _l(|x|)^2\xi ,\eta .$$
Note that for $`e^l\delta ^2|x|\delta ^2`$ this metric agrees with the above pull-back metric $`F^{}g`$.
It is convenient to think of the connected sum $`M\mathrm{\#}(S^1\times S^3)`$ as follows. Let $`M`$ be a smooth compact oriented $`4`$-manifold. Fix two points $`p_1,p_2M`$, and choose a metric $`g`$ on $`M`$ which is flat in a $`\delta `$-neighborhood of $`p_i`$. For every $`l`$ consider the $`e^{l1}\delta ^2`$-neighborhood of $`p_i`$ (with respect to $`g`$) $`B_{p_i}(e^{l1}\delta ^2)`$, and denote by $`M_l`$ the open subset of $`M`$ given by the complement of $`\overline{B_{p_1}(e^{l1}\delta ^2)}\overline{B_{p_2}(e^{l1}\delta ^2)}`$. If we denote by $`T_i=T_i(e^l\delta ^2,e^{l1}\delta ^2)`$ the annulus centered at $`p_i`$ with radii $`e^{l1}\delta ^2`$ and $`e^l\delta ^2`$, it is easy to see that there exist a diffeomorphism (orientation reversing) that takes $`T_1`$ into $`T_2`$ and if we define $`g_l=\lambda _l^2g`$, such diffeomorphism becomes a $`g_l`$-isometry. Since we have observed that $`T_1`$ and $`T_2`$ are $`g_l`$-isometric we can identify $`T_1`$ with $`T_2`$ , and call them $`T_l`$, to obtain a Riemannian manifold $`(M\mathrm{\#}_l(S^1\times S^3),g_l)`$. This manifold is simply the manifold $`M`$ with two cylindrical ends of length $`l`$ obtained by conformally rescaling the metric $`g`$ and identifying the annuli. It is easy to see that such manifold is diffeomorphic to the connected sum $`M\mathrm{\#}(S^1\times S^3)`$.
Even though the process above described can be realized on any smooth $`4`$-manifold the following results are only valid when $`M`$ is a Kähler surface, because to prove them, we (strongly) use that on a given conformal class of metrics, the moduli spaces of solutions of the SW-equations for any two representatives are diffeomorphic, and this was proved for Kähler surfaces (see proposition 6).
Our next task is to explain how a $`\mathrm{Spin}^c`$-structure on $`M`$ transforms into a $`\mathrm{Spin}^c`$-structure on $`M\mathrm{\#}(S^1\times S^3)`$ under the process above described. The following Proposition will be very useful to explain it.
###### Proposition 12.
There is a canonical projection map $`\pi :M\mathrm{\#}(S^1\times S^3)M`$. It has the following properties:
1. The induced maps in cohomology
$$\pi ^{}:H^i(M;𝔽)H^i(M\mathrm{\#}(S^1\times S^3);𝔽)$$
are injective. Here $`𝔽=_2`$ or $``$. Moreover for $`i=0,2,4`$, $`\pi ^{}`$ is an isomorphism.
2. $`\pi ^{}(w_2(M))=w_2(M\mathrm{\#}(S^1\times S^3))`$.
We will denote the $`\mathrm{Spin}^c`$-structure obtained in the above proposition by $`𝔠_{0,1}`$. It is not difficult to show that the formal dimension of the moduli space associated to $`𝔠_{0,1}`$ is $`d(𝔠_{0,1})=d(𝔠)+1`$.
To explain the increment in the dimension above we need to recall the concept of *holonomy*. Let $`P_GM`$ be a principal $`G`$-bundle over $`M`$, with a connection $`A`$. Let $`xM`$ and denote by $`C(x)`$ the loop space at $`x`$. For each $`\gamma C(x)`$ the parallel displacement along $`\gamma `$ is an isomorphism of the fiber $`G`$ onto itself and we will denote it by $`\mathrm{hol}_\gamma (A)`$. The set of all such isomorphisms forms a group, the *holonomy group of $`A`$ with reference point $`x`$*.
Once and for all for each $`l>0`$ we will choose $`p_lT_1`$, $`q_lT_2`$ and a path $`\mathrm{\Gamma }_l:IM`$ from $`p_l`$ to $`q_l`$ such that after identifying $`T_1`$ with $`T_2`$ we obtain and embedding $`\gamma _l:S^1M\mathrm{\#}_l(S^1\times S^3)`$. It is not difficult to observe that for all $`l>0`$ $`[\gamma _l]0\pi _1(M\mathrm{\#}(S^1\times S^3))`$, and in fact $`\gamma _l`$ represents the $`S^1`$ factor of the connected sum.
If $`A`$ is a $`U(1)`$-connection on the determinant line bundle $`L_𝔠`$, we can trivialize $`L_𝔠`$ along $`\mathrm{\Gamma }_l`$ so that the parallel transport along $`\mathrm{\Gamma }_l`$ induces the identity from the fiber at $`p_l`$ to the fiber at $`q_l`$. When we identify $`T_1`$ with $`T_2`$ we still have the extra degree of freedom of how to identify the fiber at $`p_l`$ with the fiber at $`q_l`$, and this is measured by $`\mathrm{hol}_\gamma (A)`$, where $`A`$ is the *glued* connection. If we change of gauge, $`\mathrm{hol}_\gamma (A)`$ remains unchanged because the structure group $`U(1)`$ is Abelian. In this section we will prove that when $`M`$ is a Kähler surface then every solution to the Seiberg-Witten equations for a $`\mathrm{Spin}^c`$-structure $`𝔠`$, induces an $`S^1`$ family of solutions to the SW-equations for the $`\mathrm{Spin}^c`$-structure $`𝔠_{0,1}`$ on $`M\mathrm{\#}(S^1\times S^3)`$.
We can *glue* a solution $`(A_{\mathrm{}},\varphi _{\mathrm{}})`$ of the SW-equations on $`(M_{\mathrm{}},g_{\mathrm{}})`$ to produce a solution $`(A_l,\varphi _l)`$ of the following set of equations on $`(M\mathrm{\#}_l(S^1\times S^3),g_l)`$
$`D_{A_l}\varphi _l`$ $`=\mu (A_l,\varphi _l)=\mu _l`$
$`F_{A_l}^+q(\varphi _l)`$ $`=\nu (A_l,\varphi _l)=\nu _l,`$
where $`(\mu _l,\nu _l)𝒮(𝔠)\times \mathrm{\Omega }_+^2(M\mathrm{\#}_l(S^1\times S^3);i)`$. It is not difficult to see that
$`(\mu _l,\nu _l)L_1^2(M\mathrm{\#}_l(S^1\times S^3),g_l)`$
$`\underset{l\mathrm{}}{lim}(\mu _l,\nu _l)_{2,1}=0,`$
###### Definition 7.
We will denote by $`_\theta (𝔠_{0,1})(𝔠_{0,1})`$ the solution subspace of the SW-equations satisfying the extra condition
$$\mathrm{hol}_\gamma (A)=\theta ,$$
and by $`SW_\theta (𝔠_{0,1})`$ the cobordism invariant associated to this moduli space (counting solutions with appropriate sign). Note that the condition $`\mathrm{hol}_\gamma (A)=\theta `$ reduces the dimension of the moduli space by one.
###### Proposition 13.
Let $`(M\mathrm{\#}_l(S^1\times S^3),g_l)`$ be the connected sum of $`M`$ with $`S^1\times S^3`$ with a neck of length $`l`$. For every $`\theta S^1`$ and for every $`l0`$, there exists some generic perturbation $`\eta _l\mathrm{\Omega }_+^2(M\mathrm{\#}(S^1\times S^3);i)`$ with $`\mathrm{supp}\eta _lT_l`$ such that $`SW_{\theta ,l}^1(0,\eta _l)\mathrm{}`$, where $`SW_{\theta ,l}(A,\varphi )=(D_A\varphi ,F_A^+q(\varphi ))`$ and $`\mathrm{hol}_\gamma (A)=\theta `$.
###### Proof.
Observe that the condition of $`\eta _l`$ having $`\mathrm{supp}\eta _lT_l`$ is not much of a restriction at all, because the space of such $`2`$-forms is open and the set of generic perturbations is dense (see ).
Suppose otherwise, there exists some $`\theta S^1`$ such that for every $`l0`$ we have $`SW_{\theta ,l}^1(0,\eta _l)=\mathrm{}`$. This would imply that $`SW_{\mathrm{}}^1(0,0)=\mathrm{}`$ since we have seen (see Corollary 11) that $`(M\mathrm{\#}_l(S^1\times S^3),g_l)(M_{\mathrm{}},g_{\mathrm{}})`$, but this is a contradiction because we have proven (see Corollary 11), that $`SW_{\mathrm{}}^1(0,0)\mathrm{}`$. ∎
###### Definition 8.
We will say that $`(\stackrel{~}{A}_l,\stackrel{~}{\varphi }_l)`$ on $`M_{\mathrm{}}`$, $`𝒞^0`$-extends a solution $`(A_l,\varphi _l)`$ of $`SW_{\theta ,l}(A,\varphi )=(0,\eta _l)`$ on $`M\mathrm{\#}_l(S^1\times S^3)`$ if
$`(\stackrel{~}{A}_l,\stackrel{~}{\varphi }_l)|_{M_l}`$ $`(A_l,\varphi _l)\text{ and}`$
$`(\stackrel{~}{A}_l(t,x),\stackrel{~}{\varphi }_l(t,x))`$ $`=(A_l(x),e^{2ϵt}\varphi _l(x))`$
$`\text{ for }(t,x)[l,\mathrm{})\times S^3.`$
###### Remark.
Note that $`(\stackrel{~}{A}_l,\stackrel{~}{\varphi }_l)L_{0,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$. From now on we will fix a $`U(1)`$-connection $`A`$ on $`L_𝔠`$.
###### Lemma 14.
If for every $`l0`$ there exist two different irreducible solutions $`[A_l^1,\varphi _l^1]`$ and $`[A_l^2,\varphi _l^2]`$ of
$`D_A\varphi `$ $`=0`$
$`F_A^+q(\varphi )`$ $`=\eta _l`$
on $`M\mathrm{\#}_l(S^1\times S^3)`$ for some generic perturbations $`\eta _l`$, then
$$(C_l,\psi _l)=(\stackrel{~}{\varphi }_l^1\stackrel{~}{\varphi }_l^2_{2,0,ϵ}(\stackrel{~}{A}_l^1\stackrel{~}{A}_l^2),\frac{1}{\stackrel{~}{\varphi }_l^1\stackrel{~}{\varphi }_l^2_{2,0,ϵ}}(\stackrel{~}{\varphi }_l^1\stackrel{~}{\varphi }_l^2))$$
satisfies
$`(C_l,\psi _l)(C,\psi )`$ $`L_{1,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$
$`\psi _{2,0,ϵ}`$ $`=1,`$
where $`(\stackrel{~}{A}_l^i,\stackrel{~}{\varphi }_l^i)`$ $`𝒞^0`$-extends $`(A_l^i,\varphi _l^i)`$ to $`(M_{\mathrm{}},g_{\mathrm{}})`$ for $`i=1,2`$, and $`(A_l^i,\varphi _l^i)`$ are the unique representatives obtained by the gauge fixing condition $`\delta (A_l^iA)=0`$.
###### Lemma 15.
The same hypothesis as before. If $`(\stackrel{~}{A}_l^i,\stackrel{~}{\varphi }_l^i)(A_{\mathrm{}},\varphi _{\mathrm{}})`$ in the $`L_{1,ϵ}^2(M_{\mathrm{}},g_{\mathrm{}})`$ topology, then we have
$$\underset{l\mathrm{}}{lim}D(SW_{\theta ,l})_{(\stackrel{~}{A}_l^1,\stackrel{~}{\varphi }_l^1)}(C_l,\psi _l)_{2,0,ϵ}=0$$
###### Remark.
The proof of the previous two lemmas is not difficult but technical (a straightforward computation) so we will omit the details.
###### Proposition 16.
For every $`\theta S^1`$, $`SW_\theta (𝔠_{0,1})=1`$.
###### Proof.
Assume that $`SW_\theta (𝔠_{0,1})\pm 1`$. Proposition 13 implies, for $`l0`$ there exist (at least) two different irreducible solutions $`(A_l^i,\varphi _l^i)`$, $`i=1,2`$ on $`(M\mathrm{\#}_l(S^1\times S^3),g_l)`$. By lemma 14 and lemma 15 we would have an element of $`\mathrm{ker}DSW_{\mathrm{}}`$ at $`(A_{\mathrm{}},\varphi _{\mathrm{}})`$ the unique solution on $`(M_{\mathrm{}},g_{\mathrm{}})`$, obtained in corollary 11. But this is a contradiction since $`(A_{\mathrm{}},\varphi _{\mathrm{}})`$ is a smooth point. The same kind of argument shows that $`SW_\theta (𝔠_{0,1})=1`$ since $`SW_{\mathrm{}}(𝔠)=1`$. ∎
## 6 Cohomology of $`^{}(𝔠)`$
In this section we will build cohomology classes for $`^{}(𝔠)`$ in order to detect $`B`$-classes (see Definition 5). To describe the cohomology of $`^{}(𝔠)`$ we have to introduce the concept of *universal family of $`SW`$-connections* associated to a $`\mathrm{Spin}^c`$ structure $`𝔠`$, parameterized by $`^{}(𝔠)`$. A $`SW`$-connection is simply a pair $`(A,\varphi )`$, where $`A`$ is a $`U(1)`$-connection on $`L_𝔠`$ and $`0\varphi S^+(𝔠)`$.
A cohomology class $`\beta H^i(^{}(𝔠);)`$ can be thought of as a homomorphism $`\beta :H_i(^{}(𝔠);)`$, and the elements of $`H_i(^{}(𝔠);)`$ can be thought of as homotopic classes of maps $`f:T^{}(𝔠)`$, where $`T`$ is a compact space. The maps $`f:T^{}(𝔠)`$ are naturally interpreted in terms of families of $`SW`$-connections.
###### Definition 9.
*A family of $`SW`$-connections in a bundle $`L_𝔠M`$ parametrized by a space $`T`$* is a bundle $`LT\times M`$ with the property that each slice $`L_t=L|_{\{t\}\times M}`$ is isomorphic to $`L_𝔠`$, together with a $`SW`$-connection $`(A_\varphi )_t=(A_t,\varphi _t)`$ in $`L_t`$, forming a family $`A_\varphi =\{(A_\varphi )_t\}`$.
Let $`p_2:𝒞^{}(𝔠)\times MM`$ be the projection onto the second factor and let $`_𝔠𝒞^{}(𝔠)\times M`$ be the pull-back line bundle, $`p_2^{}L_𝔠`$. Then $`_𝔠`$ carries a tautological family of $`SW`$-connections $`A_\varphi `$, in which the $`SW`$-connection on the slice $`_𝔠|_{\{(A,\varphi )\}}`$ over $`\{(A,\varphi )\}\times M`$ is $`(p_2^{}(A),p_2^{}(\varphi ))`$. The group $`𝒢(𝔠)`$ acts freely on $`𝒞^{}(𝔠)\times M`$ as well as on $`_𝔠=𝒞^{}(𝔠)\times L_𝔠`$, and there is therefore a quotient bundle
$`𝕃_𝔠`$ $`^{}(𝔠)\times M`$
$`𝕃_𝔠`$ $`=_𝔠/𝒢(𝔠).`$
The family of $`SW`$-connections $`A_\varphi `$ is preserved by $`𝒢(𝔠)`$, so $`𝕃_𝔠`$ carries an inherited family of $`SW`$-connections $`𝔸_\varphi `$. This is the *universal family of $`SW`$-connections in $`L_𝔠M`$ parameterized by $`^{}(𝔠)`$*.
If a family of $`SW`$-connections is parameterized by a space $`T`$ and carried by a bundle $`LT\times M`$, there is an associated map $`f:T^{}(𝔠)`$ given by
$$f(t)=[A_t,\varphi _t].$$
Conversely, given $`f:T^{}(𝔠)`$ there is a corresponding pull-back family of connections carried by $`(f\times I)^{}𝕃_𝔠`$. These two constructions are inverses of one another: if $`f`$ is determined by the above equation, then for each $`t`$ there is a *unique* isomorphism $`\psi _t`$ between the $`SW`$-connections in $`L_t`$ and $`(f\times I)^{}(𝕃_𝔠)_t`$, and as $`t`$ varies these fit together to form an isomorphism $`\psi :L(f\times I)^{}𝕃_𝔠`$ between these two families. (The uniqueness of $`\psi _t`$ results from the fact that $`𝒢(𝔠)`$ acts freely on $`𝒞^{}(𝔠)`$). Thus:
###### Lemma 17.
The maps $`f:T^{}(𝔠)`$ are in one-to-one correspondence with families of $`SW`$-connections on $`M`$ parameterized by $`T`$, and this correspondence is obtained by pulling back from the universal family $`(𝕃_𝔠,𝔸_\varphi )`$.
###### Remark.
Let $`\{\gamma _i\}`$ be fixed representatives for the generators of the free part of $`H_1(M;)`$. If $`f_1,f_2:T^{}(𝔠)`$ are homotopic, the corresponding bundles $`L_1`$ and $`L_2`$ are isomorphic, and the corresponding holonomy maps $`h_1:T(S^1)^{b_1}`$ and $`h_2:T(S^1)^{b_1}`$ are homotopic, where the holonomy map is defined as $`h_i(t)=(\mathrm{hol}_{\gamma _1}(f_i(t)),\mathrm{},\mathrm{hol}_{\gamma _{b_1}}(f_i(t)))`$.
There is a general construction which produces cohomology classes in $`^{}(𝔠)`$, using the slant-product pairing
$$/:H^{di}(^{}(𝔠);)\times H_i(M;)H^i(^{}(𝔠);).$$
We have built over $`^{}(𝔠)\times M`$ a line bundle $`𝕃_𝔠`$, so we can define a map
$$\mu :H_i(X;)H^{2i}(^{}(𝔠);)$$
by
$$\mu (\alpha )=c_1(𝕃_𝔠)/\alpha .$$
If $`T`$ is any $`(2i)`$-cycle in $`^{}(𝔠)`$, the class $`\mu (\alpha )`$ can be evaluated on $`T`$ using the formula
$$\mu (\alpha ),T_{^{}(𝔠)}=c_1(𝕃_𝔠),T\times \alpha _{^{}(𝔠)\times M},$$
which expresses the fact that the slant product is the adjoint of the cross-product homomorphism. Next we will describe another way to build cohomology classes.
###### Definition 10.
A closed curve $`\gamma :S^1M`$ induces a *holonomy map*
$$\mathrm{hol}_\gamma :^{}(𝔠)S^1$$
defined as the holonomy of the $`SW`$-connections $`A_\varphi `$ along $`\gamma `$. The pull-back of the canonical class $`d\theta `$ of $`S^1`$ defines a cohomology class on $`H^1(^{}(𝔠);)`$ which we will call the *holonomy class along $`\gamma `$*.
###### Proposition 18.
The cohomology groups of $`^{}(𝔠)`$ are generated by the image of the map $`\mu :H_i(X;)H^{2i}(^{}(𝔠);)`$. Moreover, given $`\gamma H_1(M;)`$, $`\mu (\gamma )`$ is the holonomy class along $`\gamma `$, $`\mathrm{hol}_\gamma ^{}(d\theta )`$.
###### Proof.
First we will prove that if $`\{\gamma _i\}`$ are fixed representatives for the generators for the free part of $`H_1(M;)`$ then $`\{\mu (\gamma _i)\}`$ generates $`H^1(^{}(𝔠);)`$. It is enough to prove that for every $`i`$ we can find $`\beta _i:S^1^{}(𝔠)`$ such that $`\mu (\gamma _i),\beta _i|_{^{}(𝔠)}=1`$. Consider the line bundle $`\gamma _i^{}L_𝔠S^1`$, and observe that there is no obstruction to extend it to a line bundle $`LS^1\times S^1`$ such that $`\mathrm{deg}L=c_1(L),S^1\times S^1=1`$. Let $`A_i`$ be a $`U(1)`$-connection on $`L`$ and consider the map
$$\mathrm{hol}_{\times S^1}(A_i):S^1S^1.$$
It is not difficult to see that $`\mathrm{deg}L=\mathrm{deg}(\mathrm{hol}_{\times S^1}(A_i))`$. After extending $`A_i(t,\gamma _i)`$ to a $`U(1)`$-connection on $`L_𝔠M`$ for each $`t`$, we obtain (see remark below lemma 17) our desired maps $`\beta _i:S^1^{}(𝔠)`$.
To prove the last statement we proceed as follows: let $`\alpha :S^1^{}(𝔠)`$,
$`\mu (\gamma _i),\alpha _{^{}(𝔠)}`$ $`=c_1(𝕃_𝔠),\alpha \times \gamma _i_{^{}(𝔠)\times M}`$
$`=c_1((\alpha \times \gamma _i)^{}(𝕃_𝔠)),S^1\times S^1`$
$`=\mathrm{deg}(\mathrm{hol}_{\times S^1}(A_i):S^1S^1)`$
$`=\mathrm{deg}(\mathrm{hol}_{\gamma _i}\alpha :S^1S^1)`$
$`=\mathrm{deg}_{\beta _i}^{}(d\theta ),\alpha _{^{}(𝔠)}.`$
Finally we have to show that if $`xM`$ then $`\mu (x)`$ generates the cohomology of the $`^{\mathrm{}}`$ factor. Since $`Map(M,S^1)_o`$ acts freely on $`𝒞^{}(𝔠)`$, then it is easy to show that $`𝕃_𝔠|_{^{}(𝔠)}𝒞^{}(𝔠)/𝒢_0(𝔠)`$, where $`𝒢_0(𝔠)`$ is the kernel of the homomorphism $`𝒢(𝔠)S^1`$ given by evaluating on the fiber over $`x`$. ∎
## 7 Applications
C. LeBrun showed that under some mild conditions on $`M`$, $`M\mathrm{\#}k\overline{^2}`$ does not admit Einstein metrics. The precise statement is the following:
###### Theorem (C. LeBrun).
Let $`M`$ be a smooth compact oriented $`4`$-manifold with $`2e+3\sigma >0`$. Assume, moreover, that $`M`$ has a non-trivial Seiberg-Witten invariant. If $`k\frac{25}{57}(2e+3\sigma )`$ then $`M\mathrm{\#}k\overline{^2}`$ does not admit an Einstein metric.
###### Remark.
The proof of this theorem only requires that $`M`$ has a $`\mathrm{Spin}^c`$-structure $`𝔠`$ that is a $`B`$-class.
###### Theorem A.
Let $`(M,𝔠)`$ be a smooth compact Kähler surface with a $`\mathrm{Spin}^c`$-structure $`𝔠`$. There is a canonical $`\mathrm{Spin}^c`$ structure in the connected sum manifold $`M\mathrm{\#}(S^1\times S^3)`$ which we will denote by $`𝔠_{0,1}`$. Moreover $`d(𝔠_{0,1})=d(𝔠)+1`$. If $`𝔠`$ is a non-trivial SW-class for $`M`$ then $`𝔠_{0,1}`$ is a $`B`$-class for the connected sum $`M\mathrm{\#}(S^1\times S^3)`$.
###### Proof.
$`SW_\theta (𝔠_{0,1})`$ is a cobordism invariant for every $`\theta S^1`$. Consider the smooth cobordism induced by the family of metrics $`g_l`$ on $`M\mathrm{\#}(S^1\times S^3)`$ as $`l\mathrm{}`$ and observe (corollary 11) that $`SW_{\mathrm{}}(𝔠)=1`$. This shows that
$$\mathrm{hol}_\gamma ^{}(d\theta ),(𝔠_{0,1})|_{^{}(𝔠_{0,1})}=1,$$
where $`\gamma `$ is a representative for the $`S^1`$ factor of the connected sum. This, the definition of a $`B`$-class and Proposition 18 complete the proof. ∎
###### Corollary 19.
Let $`(M,𝔠)`$ be a smooth compact oriented Kähler surface with a $`\mathrm{Spin}^c`$-structure $`𝔠`$. There is a canonical $`\mathrm{Spin}^c`$ structure in the connected sum $`M\mathrm{\#}2(S^1\times S^3)`$ which we will denote by $`𝔠_{0,2}`$. Moreover $`d(𝔠_{0,2})=d(𝔠)+2`$. If $`𝔠`$ is a non-trivial SW-class then $`𝔠_{0,2}`$ is a $`B`$-class but has trivial Seiberg-Witten invariant.
###### Proof.
Theorem A shows that every time that we perform a connected sum with $`S^1\times S^3`$ we *add a cycle* to the moduli space, that lies entirely in the $`H^1(M\mathrm{\#}(S^1\times S^3);)/H^1(M\mathrm{\#}(S^1\times S^3);)`$ part of $`^{}(𝔠_{0,1})`$. ∎
###### Lemma 20.
Let $`(M,𝔠)`$ be a smooth compact oriented Kähler surface with a $`\mathrm{Spin}^c`$-structure $`𝔠`$ and $`2e+3\sigma >0`$. Assume that $`𝔠`$ is a non-trivial SW-class. Let $`k,l`$ be any two natural numbers. Then there is a $`B`$-class $`𝔠_{k,l}`$ on $`M_{k,l}=M\mathrm{\#}k\overline{^2}\mathrm{\#}l(S^1\times S^3)`$ such that
$$(c_1^+(𝔠_{k,l}))^2(2e+3\sigma )(M).$$
###### Proof.
First observe that $`M_{k,l}=(M\mathrm{\#}k\overline{^2})_{0,l}`$. Since $`M`$ is a Kähler surface, we know that $`M\mathrm{\#}k\overline{^2}`$ is also a Kähler surface, and its associated $`\mathrm{Spin}^c`$ structure $`𝔠_{k,0}`$ satisfies $`c_1(𝔠_{k,0})=c_1(𝔠)+_{j=1}^kE_j`$, where $`E_1,\mathrm{},E_k`$ are generators for the pull-backs to $`M\mathrm{\#}k\overline{^2}`$ of the $`k`$ copies of $`H^2(^2,)`$ so that
$$c_1^+(𝔠)E_j0,j=1,\mathrm{},k.$$
Let $`c_1(𝔠_{k,l})`$ be the first Chern class of $`(𝔠_{k,0})_{0,l}`$ which is a $`B`$-class by theorem A, and notice that $`c_1(𝔠_{k,l})=c_1(𝔠_{k,0})`$. One then has
$`(c_1^+(𝔠_{k,l}))^2`$ $`=(c_1^+(𝔠_{k,0}))^2`$
$`=\left(c_1^+(𝔠)+{\displaystyle \underset{j=1}{\overset{k}{}}}E_j^+\right)^2`$
$`=(c_1^+(𝔠))^2+2{\displaystyle \underset{j=1}{\overset{k}{}}}c_1^+(𝔠_{0,l})E_j^++({\displaystyle \underset{j=1}{\overset{k}{}}}E_j^+)^2`$
$`(c_1^+(𝔠))^2`$
$`(c_1(𝔠))^2`$
$`=(2e+3\sigma )(M).`$
LeBrun’s theorem can be generalized in the following way:
###### Theorem 21.
Let $`(M,𝔠)`$ be a smooth compact oriented Kähler surface with a $`\mathrm{Spin}^c`$-structure $`𝔠`$ and $`2e+3\sigma >0`$. Assume that $`𝔠`$ is a $`B`$-class. If $`k+4l\frac{25}{57}(2e+3\sigma )`$ then $`M_{k,l}=M\mathrm{\#}k\overline{^2}\mathrm{\#}l(S^1\times S^3)`$ does not admit an Einstein metric.
###### Proof.
The proof is the same as the one given by C. LeBrun . ∎
There exists two well known topological obstructions to the existence of Einstein metrics on a differentiable compact oriented $`4`$-manifold $`M`$.
The first one is Thorpe’s inequality (see ), that comes from the Gauss-Bonnet-Chern formula for the Euler characteristic $`e(M)`$ of $`M`$ and from the Hirzebruch formula for the signature $`\sigma (M)`$ of $`M`$, which allow us to express these two topological invariants in terms of the irreducible components of the curvature under the action of $`SO(4)`$. It can be stated in the following way
###### Theorem (N. Hitchin, J. Thorpe).
Let $`M`$ be a compact oriented manifold of dimension $`4`$. If $`e(M)<\frac{3}{2}|\sigma (M)|`$ then $`M`$ does not admit any Einstein metric. Moreover, if $`e(M)=\frac{3}{2}|\sigma (M)`$ then $`M`$ admits no Einstein metric unless it is either flat, or a $`K3`$ surface, or an Enriques surface, or the quotient of an Enriques surface by a free antiholomorphic involution.
This theorem implies a previous result of M. Berger who proved that there exists no compact Einstein $`4`$-manifold with a negative Euler characteristic. On the other hand, combining the Gauss-Bonnet-Chern formula for the Euler characteristic with Gromov’s estimation of simplicial volume $`M`$ of a Riemannian manifold $`M`$ (see ), M. Gromov obtained the following obstruction
###### Theorem (M. Gromov).
Let $`M`$ be a compact manifold of dimension $`4`$. If $`e(M)<\frac{1}{2592\pi ^2}M`$ then $`M`$ does not admit any Einstein metric.
A. Sambusetti (see ) found a topological obstruction to the existence of Einstein metrics on compact $`4`$-manifolds which admit a non-zero degree map onto some compact real or complex hyperbolic $`4`$-manifold. As a consequence, by connected sums, he produces infinitely many non-homeomorphic $`4`$-manifolds which admit no Einstein metrics. This fact is not a consequence of Hitchin-Thorpe’s or Gromov’s obstruction theorems. A. Sambusetti also proves that any Euler characteristic and signature can be simultaneously realized by these non-homeomorphic manifolds admitting no Einstein metrics.
###### Definition 11.
We say that a pair $`(m,n)^2`$ is *admissible* if there exists a smooth compact oriented $`4`$-manifold with Euler characteristic $`m`$ and signature $`n`$. In fact a necessary and sufficient condition for $`(m,n)^2`$ to be an admissible pair is that $`mnmod2`$.
To prove our last result we need the following theorem by Z. Chen (see ).
###### Theorem (Z. Chen).
Let $`x`$, $`y`$ be integers satisfying
$`{\displaystyle \frac{352}{89}}x+140.2x^{2/3}<y`$ $`<{\displaystyle \frac{18644}{2129}}x365.7x^{2/3},`$
$`x`$ $`>C,`$
where $`C`$ is a large constant. There exists a simply connected minimal surface $`M`$ of general type with $`c_1^2(M)=y`$, $`\chi (M)=x`$. Furthermore, $`M`$ can be represented by a surface admitting a hyperelliptic fibration.
###### Remark.
Recall that $`\chi (M)`$ denotes the Euler-Poincaré characteristic of the invertible sheaf $`𝒪_M`$. Using Noether’s formula we have that
$`\chi (M)`$ $`={\displaystyle \frac{c_1^2(M)+e(M)}{12}}`$
$`={\displaystyle \frac{e(M)+\sigma (M)}{4}}.`$
If $`M`$ is not a complex surface $`e(M)+\sigma (M)`$ is not necessarily a multiple of $`4`$ but it is always an even number.
###### Theorem B.
For each admissible pair $`(m,n)`$ there exist an infinite number of non-homeomorphic compact oriented $`4`$-manifolds which have Euler characteristic $`m`$ and signature $`n`$, with free fundamental group and which do not admit Einstein metric.
###### Proof.
Let $`(m_0,n_0)`$ be an admissible pair and consider the pair of integers $`(x_0^{},y_0)=(\frac{m_0+n_0}{2},2m_0+3n_0)`$. It is always possible to find (infinitely many) positive integers $`k`$ and $`l`$ such that
$`(x,y)=({\displaystyle \frac{x_0^{}+l}{2}},y_0+k)`$ $`𝒵`$
$`4l+{\displaystyle \frac{32}{57}}k`$ $`{\displaystyle \frac{25}{57}}y_0`$
where $`𝒵`$ denotes the set of $`(x,y)^2`$ that satisfy the conditions of Chen’s theorem. The reason for this last statement is that the region $`𝒵_{}`$ determine by $`(x,y)^2`$ such that
$`{\displaystyle \frac{352}{89}}x+140.2x^{2/3}<y`$ $`<{\displaystyle \frac{18644}{2129}}x365.7x^{2/3},`$
$`x`$ $`>C,`$
is open, connected and not bounded, where $`C`$ is the same constant as in Chen’s theorem.
If we denote by $`M`$ the simply connected Kähler surface with $`c_1^2=y`$ and $`\chi =x`$, then $`M_{k,l}=M\mathrm{\#}k\overline{^2}\mathrm{\#}l(S^1\times S^3)`$, is a manifold that realizes the pair $`(m_0,n_0)`$ and does not admit any Einstein metric. This last statement is a consequence of theorem 21. ∎ |
warning/0002/quant-ph0002010.html | ar5iv | text | # Inverse Time–Dependent Quantum Mechanics
\[
## Abstract
Using a Bayesian method for solving inverse quantum problems, potentials of quantum systems are reconstructed from coordinate measurements in non–stationary states. The approach is based on two basic inputs: 1. a likelihood model, providing the probabilistic description of the measurement process as given by the axioms of quantum mechanics, and 2. additional a priori information implemented in form of stochastic processes over potentials.
\]
The first step to be done when applying quantum mechanics to a real world system is the reconstruction of its Hamiltonian from observational data. Such a reconstruction, also known as inverse problem, constitutes a typical example of empirical learning. Whereas the determination of potentials from spectral and from scattering data has been studied in much detail in inverse spectral and inverse scattering theory , this Paper describes the reconstruction of potentials by measuring particle positions in coordinate space for finite quantum systems in time–dependent states. The presented method can easily be generalized to other forms of observational data.
In the last years much effort has been devoted to many other practical empirical learning problems, including, just to name a few, prediction of financial time series, medical diagnosis, and image or speech recognition. This also lead to a variety of new learning algorithms, which should in principle also be applicable to inverse quantum problems. In particular, this Paper shows how the Bayesian framework can be applied to solve problems of inverse time–dependent quantum mechanics (ITDQ). The presented method generalizes a recently introduced approach for stationary quantum systems . Compared to stationary inverse problems, the observational data in time–dependent problems are more indirectly related to potentials, making them in general more difficult to solve.
Specifically, we will study the following type of observational data: Preparing a particle in an eigenstate of the position operator with coordinates $`x_0`$ at time $`t_0`$, we let this state evolve in time according to the rules of quantum mechanics and measure its new position at time $`t_1`$, finding a value $`x_1`$. Continuing from this measured position $`x_1`$, we measure the particle position again at time $`t_2`$, and repeat this procedure until $`n`$ data points $`x_i`$ at times $`t_i`$ have been collected. We thus end up with observational data of the form $`D`$ =$`\{(x_i,\mathrm{\Delta }_i,x_{i1})|1in\}`$, where $`x_i`$ is the result of the $`i`$-th coordinate measurement, $`\mathrm{\Delta }_i`$ = $`t_it_{i1}`$ the time interval between two subsequent measurements and $`x_{i1}`$ the coordinates of the previous observation (or preparation) at time $`t_{i1}`$.
We will discuss in particular systems with time–independent Hamiltonians of the form $`H`$ = $`T+V`$, consisting of a standard kinetic energy term $`T`$ and a local potential $`V(x,x^{})`$ = $`\delta (xx^{})v(x)`$, with $`x`$ denoting the position of the particle. In that case, the aim is the reconstruction of the function $`v(x)`$ from observational data $`D`$. (The restriction to local potentials simplifies the numerical calculations. Nonlocal Hamiltonians can be reconstructed similarly.)
Setting up a Bayesian model requires the definition of two probabilities: 1. the probability $`p(D|v)`$ to measure data $`D`$ given potential $`v`$, which, for $`D`$ considered fixed, is also known as the likelihood of $`v`$, and 2. a prior probability $`p(v)`$ implementing available a priori information concerning the potential to be reconstructed.
Referring to a maximum a posteriori approximation (MAP) we understand those potentials $`v`$ to be solutions of the reconstruction problem, which maximize $`p(v|D)`$, i.e., the posterior probability of $`v`$ given all available data $`D`$. The basic relation is then Bayes’ theorem, according to which $`p(v|D)p(D|v)p(v)`$.
One possibility is to choose a parametric ansatz for the potential $`v`$. In that case, an additional prior term $`p(v)`$ is often not included (so the MAP becomes a maximum likelihood approximation). In the following, we concentrate on nonparametric approaches, which are less restrictive compared to their parametric counterparts. Their large flexibility, however, makes it essential to include (nonuniform) priors. Corresponding nonparametric priors are formulated explicitly in terms of the function $`v(x)`$ . Indeed, nonparametric priors are well known from applications to regression , classification , general density estimation , and stationary inverse quantum problems . It is the likelihood model, discussed next, which is specific for ITDQ.
According to the axioms of quantum mechanics the probability that a particle is found at position $`x_i`$ at time $`t_i`$, provided the particle has been at $`x_{i1}`$ at time $`t_{i1}`$, is given by
$$p_i=p(x_i|\mathrm{\Delta }_i,x_{i1},v)=|\varphi _i(x_i)|^2,$$
(1)
where
$$\varphi _i(x_i)=<x_i|\varphi _i>=<x_i|U_ix_{i1}>,$$
(2)
are matrix elements of the time evolution operator
$$U_i=e^{i\mathrm{\Delta }_iH},$$
(3)
setting $`\mathrm{}`$ = 1. The transition amplitudes (2) can be calculated by inserting orthonormalized eigenstates $`\psi _\alpha `$ of $`H`$, with energies $`E_\alpha `$,
$$\varphi _i(x_i)=\underset{\alpha }{}e^{i\mathrm{\Delta }_iE_\alpha }\psi _\alpha (x_i)\psi _\alpha ^{}(x_{i1}).$$
(4)
Clearly, it is straightforward to modify (1) for measuring observables different from the particle position. It is also interesting to note that the transition probabilities (1) define a Markoff process with $`W_i(xx^{})`$ = $`p(x^{}|\mathrm{\Delta }_i,x,v)`$. For real eigenfunctions $`\psi _\alpha (x)`$, i.e., for a real Hamiltonian with real boundary conditions, they obey the relation $`W_i(xx^{})`$ = $`W_i(x^{}x)`$. It follows that the detailed balance condition, $`p_{\mathrm{stat}}(x)W_i(xx^{})`$ = $`p_{\mathrm{stat}}(x^{})W_i(x^{}x)`$, is fulfilled for a uniform $`p_{\mathrm{stat}}(x)`$, which therefore represents the stationary state of the Markoff process of repeated position measurements.
Having defined the likelihood model of ITDQ, in the next step a prior for $`v`$ has to be chosen. A convenient nonparametric prior $`p(v)`$ is a Gaussian
$$p_G(v)=\left(det\frac{𝐊_0}{2\pi }\right)^{\frac{1}{2}}e^{\frac{1}{2}<vv_0|𝐊_0|vv_0>},$$
(5)
with (real symmetric, positive semi–definite) inverse covariance $`𝐊_0`$, acting in the space of potentials, and mean $`v_0(x)`$, which can be considered as a reference potential for $`v`$. Typical examples are smoothness constraints on $`v`$ which correspond to choosing differential operators for $`𝐊_0`$. Reference potentials can be made more flexible by allowing parameterized families $`v_0(x;\theta )`$. Within the context of Bayesian statistics such additional parameters $`\theta `$ are known as hyperparameters. In MAP approximation the optimal hyperparameters are determined by maximizing the posterior (7) simultaneously with respect to $`\theta `$ and $`v(x)`$ . A simplified procedure consists in using a parametric approximation $`v(\theta )`$ which maximizes the likelihood $`_ip_i(\theta )`$ as reference potential $`v_0`$ for the nonparametric reconstruction $`v(x)`$ .
If available, it is useful to include some information about the ground state energy $`E_0(v)`$, which helps to determine the depth of the potential. This can, for example, be a noisy measurement of the ground state energy which, assuming Gaussian noise, is implemented by
$$p_Ee^{\frac{\mu }{2}\left(E_0(v)\kappa \right)^2}.$$
(6)
Combining (5) and (6) with (1) for $`n`$ repeated coordinate measurements starting from an initial position $`x_0`$, we obtain for the posterior (7),
$$p(v|D)p_G(v)p_E(v)\underset{i=1}{\overset{n}{}}p_i.$$
(7)
To calculate the MAP solution $`v^{}`$ =$`\mathrm{argmax}_vp(v|D)`$ we set the functional derivative of the posterior (7), or technically more convenient of its logarithm, with respect to $`v`$, denoted $`\delta _v`$, to zero. This yields,
$$0=\delta _v\mathrm{ln}p(v|D)=\delta _v\mathrm{ln}p_G(v)+\delta _v\mathrm{ln}p_E(v)+\underset{i}{}\delta _v\mathrm{ln}p_i,$$
(8)
with
$`\delta _v\mathrm{ln}p_G(v)`$ $`=`$ $`𝐊_0(vv_0),`$ (9)
$`\delta _v\mathrm{ln}p_E(v)`$ $`=`$ $`\mu \left(E_0(v)\kappa \right)\delta _vE_0(v),`$ (10)
$`\delta _v\mathrm{ln}p_i`$ $`=`$ $`2\mathrm{R}\mathrm{e}[\varphi _i^1(x_i)\delta _v\varphi _i(x_i)].`$ (11)
The functional derivative $`\delta _v\varphi _i`$ can, according to Eq. (4), be obtained from $`\delta _v\psi _\alpha `$. The still required $`\delta _v\psi _\alpha `$ and $`\delta _vE_\alpha `$ can then be found by calculating the functional derivative of the eigenvalue equation $`H\psi _\alpha `$ = $`E_\alpha \psi _\alpha `$. Using
$$\delta _{v(x)}V(x^{},x^{\prime \prime })=\delta (xx^{})\delta (x^{}x^{\prime \prime }),$$
(12)
$`\delta _{v(x)}`$ denoting the $`x`$ component of functional derivative $`\delta _v`$, we find,
$`\delta _{v(x)}E_\alpha `$ $`=`$ $`<\psi _\alpha |\delta _{v(x)}H|\psi _\alpha >=|\psi _\alpha (x)|^2,`$ (13)
$`\delta _{v(x)}\psi _\alpha (x^{})`$ $`=`$ $`{\displaystyle \underset{\gamma \alpha }{}}{\displaystyle \frac{1}{E_\alpha E_\gamma }}\psi _\gamma (x^{})\psi _\gamma ^{}(x)\psi _\alpha (x).`$ (14)
Collecting the results, gives
$`\delta _{v(x)}\varphi _i(x_i)=\delta _{v(x)}<x_i|U_ix_{i1}>=`$ (15)
$`{\displaystyle \underset{\alpha }{}}e^{i\mathrm{\Delta }_iE_\alpha }\left[(i\mathrm{\Delta }_i|\psi _\alpha (x)|^2)\right)\psi _\alpha (x_i)\psi _\alpha ^{}(x_{i1})`$ (16)
$`+{\displaystyle \underset{\gamma \alpha }{}}{\displaystyle \frac{1}{E_\alpha E_\gamma }}\psi _\gamma (x_i)\psi _\gamma ^{}(x)\psi _\alpha (x)\psi _\alpha ^{}(x_{i1})`$ (17)
$`+{\displaystyle \underset{\gamma \alpha }{}}{\displaystyle \frac{1}{E_\alpha E_\gamma }}\psi _\gamma ^{}(x_{i1})\psi _\gamma (x)\psi _\alpha ^{}(x)\psi _\alpha (x_i)].`$ (18)
Inserting Eq. (13) for $`\alpha `$ = 0 in Eq. (10) and Eq. (18) in Eq. (11) a MAP solution for the potential $`v`$ can be found by iterating the stationarity equation (8) numerically on a lattice. Clearly, such a straightforward discretization can only be expected to work for a low–dimensional $`x`$ variable. Higher dimensional systems usually require additional approximations .
As the next step, we want to check the numerical feasibility of a nonparametric reconstruction of the potential $`v`$ for a one–dimensional quantum system. For that purpose, we choose a system with the true potential
$$v_{\mathrm{true}}(x)=\frac{c_1}{\sqrt{2\pi \sigma }}e^{\frac{(xc_2)^2}{2\sigma ^2}},$$
(19)
where $`c_1`$ = $`10`$, $`c_2`$ = $`2`$, and $`\sigma `$ = 2. An example of the time evolution of an unobserved particle in the potential $`v_{\mathrm{true}}`$ is shown in Fig. 1. As input for the reconstruction algorithm 50 data points $`x_i`$ are sampled from the corresponding true likelihoods $`p(x|\mathrm{\Delta }_i,x_{i1},v_{\mathrm{true}})`$. A corresponding path of an observed particle is shown in Fig. 2.
Besides a noisy energy measurement of the form (6) we include a Gaussian prior (5) with a smoothness related inverse covariance
$$𝐊_0(x,x^{})=\delta (xx^{})\lambda \underset{k=0}{\overset{3}{}}(1)^k\frac{\sigma _0^{2m}}{k!2^k}\left(\frac{^2}{x^2}\right)^k.$$
(20)
To obtain an adapted reference potential $`v_0`$ for the Gaussian prior, a parameterized potential of the form
$$v_0(a,b,c)=\mathrm{min}[0,a(xb)^2+c],$$
(21)
is optimized with respect to $`a`$, $`b`$, $`c`$ by maximizing the “extended likelihood” $`_i\mathrm{ln}p_i(v_0)+\mathrm{ln}p_E(v_0)`$. Finally, the stationarity equation (8) is solved by iterating according to
$`v^{(r+1)}=v^{(r)}+\eta [v_0v^{(r)}+`$ (22)
$`𝐊_0^1\{2{\displaystyle \underset{i}{\overset{n}{}}}\mathrm{Re}[\delta _v\mathrm{ln}\varphi _i(x_i)]+\delta _v(\mathrm{ln}p_G+\mathrm{ln}p_E)\}].`$ (23)
The resulting nonparametric ITDQ solution $`v_{\mathrm{ITDQ}}`$ (see Fig. 3), is a reasonable reconstruction of $`v_{\mathrm{true}}`$, and clearly better than the best parametric approximation $`v_0(a,b,c)`$. It is only the flat area near the right border where, due to missing and unrepresentative data, the reconstruction differs significantly from the true potential.
Fig. 4 compares the sum over empirical transition probabilities $`\frac{1}{n}_{i=1}^n\delta (xx_i)`$ as derived from the observational data $`D`$ with the corresponding true $`p_{\mathrm{true}}`$ = $`\frac{1}{n}_{i=1}^np(x|\mathrm{\Delta }_i,x_{i1},v_{\mathrm{true}})`$ and reconstructed $`p_{\mathrm{ITDQ}}`$ = $`\frac{1}{n}_{i=1}^np(x|\mathrm{\Delta }_i,x_{i1},v_{\mathrm{ITDQ}})`$. Due to the summation over data points with different $`x_{i1}`$, the quantities shown in Fig. 4 do not present the complete information which is available to the algorithm. Hence, Fig. 5 depicts the corresponding quantities for a fixed $`x_{i1}`$. In particular, Fig. 5 compares the reconstructed transition probability (1) with the corresponding empirical and true transition probabilities for a particle having been at time $`t_{i1}`$ at position $`x_{i1}`$ =1. The ITDQ algorithm returns an approximation for all such transition probabilities.
Figs. 4 and 5 show, that the reconstructed $`v_{\mathrm{ITDQ}}`$ tends to produce a better approximation of the empirical probabilities than the true potential $`v_{\mathrm{true}}`$. Indeed, the error on the data or negative log–likelihood, $`ϵ_D(v)`$ = $`_i\mathrm{ln}p_i(v)`$, being a canonical error measure in density estimation, is smaller for $`v_{\mathrm{ITDQ}}`$ than for $`v_{\mathrm{true}}`$. A smaller $`\lambda `$, i.e., a lower influence of the prior, produces a still smaller error $`ϵ_D(v_{\mathrm{ITDQ}})`$. At the same time, however, the reconstructed potential becomes more wiggly for smaller $`\lambda `$, being the symptom of the well known effect of “overfitting”. The (true) generalization error $`ϵ_g(v)`$ = $`𝑑x𝑑x^{}p(x)p(x^{}|x,v_{\mathrm{true}})\mathrm{ln}p(x^{}|x,v)`$ \[with uniform $`p(x)`$\], on the other hand, can never be smaller for the reconstructed $`v_{\mathrm{ITDQ}}`$ than for $`v_{\mathrm{true}}`$. As it is typical for most empirical learning problems, the generalization error $`ϵ_g(v_{\mathrm{ITDQ}})`$ shows a minimum as function of $`\lambda `$. It is this minimum which gives the optimal value for $`\lambda `$. Knowledge of the true model allows in our case to calculate the generalization error exactly. If, as usual, the true model is not known, classical cross–validation and bootstrap techniques can be used to approximate the generalization error as function of $`\lambda `$ empirically.
Alternatively to optimizing $`\lambda `$ or other hyperparameters one can integrate over them . Similarly, studying the feasibility of a Bayesian Monte Carlo approach, contrasting the MAP approach of this paper, would certainly be interesting.
In summary, this Paper has presented a method to solve inverse problems for time–dependent quantum systems. The approach, based on a Bayesian framework, is able to handle quite general types of observational data. Numerical calculations proved to be feasible for a one dimensional model. |
warning/0002/cond-mat0002448.html | ar5iv | text | # Rouse Chains with Excluded Volume Interactions: Linear Viscoelasticity
## 1 Introduction
The simplest model, within the context of Polymer Kinetic Theory, to describe the rheological behavior of dilute polymer solutions is the Rouse model . The Rouse model represents the macromolecule by a linear chain of identical beads connected by Hookean springs, and assumes that the solvent influences the motion of the beads by exerting a drag force and a Brownian force. While the Rouse model is able to explain the existence of viscoelasticity in polymer solutions by predicting a constant non-zero first normal stress difference in simple shear flow, it cannot predict several other features of dilute solution behavior, such as the existence of a nonzero second normal stress difference, the existence of shear rate dependent viscometric functions, or the correct molecular weight dependence of material functions. Over the past decade, considerable progress has been made by incorporating the effect of fluctuating hydrodynamic interaction into the Rouse model . These models are able to predict the molecular weight dependence of the material functions accurately. They also predict a nonzero second normal stress difference, and shear rate dependent viscometric functions. However, since they neglect the existence of excluded volume interactions among parts of the polymer chain, they are strictly applicable only to theta solutions.
Recently, Prakash and Öttinger examined the influence of excluded volume effects on the rheological behavior of dilute polymer solutions by representing the polymer molecule with a Hookean dumbbell model, and using a narrow Gaussian repulsive potential to describe the excluded volume interactions between the beads of the dumbbell. The narrow Gaussian potential tends to a $`\delta `$-function potential in the limit of a parameter, that describes the width of the potential, going to zero. It can therefore be used to evaluate results obtained with the singular $`\delta `$-function potential. It was shown by them that the use of a $`\delta `$-function potential between the beads, which is commonly used in static theories for polymer solutions , leads to no change in the equilibrium or dynamic properties of the dumbbell when compared to the case where no excluded volume interactions are taken into account. They also found that assuming that the non-equilibrium configurational distribution function is a Gaussian leads to accurate predictions of viscometric functions in a certain range of parameter values. These results suggest that it would be worthwhile examining longer bead-spring chains. Firstly, it is interesting to see if the problem with the $`\delta `$-function potential can be resolved when there are more beads in the bead-spring chain. Secondly, it is important to find out if the Gaussian approximation is accurate even for longer chains. The purpose of this paper is to attempt to answer these questions, in the linear viscoelastic limit, by extending the methodology developed in the earlier paper to the case of bead-spring chains. The same issues, in the context of steady shear flows at finite shear rates, will be addressed in a subsequent paper.
As in the dumbbell paper, we confine our attention to excluded volume interactions alone, and neglect the presence of hydrodynamic interactions. This clearly implies—since it is essential to include hydrodynamic interaction effects for a proper description of the dynamic behavior of dilute solutions—that the results of the present paper are not yet directly comparable with experiments. They represent a preliminary step in that direction. It is felt that the inclusion of hydrodynamic interaction would make the theory significantly more complex before the role of excluded volume interactions is properly understood. The aim of this work is to develop and carefully evaluate the Gaussian approximation for excluded volume interactions. The Gaussian approximation has already been shown to be excellent for the treatment of hydrodynamic interaction effects . If it also proves to be accurate for the treatment of excluded volume effects, then it would be extremely useful for describing the combined effects of hydrodynamic interaction and excluded volume. It should be noted that, in principle, the development of such an approximation does not pose any fundamental problems.
This paper is organized as follows. In the next section, the basic equations governing the dynamics of Rouse chains with excluded volume interactions are discussed. A retarded motion expansion for the stress tensor is derived in section 3, and exact expressions for the zero shear rate viscometric functions in simple shear flow are obtained. The implications of these results for a $`\delta `$-function excluded volume potential are then discussed. In section 4, the Brownian dynamics simulation algorithm used in this work is described. Section 5 is devoted to the development of the Gaussian approximation for the configurational distribution function. Exact expressions for linear viscoelastic properties are derived through a codeformational memory-integral expansion. In section 6, a first order perturbation expansion in the strength of the excluded volume interaction is carried out. This proves to be very helpful in understanding the nature of the Gaussian approximation. The results of the various exact and approximate treatments are compared and discussed in section 7, and the main conclusions of the paper are summarized in section 8.
## 2 Basic Equations
The instantaneous configuration of a linear bead-spring chain, which consists of $`N`$ beads connected together by $`(N1)`$ Hookean springs, is specified by the bead position vectors $`𝒓_\nu ,\nu =1,2,\mathrm{},N,`$ in a laboratory-fixed coordinate system. The Newtonian solvent, in which the chain is suspended, is assumed to have a homogeneous velocity field—that is, at position $`𝒓`$ and time $`t`$, the velocity is given by $`𝒗(𝒓,t)=𝒗_0+𝜿(t)𝒓`$, where $`𝒗_0`$ is a constant vector and $`𝜿(t)`$ is a traceless tensor.
The microscopic picture of the intra-molecular forces within the bead-spring chain is one in which the presence of excluded volume interactions between the beads causes the chain to swell, while on the other hand, the entropic retractive force of the springs draws the beads together and opposes the excluded volume force. This is modeled by writing the potential energy $`\varphi `$ of the bead-spring chain as a sum of the potential energy of the springs $`S`$, and the potential energy due to excluded volume interactions $`E`$. The potential energy $`S`$ is the sum of the potential energies of all the springs in the chain, and is given by,
$$S=\frac{1}{2}H\underset{i=1}{\overset{N1}{}}𝑸_i𝑸_i$$
(1)
where, $`H`$ is the spring constant, and $`𝑸_i=𝒓_{i+1}𝒓_i`$, is the bead connector vector between the beads $`i`$ and $`i+1`$. The excluded volume potential energy $`E`$ is found by summing the interaction energy over all pairs of beads $`\mu `$ and $`\nu `$,
$$E=\frac{1}{2}\underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}E\left(𝒓_\nu 𝒓_\mu \right)$$
(2)
where, $`E\left(𝒓_\nu 𝒓_\mu \right)`$ is a short-range function. It is usually assumed to be a $`\delta `$-function potential in static theories for polymer solutions,
$$E\left(𝒓_\nu 𝒓_\mu \right)=vk_\mathrm{B}T\delta \left(𝒓_\nu 𝒓_\mu \right)$$
(3)
where, $`v`$ is the excluded volume parameter (with dimensions of volume), $`k_\mathrm{B}`$ is Boltzmann’s constant, and $`T`$ is the absolute temperature. In this work, we regularise the $`\delta `$-function potential, and assume that $`E\left(𝒓_\nu 𝒓_\mu \right)`$ is given by a narrow Gaussian potential,
$$E\left(𝒓_\nu 𝒓_\mu \right)=\frac{vk_\mathrm{B}T}{[2\pi \stackrel{~}{d}^2]^{\frac{3}{2}}}\mathrm{exp}\left(\frac{1}{2}\frac{𝒓_{\nu \mu }^2}{\stackrel{~}{d}^2}\right)$$
(4)
where, $`\stackrel{~}{d}`$ is a parameter that describes the width of the potential (it represents, in some sense, the extent of excluded volume interactions), and $`𝒓_{\nu \mu }=𝒓_\nu 𝒓_\mu `$, is the vector between beads $`\mu `$ and $`\nu `$. In the limit $`\stackrel{~}{d}`$ tending to zero, the narrow Gaussian potential becomes a $`\delta `$-function potential.
The intra-molecular force on a bead $`\nu `$, $`𝑭_\nu ^{(\varphi )}=(\varphi /𝒓_\nu )`$, can consequently be written as, $`𝑭_\nu ^{(\varphi )}=𝑭_\nu ^{(S)}+𝑭_\nu ^{(E)}`$, where,
$`𝑭_\nu ^{(S)}`$ $`=`$ $`{\displaystyle \frac{S}{𝒓_\nu }}=H{\displaystyle \underset{k=1}{\overset{N1}{}}}\overline{B}_{k\nu }𝑸_k`$ (5)
$`𝑭_\nu ^{(E)}`$ $`=`$ $`{\displaystyle \frac{E}{𝒓_\nu }}={\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{}{𝒓_\nu }}E\left(𝒓_\nu 𝒓_\mu \right)`$ (6)
In eq 5, $`\overline{B}_{k\nu }`$ is an $`(N1)\times N`$ matrix defined by, $`\overline{B}_{k\nu }=\delta _{k+1,\nu }\delta _{k\nu }`$, with $`\delta _{k\nu }`$ denoting the Kronecker delta.
For homogeneous flows, the internal configurations of the bead-spring chain are expected to be independent of the location of the centre of mass. Consequently, it is assumed that the configurational distribution function $`\psi `$ depends only on the $`(N1)`$ bead connector vectors $`𝑸_k`$. The diffusion equation that governs $`\psi (𝑸_1,\mathrm{},𝑸_{N1},t)`$, for a system with an intra-molecular potential energy $`\varphi `$ as described above, can then be shown to be given by,
$`{\displaystyle \frac{\psi }{t}}=`$ $``$ $`{\displaystyle \underset{j=1}{\overset{N1}{}}}{\displaystyle \frac{}{𝑸_j}}\left(𝜿𝑸_j{\displaystyle \frac{H}{\zeta }}{\displaystyle \underset{k=1}{\overset{N1}{}}}A_{jk}𝑸_k+{\displaystyle \frac{1}{\zeta }}{\displaystyle \underset{\nu =1}{\overset{N}{}}}\overline{B}_{j\nu }𝑭_\nu ^{(E)}\right)\psi `$ (7)
$`+`$ $`{\displaystyle \frac{k_\mathrm{B}T}{\zeta }}{\displaystyle \underset{j,k=1}{\overset{N1}{}}}A_{jk}{\displaystyle \frac{}{𝑸_j}}{\displaystyle \frac{\psi }{𝑸_k}}`$
where, $`\zeta `$ is the bead friction coefficient (which, for spherical beads with radius $`a`$, in a solvent with viscosity $`\eta _s`$, is given by the Stokes expression: $`\zeta =6\pi \eta _sa`$), and $`A_{jk}`$ is the Rouse matrix,
$$A_{jk}=\underset{\nu =1}{\overset{N}{}}\overline{B}_{j\nu }\overline{B}_{k\nu }=\{\begin{array}{cc}2\hfill & \text{for }|jk|=0\text{,}\hfill \\ \multicolumn{2}{c}{}\\ 1\hfill & \text{for }|jk|=1\text{,}\hfill \\ \multicolumn{2}{c}{}\\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(8)
The time evolution of the average of any arbitrary quantity, carried out with the configurational distribution function $`\psi `$, can be obtained from the diffusion equation. In particular, by multiplying eq 7 by $`𝑸_j𝑸_k`$, and integrating over all configurations, the following time evolution equation for the second moments of the bead connector vectors is obtained,
$`{\displaystyle \frac{d}{dt}}𝑸_j𝑸_k`$ $`=`$ $`𝜿𝑸_j𝑸_k+𝑸_j𝑸_k𝜿^T+{\displaystyle \frac{2k_\mathrm{B}T}{\zeta }}A_{jk}\mathrm{𝟏}`$ (9)
$``$ $`{\displaystyle \frac{H}{\zeta }}{\displaystyle \underset{m=1}{\overset{N1}{}}}\left\{𝑸_j𝑸_mA_{mk}+A_{jm}𝑸_m𝑸_k\right\}+𝒀_{jk}`$
where, $`\mathrm{𝟏}`$ is the unit tensor, and,
$$𝒀_{jk}=\frac{1}{\zeta }\underset{\mu =1}{\overset{N}{}}\left\{𝑸_j𝑭_\mu ^{(E)}\overline{B}_{k\mu }+\overline{B}_{j\mu }𝑭_\mu ^{(E)}𝑸_k\right\}$$
(10)
The term $`𝒀_{jk}`$, which arises due to the presence of excluded volume interactions, does not appear in the second moment equation for the Rouse model. Due to this term, which in general involves higher order moments, eq 9, is not a closed equation for $`𝑸_j𝑸_k`$. As will be discussed in greater detail in the section on the Gaussian approximation, finding an approximate solution for the present model involves making eq 9 a closed equation for the second moments.
The polymer contribution to the stress tensor—for models with arbitrary intra-molecular potential forces but no internal constraints—is given by the Kramers expression ,
$$𝝉^p=n_\mathrm{p}H\underset{k=1}{\overset{N1}{}}𝑸_k𝑸_k+𝒁+(N1)n_\mathrm{p}k_\mathrm{B}T\mathrm{𝟏}$$
(11)
where,
$$𝒁=n_\mathrm{p}\underset{\nu =1}{\overset{N}{}}\underset{k=1}{\overset{N1}{}}B_{\nu k}𝑸_k𝑭_\nu ^{(E)}$$
(12)
Here, $`n_\mathrm{p}`$ is the number density of polymers, and $`B_{\nu k}`$ is a $`N\times (N1)`$ matrix defined by, $`B_{\nu k}=k/N\mathrm{\Theta }(k\nu )`$, with $`\mathrm{\Theta }(k\nu )`$ denoting a Heaviside step function.
It is clear from eq 11 that there are two reasons why the presence of excluded volume interactions leads to a stress tensor that is different from that obtained in the Rouse model. Firstly, there is an additional term represented by $`𝒁`$ which is the direct influence of excluded volume effects. Secondly, there is an indirect influence due to a change in the contribution of the term $`_{k=1}^{N1}𝑸_k𝑸_k`$, relative to its contribution in the Rouse case. For a $`\delta `$-function excluded volume potential, it can be shown that the direct contribution to the stress tensor is isotropic . On the other hand, for the narrow Gaussian potential, $`𝒁`$ is not isotropic unless $`\stackrel{~}{d}`$ is equal to zero. It is therefore important to use the complete form of the Kramers expression, eq 11, when carrying out simulations with an excluded volume potential that is not a $`\delta `$-function potential.
All the rheological properties of interest can be obtained once the stress tensor in eq 11 is evaluated. In the next section, a retarded motion expansion for the stress tensor is derived.
## 3 Retarded Motion Expansion
A retarded motion expansion for the stress tensor can be obtained by extending the derivation carried out previously for the dumbbell model to the case of bead-spring chains. The dumbbell model derivation was, in turn, an adaptation of a similar development for the FENE dumbbell model . The argument in all these cases rests basically on seeking a solution of the diffusion equation, eq 7, of the following form,
$$\psi (𝑸_1,\mathrm{},𝑸_{N1},t)=\psi _{\mathrm{eq}}(𝑸_1,\mathrm{},𝑸_{N1})\varphi _{\mathrm{fl}}(𝑸_1,\mathrm{},𝑸_{N1},t)$$
(13)
where, $`\psi _{\mathrm{eq}}`$ is the equilibrium distribution function given by,
$$\psi _{\mathrm{eq}}(𝑸_1,\mathrm{},𝑸_{N1})=𝒩_{\mathrm{eq}}e^{\varphi /k_\mathrm{B}T}$$
(14)
with $`𝒩_{\mathrm{eq}}`$ denoting the normalization constant, and $`\varphi _{\mathrm{fl}}`$ is the correction to $`\psi _{\mathrm{eq}}`$ due to flow—appropriately termed the flow contribution.
The governing partial differential equation for $`\varphi _{\mathrm{fl}}(𝑸_1,\mathrm{},𝑸_{N1},t)`$ can be obtained by substituting eq 13 into the diffusion equation, eq 7. It turns out that, regardless of the form of the excluded volume potential, at steady state, an exact solution to this partial differential equation can be found for all homogeneous potential flows. For more general homogeneous flows, however, one can only obtain a perturbative solution of the form,
$$\varphi _{\mathrm{fl}}(𝑸_1,\mathrm{},𝑸_{N1},t)=1+\varphi _1+\varphi _2+\varphi _3+\mathrm{}$$
(15)
where $`\varphi _k`$ is of order $`k`$ in the velocity gradient.
Partial differential equations governing each of the $`\varphi _k`$ may be derived by substituting eq 15 into the partial differential equation for $`\varphi _{\mathrm{fl}}`$ and equating terms of like order. The forms of the functions $`\varphi _k`$ can then be guessed by requiring that they fulfill certain conditions . In the present instance, we only find the form of $`\varphi _1`$, since our interest is confined to zero shear rate properties. One can show that,
$$\varphi _1=\frac{\zeta }{4k_\mathrm{B}T}\underset{m,n=1}{\overset{N1}{}}C_{mn}𝑸_m\dot{𝜸}𝑸_n$$
(16)
where, $`\dot{𝜸}`$ is the rate of strain tensor, $`\dot{𝜸}=𝒗+𝒗^T`$, and $`C_{mn}`$ is the Kramers matrix. The Kramers matrix is the inverse of the Rouse matrix, and is defined by,
$$C_{mn}=\underset{\nu =1}{\overset{N}{}}B_{\nu m}B_{\nu n}=\mathrm{min}(m,n)mn/N$$
(17)
In order to proceed further, we need to show that the present model satisfies the Giesekus expression for the stress tensor . Upon multiplying eq 7 with $`_{m,n=1}^{N1}C_{mn}𝑸_m𝑸_n`$, and integrating over all configurations, we can show that,
$`{\displaystyle \frac{d}{dt}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}C_{mn}𝑸_m𝑸_n{\displaystyle \underset{m,n=1}{\overset{N1}{}}}C_{mn}\left[𝜿𝑸_m𝑸_n+𝑸_m𝑸_n𝜿^T\right]`$
$`={\displaystyle \frac{2k_\mathrm{B}T}{\zeta }}(N1)\mathrm{𝟏}{\displaystyle \frac{2H}{\zeta }}{\displaystyle \underset{m=1}{\overset{N1}{}}}𝑸_m𝑸_m+{\displaystyle \frac{2}{\zeta }}{\displaystyle \underset{\nu =1}{\overset{N}{}}}{\displaystyle \underset{m=1}{\overset{N1}{}}}B_{\nu m}𝑸_m𝑭_\nu ^{(E)}`$ (18)
On combining this equation with eq 11 for the stress tensor, it is straight forward to see that the Giesekus expression is indeed satisfied. At steady state the Giesekus expression reduces to,
$$𝝉^p=\frac{n_\mathrm{p}\zeta }{2}\underset{m,n=1}{\overset{N1}{}}C_{mn}\left\{𝜿𝑸_m𝑸_n+𝑸_m𝑸_n𝜿^T\right\}$$
(19)
Clearly, the stress tensor at steady state can be found once the average $`𝑸_m𝑸_n`$ is evaluated. This can be done, correct to first order in velocity gradients, by using the power series expansion for $`\varphi _{\mathrm{fl}}`$, eq 15, with the specific form for $`\varphi _1`$ in eq 16. The following retarded motion expansion for the stress tensor, correct to second order in velocity gradients and valid for arbitrary homogeneous flows, is then obtained,
$`𝝉^p`$ $`=`$ $`{\displaystyle \frac{n_\mathrm{p}\zeta }{2}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}C_{mn}\left[𝜿𝑸_m𝑸_n_{\mathrm{eq}}+𝑸_m𝑸_n_{\mathrm{eq}}𝜿^T\right]`$ (20)
$``$ $`{\displaystyle \frac{n_\mathrm{p}\zeta ^2}{8k_\mathrm{B}T}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}{\displaystyle \underset{j,k=1}{\overset{N1}{}}}C_{mn}C_{jk}[𝜿𝑸_m𝑸_n(𝑸_j\dot{𝜸}𝑸_k)_{\mathrm{eq}}`$
$`+`$ $`(𝑸_j\dot{𝜸}𝑸_k)𝑸_m𝑸_n_{\mathrm{eq}}𝜿^T]+\mathrm{}`$
where, $`X_{\mathrm{eq}}`$ denotes the average of any arbitrary quantity $`X`$ with the equilibrium distribution function $`\psi _{\mathrm{eq}}`$.
One can see clearly from eq 20 that rheological properties, at small values of the velocity gradient, can be obtained by merely evaluating equilibrium averages. The special case of steady simple shear flow in the limit of zero shear rate is considered below.
### 3.1 Zero Shear Rate Viscometric Functions
Steady simple shear flows are described by a tensor $`𝜿`$ which has the following matrix representation in the laboratory-fixed coordinate system,
$$𝜿=\dot{\gamma }\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$
(21)
where $`\dot{\gamma }`$ is the constant shear rate.
The three independent material functions used to characterize such flows are the viscosity, $`\eta _p`$, and the first and second normal stress difference coefficients, $`\mathrm{\Psi }_1\mathrm{and}\mathrm{\Psi }_2`$, respectively. These functions are defined by the following relations ,
$$\tau _{xy}^p=\dot{\gamma }\eta _p;\tau _{xx}^p\tau _{yy}^p=\dot{\gamma }^2\mathrm{\Psi }_1;\tau _{yy}^p\tau _{zz}^p=\dot{\gamma }^2\mathrm{\Psi }_2$$
(22)
The components of the stress tensor in simple shear flow, for small values of the shear rate $`\dot{\gamma }`$, can be found by substituting eq 21 for the rate of strain tensor, into eq 20. This leads to,
$`\tau _{xy}^p`$ $`=`$ $`{\displaystyle \frac{n_\mathrm{p}\zeta \dot{\gamma }}{2}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}C_{mn}Y_mY_n_{\mathrm{eq}}{\displaystyle \frac{n_\mathrm{p}\zeta ^2\dot{\gamma }^2}{4k_\mathrm{B}T}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}{\displaystyle \underset{p,q=1}{\overset{N1}{}}}C_{mn}C_{pq}Y_mY_nX_pY_q_{\mathrm{eq}}`$
$`\tau _{xx}^p`$ $`=`$ $`n_\mathrm{p}\zeta \dot{\gamma }{\displaystyle \underset{m,n=1}{\overset{N1}{}}}C_{mn}X_mY_n_{\mathrm{eq}}{\displaystyle \frac{n_\mathrm{p}\zeta ^2\dot{\gamma }^2}{2k_\mathrm{B}T}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}{\displaystyle \underset{p,q=1}{\overset{N1}{}}}C_{mn}C_{pq}X_mY_nX_pY_q_{\mathrm{eq}}`$
$`\tau _{yy}^p`$ $`=`$ $`\tau _{zz}^p=0`$ (23)
where, $`(X_m,Y_m,Z_m)`$ are the Cartesian components of the bead connector vector $`𝑸_m`$.
Using the symmetry property of the potential energy $`\varphi `$, which remains unchanged when the sign of the $`Y_k`$ component of all the bead connector vectors $`𝑸_k;k=1,2,\mathrm{},(N1)`$, is reversed, we can show that, $`X_mY_n_{\mathrm{eq}}=0`$. From the definitions of the viscometric functions in eq 22, it is straight forward to show that, in the limit of zero shear rate, the following exact expressions for the zero shear rate viscometric functions are obtained.
$`\eta _{p,0}`$ $`=`$ $`{\displaystyle \frac{n_\mathrm{p}\zeta }{6}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}C_{mn}𝑸_m𝑸_n_{\mathrm{eq}}`$ (24)
$`\mathrm{\Psi }_{1,0}`$ $`=`$ $`{\displaystyle \frac{n_\mathrm{p}\zeta ^2}{2k_\mathrm{B}T}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}{\displaystyle \underset{p,q=1}{\overset{N1}{}}}C_{mn}C_{pq}X_mY_nX_pY_q_{\mathrm{eq}}`$ (25)
$`\mathrm{\Psi }_{2,0}`$ $`=`$ $`0`$ (26)
In order to derive eq 24, we have used the fact that, since $`\varphi `$ is the same function of $`X_p`$, $`Y_p`$, and $`Z_p`$ for all $`p`$, $`X_pX_q_{\mathrm{eq}}=Y_pY_q_{\mathrm{eq}}=Z_pZ_q_{\mathrm{eq}}`$. Equation 26 indicates that the inclusion of excluded volume interactions alone is not sufficient to lead to the prediction of a non-zero second normal stress difference. The proper inclusion of hydrodynamic interaction is required.
It is interesting to note, by making use of eq 24, that the mean square radius of gyration at equilibrium, which is defined as ,
$$R_g^2_{\mathrm{eq}}=\frac{1}{N}\underset{\nu =1}{\overset{N}{}}𝑑𝑸_1𝑑𝑸_2\mathrm{}𝑑𝑸_{N1}(𝒓_\nu 𝒓_c)(𝒓_\nu 𝒓_c)\psi _{\mathrm{eq}}$$
(27)
(where, $`𝒓_c`$ is the position of the center of mass), is related to the zero shear rate viscosity by,
$$\eta _{p,0}=\frac{n_\mathrm{p}\zeta }{6}NR_g^2_{\mathrm{eq}}$$
(28)
An alternative expression for the zero shear rate viscosity, which will prove very useful subsequently, can also be obtained from eq 24,
$$\eta _{p,0}=\frac{n_\mathrm{p}\zeta }{12N}\underset{\nu ,\mu =1}{\overset{N}{}}𝒓_{\nu \mu }^2_{\mathrm{eq}}$$
(29)
In order to derive eqs 28 and 29, equations which relate the bead connector vector coordinates to bead position vector coordinates, summarized for example, in Chapter 11 and Table 15.1-1 of Chapter 15 of the text book by Bird et al. , have been used.
The evaluation of the equilibrium averages in eqs 25 and 29, for various values of the parameters in the narrow Gaussian potential, and for various chain lengths $`N`$, have been carried out here with the help of Brownian dynamics simulations. More details of these simulations are given subsequently. In the special case of the extent of excluded volume interactions $`\stackrel{~}{d}`$ going to zero or infinity, we had shown earlier for a Hookean dumbbell model that the values of $`\eta _{p,0}`$ and $`\mathrm{\Psi }_{1,0}`$ remain unchanged from the values that they have in the absence of excluded volume interactions . In the next section, we consider the same limits for the more general case of bead-spring chains of arbitrary (but finite) length.
### 3.2 The Limits $`\stackrel{~}{d}0`$ and $`\stackrel{~}{d}\mathrm{}`$
The average in eq 29 can be evaluated with the distribution function $`\psi _{\mathrm{eq}}(𝑸_1,\mathrm{},𝑸_{N1})`$, or equivalently, with the distribution function $`P_{\mathrm{eq}}(𝒓_{\nu \mu })`$, which is a contracted distribution function for each vector $`𝒓_{\nu \mu }`$, and which is defined by,
$$P_{\mathrm{eq}}(𝒓_{\nu \mu })=𝑑𝑸_1𝑑𝑸_2\mathrm{}𝑑𝑸_{N1}\delta \left(𝒓_{\nu \mu }\underset{j=\mu }{\overset{\nu 1}{}}𝑸_j\right)\psi _{\mathrm{eq}}$$
(30)
We have assumed here, without loss of generality, that $`\nu >\mu `$.
In the Rouse model, as is well known, the equilibrium distribution function is Gaussian,
$$\psi _{\mathrm{eq}}^R(𝑸_1,\mathrm{},𝑸_{N1})=\underset{j=1}{\overset{N1}{}}\left(\frac{H}{2\pi k_\mathrm{B}T}\right)^{3/2}\mathrm{exp}\left(\frac{H}{2k_\mathrm{B}T}𝑸_j𝑸_j\right)$$
(31)
A superscript or subscript ‘$`R`$’ on any quantity will henceforth indicate a quantity defined or evaluated in the Rouse model. The distribution function $`P_{\mathrm{eq}}^R(𝒓_{\nu \mu })`$ can then be evaluated, by substituting eq 31 and the Fourier representation of a $`\delta `$-function, into eq 30 ,
$$P_{\mathrm{eq}}^R(𝒓_{\nu \mu })=\left(\frac{H}{2\pi k_\mathrm{B}T|\nu \mu |}\right)^{3/2}\mathrm{exp}\left(\frac{H}{2|\nu \mu |k_\mathrm{B}T}𝒓_{\nu \mu }^2\right)$$
(32)
The absolute value $`|\nu \mu |`$ indicates that this expression is valid regardless of whether $`\nu `$ or $`\mu `$ is greater. This is another well known result of the Rouse model, namely, at equilibrium, the vector $`𝒓_{\nu \mu }`$ between any two beads $`\mu `$ and $`\nu `$, also obeys a Gaussian distribution.
A similar procedure can be adopted to evaluate $`P_{\mathrm{eq}}(𝒓_{\nu \mu })`$, in the presence of excluded volume interactions, by substituting eq 14 and the Fourier representation of a $`\delta `$-function, into eq 30. We show in appendix A that,
$$\underset{\genfrac{}{}{0pt}{}{\stackrel{~}{d}0}{\mathrm{or},\stackrel{~}{d}\mathrm{}}}{lim}P_{\mathrm{eq}}(𝒓_{\nu \mu })=P_{\mathrm{eq}}^R(𝒓_{\nu \mu })$$
(33)
As a result, for all quantities $`X(𝒓_{\nu \mu })`$, such that the product $`X(𝒓_{\nu \mu })P_{\mathrm{eq}}(𝒓_{\nu \mu })`$ remains bounded for all $`𝒓_{\nu \mu }`$, $`lim_{\genfrac{}{}{0pt}{}{\stackrel{~}{d}0}{\mathrm{or},\stackrel{~}{d}\mathrm{}}}X(𝒓_{\nu \mu })_{\mathrm{eq}}=X(𝒓_{\nu \mu })_{\mathrm{eq}}^R`$. It follows from eq 29 that,
$$\underset{\genfrac{}{}{0pt}{}{\stackrel{~}{d}0}{\mathrm{or},\stackrel{~}{d}\mathrm{}}}{lim}\eta _{p,0}=\eta _{p,0}^R$$
(34)
Thus, the polymer contribution to the viscosities in the limit of zero shear rate, for chains of arbitrary (but finite) length, in (i) the presence of $`\delta `$-function excluded volume interactions, and (ii) the absence of excluded volume interactions (the Rouse model), are identical to each other. Brownian dynamics simulations, details of which are given in the section below, indicate that this is also true for the first normal stress difference coefficients.
## 4 Brownian Dynamics Simulations
The equilibrium averages in eqs 25 and 29, as mentioned above, can be evaluated with the help of Brownian dynamics simulations. As a result, exact numerical predictions of the zero shear rate viscometric functions can be obtained. Brownian dynamics simulations basically involve the numerical solution of the Ito stochastic differential equation that corresponds to the diffusion equation, eq 7. Using standard methods to transcribe a Fokker-Planck equation to a stochastic differential equation, one can show that eq 7 is equivalent to the following system of $`(N1)`$ stochastic differential equations for the connector vectors $`𝑸_j`$,
$$d𝑸_j=\left\{𝜿𝑸_j\frac{1}{\zeta }\underset{k=1}{\overset{N1}{}}A_{jk}\frac{\varphi }{𝑸_k}\right\}dt+\underset{\nu =1}{\overset{N}{}}\sqrt{\frac{2k_\mathrm{B}T}{\zeta }}\overline{B}_{j\nu }d𝑾_\nu $$
(35)
where, $`𝑾_\nu `$ is a $`3N`$ dimensional Wiener process.
A second order predictor-corrector algorithm with time-step extrapolation was used for the numerical solution of eq 35. Steady-state expectations at equilibrium were obtained by setting $`𝜿=0`$, and simulating a single long trajectory. This is justified based on the assumption of ergodicity .
## 5 The Gaussian Approximation
A crucial step in the calculation of the rheological properties predicted by the present model is the evaluation of the complex moments that occur in Kramers expression. The Gaussian approximation—which has previously been shown to be useful in the treatment of hydrodynamic interaction and internal viscosity effects —consists essentially of reducing complex higher order moments to functions of only second order moments by assuming that the non-equilibrium configurational distribution function is a Gaussian distribution, and subsequently, evaluating these second order moments by integrating a time evolution equation.
For the narrow Gaussian potential, the complex moment $`𝑸_k𝑭_\mu ^{(E)}`$, which appears in the quantity $`𝒁`$ on the right hand side of Kramers expression, eq 11, can be rewritten in terms of averages of the form: $`𝑸_k𝑸_nE\left(𝒓_\nu 𝒓_\mu \right)`$. Assuming that $`\psi `$ is a Gaussian distribution of the form,
$$\psi (𝑸_1,\mathrm{},𝑸_{N1},t)=𝒩(t)\mathrm{exp}\left[\frac{1}{2}\underset{j,k}{}𝑸_j(𝝈^1)_{jk}𝑸_k\right]$$
(36)
where, the $`(N1)\times (N1)`$ matrix of tensor components $`𝝈_{jk}`$ (with $`𝝈_{jk}=𝑸_j𝑸_k`$ and $`𝝈_{jk}=𝝈_{kj}^T`$) uniquely characterizes the Gaussian distribution and $`𝒩(t)`$ is the normalization factor, and using general decomposition rules for the moments of a Gaussian distribution , one can show that,
$`𝑸_m𝑸_nE(𝒓_\nu 𝒓_\mu )={\displaystyle \frac{vk_\mathrm{B}T}{\left(2\pi \right)^{3/2}}}{\displaystyle \frac{1}{\sqrt{det\left([\stackrel{~}{d}^2\mathrm{𝟏}+𝒓_{\nu \mu }𝒓_{\nu \mu }]\right)}}}\times `$
$`\left\{𝑸_m𝑸_n𝑸_m𝒓_{\nu \mu }\left[\stackrel{~}{d}^2\mathrm{𝟏}+𝒓_{\nu \mu }𝒓_{\nu \mu }\right]^1𝒓_{\nu \mu }𝑸_n\right\}`$ (37)
The vector $`𝒓_{\nu \mu }`$ also obeys a Gaussian distribution since it is a sum of Gaussian distributed bead connector vectors. As a result, the right hand side of eq 37 involves only second moments, and averages which can be evaluated by Gaussian integrals.
In the Gaussian approximation therefore, Kramers expression for the stress tensor can be rewritten as,
$$𝝉^p=n_\mathrm{p}H\underset{k=1}{\overset{N1}{}}𝝈_{kk}+𝒁+(N1)n_\mathrm{p}k_\mathrm{B}T\mathrm{𝟏}$$
(38)
where,
$$𝒁=\frac{1}{2}z^{}n_\mathrm{p}k_\mathrm{B}T\underset{\genfrac{}{}{0pt}{}{\nu ,\mu =1}{\nu \mu }}{\overset{N}{}}\widehat{𝝈}_{\nu \mu }𝚷(\widehat{𝝈}_{\nu \mu })$$
(39)
In eq 39, the function $`𝚷(\widehat{𝝈}_{\nu \mu })`$ is given by,
$$𝚷(\widehat{𝝈}_{\nu \mu })=\frac{\left[d_{}^{}{}_{}{}^{2}\mathrm{𝟏}+\widehat{𝝈}_{\nu \mu }\right]^1}{\sqrt{det\left([d_{}^{}{}_{}{}^{2}\mathrm{𝟏}+\widehat{𝝈}_{\nu \mu }]\right)}}$$
(40)
with the tensors $`\widehat{𝝈}_{\nu \mu }`$ defined by,
$$\widehat{𝝈}_{\nu \mu }=\widehat{𝝈}_{\nu \mu }^T=\widehat{𝝈}_{\mu \nu }=\frac{H}{k_BT}\underset{j,k=\mathrm{min}(\mu ,\nu )}{\overset{\mathrm{max}(\mu ,\nu )1}{}}𝝈_{jk}$$
(41)
The quantities $`z^{}`$ and $`d^{}`$ are non-dimensional versions of the two parameters, $`v`$ and $`\stackrel{~}{d}`$, which characterize the narrow Gaussian potential. They are defined by,
$$z^{}=v\left(\frac{H}{2\pi k_BT}\right)^{\frac{3}{2}};d^{}=\stackrel{~}{d}\sqrt{\frac{H}{k_BT}}$$
(42)
While $`z^{}`$ measures the strength of the excluded volume interaction, $`d^{}`$ is a measure of the extent of excluded volume interaction.
In the limit of $`d^{}0`$, it is straight forward to see that the tensor $`𝒁`$ becomes isotropic. As a result, the direct contribution to the stress tensor has no influence on the rheological properties of the polymer solution only when a $`\delta `$-function potential is used to represent excluded volume interactions.
All that remains to be done in order to evaluate the stress tensor is to find the components of the covariance matrix $`𝝈_{jk}`$. A system of $`9(N1)^2`$ coupled ordinary differential equations for $`𝝈_{jk}`$ can be obtained from the time evolution equation for the second moments, eq 9. As mentioned earlier, in the presence of excluded volume interactions, eq 9 also involves higher order moments due to the occurrence of the term $`𝒀_{jk}`$, and consequently, it is not in general a closed equation for the second moments. However, these higher order moments can also be reduced to second order moments with the help of the decomposition result, eq 37. In the Gaussian approximation, the second moment equation can therefore be rewritten as,
$$\frac{d}{dt}𝝈_{jk}=𝜿𝝈_{jk}+𝝈_{jk}𝜿^T+\frac{2k_\mathrm{B}T}{\zeta }A_{jk}\mathrm{𝟏}\frac{H}{\zeta }\underset{m=1}{\overset{N1}{}}\left[𝝈_{jm}A_{mk}+A_{jm}𝝈_{mk}\right]+𝒀_{jk}$$
(43)
where,
$$𝒀_{jk}=z^{}\left(\frac{H}{\zeta }\right)\underset{m=1}{\overset{N1}{}}\left[𝝈_{jm}𝚫_{km}+𝚫_{jm}𝝈_{mk}\right]$$
(44)
In eq 44, the $`(N1)\times (N1)`$ matrix of tensor components $`𝚫_{jm}`$ is defined by,
$$𝚫_{jm}=\underset{\mu =1}{\overset{N}{}}\left\{(B_{j+1,m}B_{\mu m})𝚷(\widehat{𝝈}_{j+1,\mu })(B_{jm}B_{\mu m})𝚷(\widehat{𝝈}_{j\mu })\right\}$$
(45)
For any homogeneous flow, rheological properties predicted by the Gaussian approximation can be obtained by appropriately choosing the tensor $`𝜿`$, solving the differential equations, eqs 43, for $`𝝈_{jk}`$, and substituting the result into Kramers expression, eq 38. In this paper, we confine attention to the prediction of linear viscoelastic properties, namely, material functions in small amplitude oscillatory shear flow, and zero shear rate viscometric functions.
Linear viscoelastic properties predicted by the Gaussian approximation can be obtained by constructing a codeformational memory-integral expansion. This is done by expanding the tensors $`𝝈_{jk}`$, in terms of deviations from their isotropic equilibrium solution, up to first order in velocity gradient,
$$𝝈_{jk}=f_{jk}\mathrm{𝟏}+\mathit{ϵ}_{jk}+\mathrm{}$$
(46)
where, the tensors $`f_{jk}\mathrm{𝟏}`$ represent equilibrium second moments in the Gaussian approximation, and the tensors $`\mathit{ϵ}_{jk}`$ are the first order corrections. Since the details of the calculation are not very illuminating, they are given in appendix B, and only the results are summarized below.
The first order codeformational memory-integral expansion derived by the above procedure has the form,
$$𝝉^p=_{\mathrm{}}^t𝑑sG(ts)𝜸_{[1]}(t,s)$$
(47)
where, $`𝜸_{[1]}`$ is the codeformational rate-of-strain tensor , and the memory function $`G(t)`$ is given by eq 85 in appendix B. This expression can now be used to obtain exact expressions for material functions in small amplitude oscillatory shear flow, and for the zero shear rate viscosity and first normal stress difference coefficient in steady shear flow, as shown below.
Small amplitude oscillatory shear flow is characterized by a tensor $`𝜿(t)`$ given by,
$$𝜿(t)=\dot{\gamma }_0\mathrm{cos}\omega t\left(\begin{array}{ccc}0& 1& 0\\ 0& 0& 0\\ 0& 0& 0\end{array}\right)$$
(48)
where, $`\dot{\gamma }_0`$ is the amplitude, and $`\omega `$ is the frequency of oscillations in the plane of flow. The $`yx`$ component of the polymer contribution to the shear stress is then defined by ,
$$\tau _{yx}^p=\eta ^{}(\omega )\dot{\gamma }_0\mathrm{cos}\omega t\eta ^{\prime \prime }(\omega )\dot{\gamma }_0\mathrm{sin}\omega t$$
(49)
where $`\eta ^{}`$ and $`\eta ^{\prime \prime }`$ are the material functions characterizing oscillatory shear flow. They can be represented in a combined form as the complex viscosity, $`\eta ^{}=\eta ^{}i\eta ^{\prime \prime }`$, and $`\eta ^{}`$ can be found, in terms of the relaxation modulus, from the expression
$$\eta ^{}=_0^{\mathrm{}}G(s)e^{i\omega s}𝑑s$$
(50)
Upon substituting eq 85 for the memory function $`G(s)`$ into eq 50, one obtains the predictions of the Gaussian approximation for $`\eta ^{}`$ and $`\eta ^{\prime \prime }`$. These are given by eqs 88 in appendix B.
The zero shear rate viscosity $`\eta _{p,0}`$ and the zero shear rate first normal stress difference coefficient $`\mathrm{\Psi }_{1,0}`$, can be obtained from the complex viscosity in the limit of vanishing frequency,
$$\eta _{p,0}=\underset{\omega 0}{lim}\eta ^{}(\omega );\mathrm{\Psi }_{1,0}=\underset{\omega 0}{lim}\frac{2\eta ^{\prime \prime }(\omega )}{\omega }$$
(51)
The predictions of the zero shear rate viscometric functions by the Gaussian approximation are given by eqs 90 and 91 in appendix B. They are compared with the exact results, eqs 24 and 25, evaluated by Brownian dynamics simulations, in section 7 below.
## 6 First Order Perturbation Expansion
The retarded motion expansion, eq 20, which was obtained by carrying out a perturbation expansion of the distribution function $`\psi `$, in terms of velocity gradients, is valid for arbitrary strength of the excluded volume interaction. In this section, using arguments similar to those in the papers by Öttinger and co-workers , we derive a perturbation expansion of $`𝝉^p`$ in the strength of excluded volume interaction, which is valid for arbitrary shear rates. A significant benefit of the perturbation expansion will be a better understanding of the nature of the Gaussian approximation.
The distribution function $`\psi `$ may be written as $`\psi _R+\psi _z^{}`$, where $`\psi _R`$ is the distribution function in the absence of excluded volume, i.e. in the Rouse model, and $`\psi _z^{}`$ is the correction to first order in the strength of the excluded volume interaction. Since $`\psi _R`$ is Gaussian, it has the form given by eq 36, with $`𝒩(t)`$ replaced by $`𝒩_R(t)`$, and $`𝝈_{jk}`$ replaced by $`𝝈_{jk}^R=𝑸_j𝑸_k_R`$. The second moments $`𝑸_j𝑸_k`$ can then be expanded to first order as, $`𝑸_j𝑸_k=𝝈_{jk}^R+𝑸_j𝑸_k_z^{}`$.
On substituting this expansion into eq 9, and equating terms of like order, the second moment equation can be separated into two equations, a zeroth-order equation and a first-order equation. The zeroth-order equation, which is the second moment equation of the Rouse model, is linear in $`𝝈_{jk}^R`$, and has the following explicit solution,
$$𝝈_{jk}^R=\frac{k_\mathrm{B}T}{H}\left\{\delta _{jk}\mathrm{𝟏}+_{\mathrm{}}^t𝑑s\left[\frac{2H}{\zeta }(ts)A\right]_{jk}𝜸_{[1]}(t,s)\right\}$$
(52)
where, $``$ is an exponential operator. Properties of exponential operators that operate on $`(N1)^2\times (N1)^2`$ matrices are discussed in appendix B. The exponential operators used in this section have similar properties, but operate on $`(N1)\times (N1)`$ matrices.
The first-order second moment equation has the form,
$`{\displaystyle \frac{d}{dt}}𝑸_j𝑸_k_z^{}`$ $`=`$ $`𝜿𝑸_j𝑸_k_z^{}+𝑸_j𝑸_k_z^{}𝜿^T`$ (53)
$``$ $`\left({\displaystyle \frac{H}{\zeta }}\right){\displaystyle \underset{m=1}{\overset{N1}{}}}\left\{𝑸_j𝑸_m_z^{}A_{mk}+A_{jm}𝑸_m𝑸_k_z^{}\right\}+𝒀_{jk}^R`$
where, $`𝒀_{jk}^R`$ is given by eq 10, with the averages on the right hand side evaluated with $`\psi _R`$, i.e., $`\mathrm{}`$ are replaced with $`\mathrm{}_R`$. Since $`\psi _R`$ is a Gaussian distribution, the decomposition result, eq 37, with $`\mathrm{}`$ replaced with $`\mathrm{}_R`$, can be used to reduce $`𝒀_{jk}^R`$ to a function of second moments alone. This leads to,
$$𝒀_{jk}^R=z^{}\left(\frac{H}{\zeta }\right)\underset{m=1}{\overset{N1}{}}\left[𝝈_{jm}^R𝚫_{km}^R+𝚫_{jm}^R𝝈_{mk}^R\right]$$
(54)
In eq 54, $`𝚫_{jm}^R`$ is given by eq 45, with $`𝝈_{jk}`$ replaced by $`𝝈_{jk}^R`$ in the definition of $`\widehat{𝝈}_{\mu \nu }`$ on the right hand side. Equation 53 is a system of linear inhomogeneous ordinary differential equations, whose solution is,
$`𝑸_j𝑸_k_z^{}`$ $`=`$ $`{\displaystyle \underset{r,s=1}{\overset{N1}{}}}{\displaystyle _{\mathrm{}}^t}ds[{\displaystyle \frac{H}{\zeta }}(ts)A]_{jr}𝑬(t,s)𝒀_{rs}^R(s)𝑬{}_{}{}^{T}(t,s)`$ (55)
$`\times `$ $`\left[{\displaystyle \frac{H}{\zeta }}(ts)A\right]_{sk}`$
where, $`𝑬`$ is the displacement gradient tensor .
It is immediately clear from eq 53 that the Gaussian approximation is exact to first order in the strength of excluded volume interaction. This follows from the fact that it could have also been derived by expanding eq 43 to first order in $`z^{}`$. It will be seen later that this property of the Gaussian approximation, is helpful in elucidating its nature.
The first order perturbation expansion for the stress tensor can be obtained by expanding Kramers expression, eq 11, to first order in $`z^{}`$. After reducing complex moments evaluated with the Rouse distribution function to second moments, the stress tensor can be shown to depend only on second moments through the relation,
$$𝝉^p=n_\mathrm{p}H\underset{k=1}{\overset{N1}{}}𝝈_{kk}^Rn_\mathrm{p}H\underset{k=1}{\overset{N1}{}}𝑸_k𝑸_k_z^{}+𝒁^R+(N1)n_\mathrm{p}k_\mathrm{B}T\mathrm{𝟏}$$
(56)
where, $`𝒁^R`$ is given by eq 39, with $`𝝈_{jk}`$ replaced by $`𝝈_{jk}^R`$ in the definition of $`\widehat{𝝈}_{\mu \nu }`$ on the right hand side. Equations 52 and 55 may then be used to derive the following first order perturbation expansion for the stress tensor in arbitrary homogeneous flows,
$`𝝉^p`$ $`=`$ $`n_\mathrm{p}k_BT{\displaystyle \underset{r,s=1}{\overset{N1}{}}}{\displaystyle _{\mathrm{}}^t}ds[{\displaystyle \frac{2H}{\zeta }}(ts)A]_{sr}𝑬(t,s)\{(𝜿(s)+𝜿^T(s))\delta _{rs}`$ (57)
$`+`$ $`\left({\displaystyle \frac{H}{k_\mathrm{B}T}}\right)𝒀_{rs}^R(s)\}𝑬^T(t,s)+𝒁^R+(N1)n_\mathrm{p}k_\mathrm{B}T\mathrm{𝟏}`$
Note that $`𝒁^R`$, the direct contribution to the stress tensor, is isotropic only in the limit $`d^{}0`$. We now consider the special case of steady shear flow, and obtain the zero shear rate viscometric functions.
### 6.1 Steady Shear Flow
In order to obtain zero shear rate viscometric functions correct to first order in $`z^{}`$, it is necessary to evaluate the time integrals in eqs 52 and 57, and to evaluate the quantities $`𝒀_{jk}^R`$ and $`𝒁^R`$ in steady shear flow. The results of these calculations are given below, while the details are given in appendix C.
The excluded volume contributions to the zero shear rate viscometric functions (correct to first order in $`z^{}`$) obtained by setting $`\dot{\gamma }`$ equal to zero in eqs 96 to 98 of appendix C are,
$`\eta _{p,0}^{(E)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\lambda _H^2z^{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{1}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{7/2}}}\left[S_{\mu \nu }^{(0)}S_{\mu \nu }^{(1)}+d_{}^{}{}_{}{}^{2}S_{\mu \nu }^{(1)}\right]`$ (58)
$`\mathrm{\Psi }_{1,0}^{(E)}`$ $`=`$ $`\lambda _H^2z^{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{1}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{7/2}}}\left[\mathrm{\hspace{0.17em}2}S_{\mu \nu }^{(2)}\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)S_{\mu \nu }^{(1)}S_{\mu \nu }^{(1)}\right]`$ (59)
$`\mathrm{\Psi }_{2,0}^{(E)}`$ $`=`$ $`0`$ (60)
where, the time constant $`\lambda _H=(\zeta /4H)`$ has been introduced previously in appendix B, and the quantities $`S_{\mu \nu }^{(m)}`$, which occur in these functions and which were introduced earlier by Öttinger , are defined by,
$$S_{\mu \nu }^{(m)}=2^m\underset{j,k=\mathrm{min}(\mu ,\nu )}{\overset{\mathrm{max}(\mu ,\nu )1}{}}C_{jk}^m$$
(61)
The first order perturbation theory predictions of the zero shear rate viscometric functions given above are compared with exact Brownian dynamics simulations and the Gaussian approximation in section 7. We first, however, examine the role of the parameters $`d^{}`$ and $`z^{}`$ in the present model, by considering the end-to-end vector at equilibrium in the limit of large $`N`$.
### 6.2 The Equilibrium End-to-End Vector For Large Values of $`N`$
The second moment of the end-to-end vector $`𝒓`$ at equilibrium is given by the expression,
$$𝒓𝒓_{\mathrm{eq}}=\underset{j,k=1}{\overset{N1}{}}𝑸_j𝑸_k_{\mathrm{eq}}$$
(62)
For the Rouse model, $`𝝈_{jk}^R|_{\mathrm{eq}}=\left(k_\mathrm{B}T/H\right)\delta _{jk}\mathrm{𝟏}`$. One can show, from eq 55, that the first order correction to the second moments has the following form at equilibrium,
$$𝑸_j𝑸_k_z^{}|_{\mathrm{eq}}=\left(\frac{\zeta }{H}\right)\underset{r,s=1}{\overset{N1}{}}R_{jk,rs}^1𝒀_{rs}^R$$
(63)
where, the $`(N1)^2\times (N1)^2`$ matrix $`R_{jk,mn}`$ is defined by, $`R_{jk,mn}=A_{jm}\delta _{kn}+\delta _{jm}A_{kn}`$, and $`𝒀_{jk}^R`$ has the form,
$$𝒀_{jk}^R|_{\mathrm{eq}}=z^{}\frac{k_\mathrm{B}T}{2\zeta }\underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}\frac{1}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{5/2}}\left(\underset{m,n=1}{\overset{N1}{}}\theta (\mu ,m,n,\nu )R_{jk,mn}\right)\mathrm{𝟏}$$
(64)
Note that the function $`\theta (\mu ,m,n,\nu )`$ has been introduced previously in appendix B (see eq 83). It follows that the mean square end-to-end vector at equilibrium, correct to first order in $`z^{}`$, is given by,
$$𝒓^2_{\mathrm{eq}}=\frac{3k_\mathrm{B}T}{H}\left[(N1)+\frac{1}{2}z^{}\underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}\frac{|\mu \nu |^2}{\left(d_{}^{}{}_{}{}^{2}+|\mu \nu |\right)^{5/2}}\right]$$
(65)
We now consider the limit of a large number of beads, $`N`$. In this limit, the sums in eq 65 can be replaced by integrals. Introducing the following variables,
$$x=\frac{\mu }{N};y=\frac{\nu }{N};d=\frac{d^{}}{\sqrt{N}}$$
(66)
and exploiting the symmetry in $`x`$ and $`y`$, we obtain,
$$𝒓^2_{\mathrm{eq}}=\frac{3k_\mathrm{B}T}{H}N\left\{1+z^{}\sqrt{N}\underset{x>y+c}{_0^1𝑑x_0^x𝑑y}\frac{(xy)^2}{\left(d^2+xy\right)^{5/2}}\right\}$$
(67)
where, $`c`$ is a cutoff parameter of order $`1/N`$ which accounts for the fact that $`\mu \nu `$.
It is clear from eq 67 that the excluded volume corrections to the Rouse end-to-end vector are proportional to $`z^{}\sqrt{N}`$. As a result, the proper perturbation parameter to choose is $`zz^{}\sqrt{N}`$, and not $`z^{}`$. This is a very well known result of the theory of polymer solutions , and indicates that a perturbation expansion in $`z^{}`$ becomes useless for long chains.
The integrals in eq 67 can be performed analytically. However, we are interested only in the form of eq 67, which leads to a very valuable insight. Defining the quantity $`\alpha `$, which is frequently used to represent the swelling of the polymer chain at equilibrium due to excluded volume effects,
$$\alpha ^2=\frac{𝒓^2_{\mathrm{eq}}}{𝒓^2_{\mathrm{eq}}^R}$$
(68)
we can see that in the limit of long chains, $`\alpha =\alpha (z,d)`$. In other words, $`\alpha `$ depends asymptotically only on the parameters $`z`$ and $`d`$, and not on the chain length $`N`$. We shall see later that this insight is very useful in understanding the results of Brownian dynamics simulations, and the Gaussian approximation.
## 7 Equilibrium Swelling and Zero Shear Rate <br>Viscometric Functions
The prediction of equilibrium properties and zero shear rate viscometric functions, by Brownian dynamics simulations, the Gaussian approximation and the first order perturbation expansion, are compared in this section. Before doing so, it is appropriate to note that an equilibrium property, frequently defined in static theories of polymer solutions, namely, the swelling of the radius of gyration, $`\alpha _g^2`$, can be found from the following expression,
$$\alpha _g^2=\frac{R_g^2_{\mathrm{eq}}}{R_g^2_{\mathrm{eq}}^R}=\frac{\eta _{p,0}}{\eta _{p,0}^R}$$
(69)
because of the relation between the radius of gyration and the zero shear rate viscosity, eq 28. Plots of $`\alpha _g^2`$ in this section must, therefore, also be seen as plots of the ratio of the zero shear rate viscosity in the presence of excluded volume interactions to the zero shear rate viscosity in the Rouse model.
Figures 1 to 3 are plots of $`\alpha ^2`$, $`\alpha _g^2`$ and $`(\mathrm{\Psi }_{1,0}/\mathrm{\Psi }_{1,0}^R)`$ versus $`d^{}`$, respectively, at a constant value of $`z^{}=0.5`$, for three different chain lengths, $`N=3`$, $`N=6`$ and $`N=12`$. The squares, circles and triangles are exact results of Brownian dynamics simulations for the narrow Gaussian potential, the dashed lines are the predictions of the Gaussian approximation, and the dot-dashed curves are the predictions of the first order perturbation theory.
In the limit $`d^{}0`$ and for large values of $`d^{}`$, for all the values of chain length $`N`$, the Brownian dynamics simulations reveal that equilibrium and zero shear rate properties tend to Rouse model values. In the case of $`\alpha ^2`$ and $`\alpha _g^2`$, this is expected because of the rigorous result, eq 33. An immediate implication of this behavior is that, for chains of arbitrary but finite length, it is not fruitful to use a $`\delta `$-function potential to represent excluded volume interactions. On the other hand, the figures seem to suggest that a finite range of excluded volume interaction is required to cause an increase from Rouse model values. Both the first order perturbation theory and the Gaussian approximation predict a significant change from Rouse model values in the limit $`d^{}0`$. In the case of a dumbbell model, we were able to rigorously understand the origin of these spurious results . The incorrect term-by-term integration of a series that was not uniformly convergent was found to be the source of the problem. Since first order perturbation theory is the basis for renormalisation group calculations, the invalidity of the $`\delta `$-function potential, which is frequently used in these calculations, is at first sight worrisome. However, we shall see below that the use of a $`\delta `$-function potential may be legitimate when both the limits, $`N\mathrm{}`$ and $`d^{}0`$, are considered.
Figures 1 to 3 show that there exists a threshold value of $`d^{}`$ at which, the results of the Gaussian approximation and the first order perturbation theory, first agree with exact Brownian dynamics simulations. This is consistent with the first order perturbation theory predictions of the end-to-end vector, eq 65, and the viscometric functions, eqs 58 and 59, which reveal that, excluded volume corrections to the Rouse model decrease with increasing values of $`d^{}`$. The threshold value of $`d^{}`$ at which the approximations become accurate, increases as $`N`$ increases. This is a consequence of the well known result, which was demonstrated in section 6, that excluded volume corrections scale as $`z^{}\sqrt{N}`$. Note, however, that the Gaussian approximation always becomes accurate at a smaller threshold value of $`d^{}`$ than the first order perturbation theory. The Gaussian approximation, while being a non-perturbative approximation, is nevertheless, exact to first order in $`z^{}`$. Consequently, as mentioned earlier, it might be considered to consist of an infinite number of higher order terms, and can be expected to be more accurate than the results of the first order perturbation theory.
All the results in Figures 1 to 3 are entirely consistent with the results obtained earlier with a dumbbell model for the polymer molecule. However, in the case of the dumbbell model, the dependence of the quality of the approximations on the chain length $`N`$, could not be examined. The results in Figures 1 to 3 seem to suggest that the Gaussian approximation has a rather limited validity, since for a given value of $`z^{}`$ and $`d^{}`$, it gets progressively worse as the chain length $`N`$ increases. This is in fact not a realistic picture—as is revealed below—when the data is reinterpreted in terms of a different set of coordinates.
Figures 4 to 6 are plots of $`\alpha ^2`$, $`\alpha _g^2`$ and $`(\mathrm{\Psi }_{1,0}/\mathrm{\Psi }_{1,0}^R)`$ versus $`d=(d^{}/\sqrt{N})`$, respectively, at a constant value of $`z=z^{}\sqrt{N}=1.0`$, for three different chain lengths, $`N=6`$, $`N=12`$ and $`N=24`$. Before we discuss the figures, it is appropriate to make a few remarks about the choice of the variables in terms of which the data are displayed. Firstly, we choose $`z`$ as the measure of the strength of excluded volume interaction because perturbation theory clearly reveals that excluded volume corrections scale as $`z^{}\sqrt{N}`$. A constant value of $`z`$, as $`N`$ increases, implies that $`z^{}`$ must simultaneously decrease in order to keep the relative role of excluded volume interactions the same. Secondly, we choose the $`x`$-axis coordinate as $`d=(d^{}/\sqrt{N})`$, because, as was shown in section 6, perturbation theory in the limit of long chains indicates that when the data is displayed in terms of $`d`$ and $`z`$, all the curves should collapse on to a single curve as $`N\mathrm{}`$. The parameter $`d`$ may be considered to be the extent of excluded volume interaction, measured as a fraction of the unperturbed (i.e., Rouse) root mean square end-to-end vector $`\sqrt{𝒓^2_{\mathrm{eq}}^R}`$.
We first discuss the results of exact Brownian dynamics simulations displayed in Figures 4 to 6. As in Figures 1 to 3, all the properties start at Rouse values at $`d=0`$, go through a maximum as $`d`$ increases, and then finally decrease once more towards Rouse values with the continued increase of $`d`$. However, as the chain length increases, the values seem to rise increasingly more rapidly from the Rouse values at $`d=0`$, towards the maximum value. In other words, the slope at the origin seems to be getting steeper as $`N`$ increases. Indeed, the trend of the data leads one to speculate that, in the limit $`N\mathrm{}`$, the data might be singular at $`d=0`$, and consequently legitimize, in this limit, the use of a $`\delta `$-function excluded volume potential. This conclusion is of course only speculative, and needs to be established more rigorously. It has not been possible to examine more closely, with the help of Brownian dynamics simulations, the behavior at small values of $`d`$ for larger values of $`N`$, because of the excessive CPU time that is required. In terms of the non-dimensional time $`t^{}=(t/\lambda _H)`$, for $`N=24`$, a run for two non-dimensional time steps $`\mathrm{\Delta }t^{}=0.1`$ and $`\mathrm{\Delta }t^{}=0.08`$, required roughly 65 hours of CPU time on a SGI Origin2000 with a 195 MHz processor.
When viewed in terms of $`z`$ and $`d`$, the Gaussian approximation is revealed to be far more satisfactory than appeared at first sight in Figures 1 to 3. Indeed, for relatively small values of $`d`$, where the Gaussian approximation is inaccurate at small values of chain length, the Gaussian approximation seems to be becoming more accurate as $`N`$ increases. One might expect that as $`N\mathrm{}`$, the Gaussian approximation becomes accurate for an increasingly larger range of $`d`$ values. However, as will perhaps become clearer with the results discussed below, it appears that, for a given value of $`z`$, there exists a threshold value of $`d`$, below which the Gaussian approximation will be inaccurate, no matter how large a choice of $`N`$ is made. The reason for this behavior is related to a feature that is just noticeable in these figures—curves for different values of $`N`$ appear to be converging to an asymptote. This feature will become much clearer in Figure 7, and will be discussed in greater detail below.
For the sake of clarity, the predictions of the first order perturbation theory are not displayed in Figures 4 to 6. In contrast to the situation in Figures 1 to 3, where the accuracy of the first order perturbation theory becomes progressively worse as $`N`$ increases, its accuracy appears frozen when viewed in terms of $`z`$ and $`d`$. In other words, for different—sufficiently large—values of $`N`$, the first order perturbation theory first becomes accurate at the same threshold value of $`d`$. As in the case of the predictions of the Gaussian approximation, curves for different values of $`N`$ appear to be converging to a common asymptote. This can be seen clearly in Figure 7.
Figure 7 displays plots of $`\alpha ^2`$ versus $`d`$, for different chain lengths, at a constant value of $`z=1`$. It clearly reveals the fact that, both in the Gaussian approximation and in the first order perturbation theory, curves for different values of $`N`$ collapse on to a single curve in the limit $`N\mathrm{}`$. A similar approach to an asymptotic limit is observed as $`N\mathrm{}`$, in the predictions of $`\alpha _g^2`$ and $`(\mathrm{\Psi }_{1,0}/\mathrm{\Psi }_{1,0}^R)`$ by both the approximations, when they are plotted versus $`d`$. The results of Brownian dynamics simulations for $`N=24`$ are also plotted in Figure 7. They indicate that for $`z=1`$, asymptotic values have already been reached by Brownian dynamics simulations, at this relatively small value of $`N`$, for $`d0.3`$. One expects that as $`N`$ increases, asymptotic values will be reached for smaller and smaller values of $`d`$.
The asymptotic values predicted by the first order perturbation theory were obtained by carrying out the integrals in eq 67 analytically. It is worth noting that the convergence to the asymptotic value is quite slow as $`d0`$. On the other hand, the asymptotic values predicted by the Gaussian approximation were obtained by a numerical procedure, as discussed below.
In the Gaussian approximation, calculation of the equilibrium and zero shear rate quantities requires the evaluation of the equilibrium moments $`f_{jk}`$. These are found here, as mentioned in appendix B, by numerical integration of the system of ordinary differential equations, eq 77, using a simple Euler scheme, until steady state is reached (the symmetry in $`j`$ and $`k`$ is used to reduce the number of equations by a factor of two). In addition, the evaluation of $`\eta _{p,0}`$ and $`\mathrm{\Psi }_{1,0}`$ requires the inversion of the $`(N1)^2\times (N1)^2`$ matrix $`\overline{A}_{jk,mn}`$ (see eqs 90 and 91). As a result, the CPU time scales as $`N^6`$, and makes the task of generating data for large values of $`N`$ extremely computationally intensive. We have explored the predictions of chains up to a maximum of $`N=40`$, since for this value of $`N`$, a single run on the SGI Origin2000 computer required approximately 54 hours of CPU time. The asymptotic values in Figure 7 were obtained by the following procedure. For $`z=1`$, equilibrium and zero shear rate data, consisting of property values at different pairs of values $`(d,N)`$, were first compiled by performing a large number of runs for various values of $`N`$ as a function of $`d`$. A specific value of $`d`$ was then chosen, and assuming that the various properties were functions of $`1/\sqrt{N}`$, the values for different $`N`$ were extrapolated to the limit $`N\mathrm{}`$ using a rational function extrapolation algorithm . The choice of $`1/\sqrt{N}`$ is motivated by the fact that the leading correction to the integrals in eq 67 is of order $`1/\sqrt{N}`$ .
The quality of Gaussian approximation as a function of the variable $`z`$, for the quantities $`\alpha _g^2`$ and $`(\mathrm{\Psi }_{1,0}/\mathrm{\Psi }_{1,0}^R)`$, is displayed in Figures 8 and 9, respectively. The behavior of $`\alpha ^2`$ has not been displayed as it is very similar to that of $`\alpha _g^2`$. It is clear from these figures that for a given value of $`N`$, the threshold value of $`d`$ beyond which the Gaussian approximation is accurate increases as $`z`$ increases. On the other hand, as in the case of $`z=1`$, for a fixed value of $`z`$, the accuracy of the Gaussian approximation seems to be increasing with $`N`$, for small values of $`d`$. There is, however, clearly a limit to this accuracy. As $`N`$ becomes large, the results of the exact Brownian dynamics simulations and the Gaussian approximation approach asymptotic values, and consequently, no further change can be noticed with changing $`N`$. Figures 8 and 9 seem to indicate that at small values of $`d`$, while the asymptotic values of Brownian dynamics simulations lie below the asymptotic values of the Gaussian approximation for $`\alpha _g^2`$, the opposite is true for $`(\mathrm{\Psi }_{1,0}/\mathrm{\Psi }_{1,0}^R)`$. A clearer picture would be obtained if it were possible to carry out Brownian dynamics simulations with larger values of $`N`$.
Typical experimental values of $`z`$ lie in the range $`0z15`$ . As we have seen above, for large enough values of $`d`$, the Gaussian approximation remains accurate for a significant fraction of values of $`z`$ in this range. Since corrections to the Rouse model due to excluded volume interactions decrease with increasing shear rate, we can anticipate that the accuracy of the Gaussian approximation will improve as the shear rate increases. Furthermore, since the Gaussian approximation is extremely accurate for the treatment of hydrodynamic interaction effects, and since hydrodynamic interaction is likely to be the dominant effect in a combined theory of hydrodynamic interaction and excluded volume , it is perhaps fair to say that the results obtained so far clearly indicate the practical usefulness of the Gaussian approximation.
## 8 Conclusions
The influence of excluded volume interactions on the linear viscoelastic properties of a dilute polymer solution has been studied with the help of a narrow Gaussian excluded volume potential that acts between pairs of beads in a bead-spring chain model for the polymer molecule. Exact predictions of the model have been obtained by carrying out Brownian dynamics simulations, and approximate predictions have been obtained by two methods—firstly, by carrying out a first order perturbation expansion in the strength of excluded volume interaction, and secondly, by introducing a Gaussian approximation for the configurational distribution function.
The most appropriate way to represent the results of model calculations has been shown to be in terms of a suitably normalized strength of excluded volume interaction $`z`$, and a suitably normalized extent of excluded volume interaction $`d`$. When the results are viewed in terms of these variables, the following conclusions can be drawn:
1. The use of a $`\delta `$-function excluded volume potential (which is the narrow Gaussian excluded volume potential in the limit $`d0`$) is not fruitful for chains with an arbitrary, but finite, number of beads $`N`$, because it leads to the prediction of properties identical to the Rouse model. The narrow Gaussian potential with a finite, non-zero, extent of interaction $`d`$, on the other hand, causes a swelling of the polymer chain at equilibrium, and an increase in the zero shear rate properties from their Rouse model values.
2. Curves for different—but sufficiently large—values of chain length $`N`$, collapse on to a unique asymptotic curve in the limit $`N\mathrm{}`$. The manner in which the results of Brownian dynamics simulations approach the asymptotic behavior indicates that there might be a singularity at $`d=0`$, and consequently, the use of a $`\delta `$-function excluded volume potential might be justified in the limit of infinite chain length.
3. The accuracy of the first order perturbation expansion becomes independent of $`N`$ for large $`N`$. For a given value of $`z`$, there exists a threshold value of $`d`$ beyond which the results of the first order perturbation theory agree with the exact results of Brownian dynamics simulations.
4. As in the case of the first order perturbation expansion, there exists a threshold value of $`d`$ beyond which the results of the Gaussian approximation agree with exact results. For a given value of $`z`$, this threshold value for accuracy is smaller than the threshold value in the first order perturbation theory. The accuracy of the Gaussian approximation decreases with increasing values of $`z`$.
Explicit expressions for the end-to-end vector and the viscometric functions in terms of the model parameters, obtained by carrying out a first order perturbation expansion, enable one to understand the behavior of the Gaussian approximation. This is because the Gaussian approximation is shown here to be exact to first order in $`z`$.
The accuracy of the Gaussian approximation, for a given value of $`z`$ and $`d`$, is expected to improve as the shear rate increases. This follows from the fact that corrections to the Rouse model, due to excluded volume interactions, decrease with increasing shear rate. Viscometric functions at non-zero shear rate predicted by Brownian dynamics simulations, the Gaussian approximation, and the first order perturbation expansion, will be compared in a subsequent publication.
The advantage of using the narrow Gaussian potential to represent excluded volume interactions is that the accuracy of approximate solutions can be assessed by comparison with exact results. This is in contrast with the situation for approximate renormalisation group calculations based on a $`\delta `$-function excluded volume potential, whose accuracy can only be judged by comparison with experiment, or with Monte Carlo simulations based on a different excluded volume potential.
The results obtained here indicate that the Gaussian approximation is an accurate approximation for describing excluded volume interactions, albeit within a certain range of parameter values. Since the usefulness of the Gaussian approximation has already been established for the treatment of hydrodynamic interactions, it is clearly worthwhile to examine the quality of the Gaussian approximation in a model for the combined effects of hydrodynamic interaction and excluded volume.
Acknowledgement. Support for this work through a grant III. 5(5)/98-ET from the Department of Science and Technology, India, is acknowledged. A significant part of this work was carried out while the author was an Alexander von Humboldt fellow at the Department of Mathematics, University of Kaiserslautern, Germany. The author would like to thank Professor H. C. Öttinger for carefully reading the manuscript, and for helpful suggestions. Thanks are also due to the High Performance Computational Facility at the University of Kaiserslautern for providing the use of their computers.
## Appendix A $`P_{\mathrm{eq}}(𝒓_{\nu \mu })`$ in the limit $`\stackrel{~}{d}0`$ or $`\stackrel{~}{d}\mathrm{}`$
Upon substituting eq 14 and the Fourier representation of a $`\delta `$-function into eq 30, and rearranging terms, one obtains,
$$P_{\mathrm{eq}}(𝒓_{\nu \mu })=\frac{1}{(2\pi )^3}𝑑𝒌e^{i𝒓_{\nu \mu }𝒌}\left\{𝒩_{\mathrm{eq}}𝑑𝑸_1\mathrm{}𝑑𝑸_{N1}e^{\left[\left(\varphi /k_\mathrm{B}T\right)+i\left(_{j=\mu }^{\nu 1}𝑸_j\right)𝒌\right]}\right\}$$
(70)
We now consider the integral within braces on the right hand side of eq 70, and take up the integration over the bead connector vector $`𝑸_1`$. Separating out all the terms containing the vector $`𝑸_1`$, we can rewrite this integral as,
$`𝒩_{\mathrm{eq}}{\displaystyle }d𝑸_2\mathrm{}d𝑸_{N1}\mathrm{exp}[{\displaystyle \frac{H}{2k_\mathrm{B}T}}{\displaystyle \underset{j=2}{\overset{N1}{}}}𝑸_j^2i\left({\displaystyle \underset{j=\mu }{\overset{\nu 1}{}}}(1\delta _{1j})𝑸_j\right)𝒌`$
$`{\displaystyle \frac{1}{2k_\mathrm{B}T}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\alpha ,\beta =2}{\alpha \beta }}{\overset{N}{}}}E(𝒓_\alpha 𝒓_\beta )\left]\right\{{\displaystyle }d𝑸_1\mathrm{exp}[{\displaystyle \frac{H}{2k_\mathrm{B}T}}𝑸_1^2i\delta _{1\mu }𝑸_1𝒌`$
$`{\displaystyle \frac{1}{k_\mathrm{B}T}}[E(𝒓_1𝒓_2)+E(𝒓_1𝒓_3)+\mathrm{}+E(𝒓_1𝒓_N)]]\}`$ (71)
where, a typical term of the excluded volume potential contribution to the $`𝑸_1`$ integral, has the form,
$`E\left(𝒓_1𝒓_\beta \right)`$ $`=`$ $`{\displaystyle \frac{vk_\mathrm{B}T}{[2\pi \stackrel{~}{d}^2]^{\frac{3}{2}}}}\mathrm{exp}\{{\displaystyle \frac{1}{2\stackrel{~}{d}^2}}(𝑸_1^2+2𝑸_1𝒓_{\beta 2}+𝑸_2^2+2𝑸_2𝒓_{\beta 3}+\mathrm{}`$
$`+`$ $`𝑸_{\beta 2}^2+2𝑸_{\beta 2}𝒓_{\beta ,\beta 1}+𝑸_{\beta 1}^2)\}`$
We now convert the $`𝑸_1`$ integral into spherical coordinates. In order to do so, we need to choose a reference vector to fix a direction in space. In the $`𝑸_1`$ integration, all the other vectors, $`𝑸_2,\mathrm{},𝑸_{N1}`$ and $`𝒌`$ are fixed. Without loss of generality, we choose the fixed vector as $`𝑸_2`$, denote its direction as the $`z`$ direction, and choose, in the plane perpendicular to $`𝑸_2`$, an arbitrary pair of orthogonal directions as the $`x`$ and $`y`$ axes. Let, $`\theta _1`$, $`\theta _{\beta 2}`$, and $`\theta _k`$ represent the angles that the vectors $`𝑸_1`$, $`𝒓_{\beta 2}`$ and $`𝒌`$ make with the $`z`$ direction, respectively. Similarly, let $`\varphi _1`$, $`\varphi _{\beta 2}`$, and $`\varphi _k`$ represent the angles that the projections of these vectors on the $`xy`$ plane, make with the $`x`$ direction. Then,
$$𝑸_1𝒓_{\beta 2}=Q_1r_{\beta 2}F_{\beta 2}(\theta _1,\varphi _1)$$
where, $`Q_1`$ and $`r_{\beta 2}`$ represent the magnitudes of $`𝑸_1`$ and $`𝒓_{\beta 2}`$, respectively, and,
$$F_{\beta 2}(\theta _1,\varphi _1)=\mathrm{sin}\theta _1\mathrm{sin}\theta _{\beta 2}(\mathrm{cos}\varphi _1\mathrm{cos}\varphi _{\beta 2}+\mathrm{sin}\varphi _1\mathrm{sin}\varphi _{\beta 2})+\mathrm{cos}\theta _1\mathrm{cos}\theta _{\beta 2}$$
Defining the function $`F_k(\theta _1,\varphi _1)`$ similarly, we can rewrite the $`𝑸_1`$ integral in expression 71, in terms of spherical coordinates as,
$`I_{Q_1}={\displaystyle _0^{\mathrm{}}}dQ_1{\displaystyle _0^{2\pi }}d\theta _1{\displaystyle _0^\pi }d\varphi _1Q_1^2\mathrm{sin}\theta _1\mathrm{exp}[{\displaystyle \frac{H}{2k_\mathrm{B}T}}Q_1^2i\delta _{1\mu }Q_1kF_k(\theta _1,\varphi _1)]\times `$
$`\mathrm{exp}\{{\displaystyle \frac{1}{k_\mathrm{B}T}}[{\displaystyle \frac{vk_\mathrm{B}T}{\left(2\pi \stackrel{~}{d}^2\right)^{3/2}}}\{e^{\frac{1}{2\stackrel{~}{d}^2}Q_1^2}+e^{\frac{1}{2\stackrel{~}{d}^2}\left(Q_1^2+2Q_1r_{32}F_{32}(\theta _1,\varphi _1)\right)}e^{\frac{1}{2\stackrel{~}{d}^2}𝑸_2^2}`$
$`+\mathrm{}+e^{\frac{1}{2\stackrel{~}{d}^2}\left(Q_1^2+2Q_1r_{N2}F_{N2}(\theta _1,\varphi _1)\right)}e^{\frac{1}{2\stackrel{~}{d}^2}\left(𝑸_2^2+2𝑸_2𝒓_{N3}+\mathrm{}+𝑸_{N1}^2\right)}\}]\}`$ (72)
For $`Q_1=0`$, the integrand is identically zero. For $`Q_10`$, in the limit $`\stackrel{~}{d}0`$ or $`\stackrel{~}{d}\mathrm{}`$, the integrand tends to,
$$Q_1^2\mathrm{sin}\theta _1\mathrm{exp}\left\{\frac{H}{2k_\mathrm{B}T}Q_1^2i\delta _{1\mu }Q_1kF_k(\theta _1,\varphi _1)\right\}$$
The integrand is also a bounded function of $`Q_1`$ for all values $`\stackrel{~}{d}`$.
An argument similar to the one above can be carried out for each of the remaining integrations over $`𝑸_2,\mathrm{},𝑸_{N1}`$. It follows that,
$`\underset{\genfrac{}{}{0pt}{}{\stackrel{~}{d}0}{\mathrm{or},\stackrel{~}{d}\mathrm{}}}{lim}{\displaystyle 𝑑𝑸_1\mathrm{}𝑑𝑸_{N1}\mathrm{exp}\left\{\left(\frac{\varphi }{k_\mathrm{B}T}\right)i\left(\underset{j=\mu }{\overset{\nu 1}{}}𝑸_j\right)𝒌\right\}}`$
$`={\displaystyle 𝑑𝑸_1\mathrm{}𝑑𝑸_{N1}\mathrm{exp}\left\{\left(\frac{H}{2k_\mathrm{B}T}\right)\underset{j=1}{\overset{N1}{}}𝑸_j^2i\left(\underset{j=\mu }{\overset{\nu 1}{}}𝑸_j\right)𝒌\right\}}`$ (73)
With regard to the normalization factor $`𝒩_{\mathrm{eq}}`$, since,
$$𝒩_{\mathrm{eq}}=\left[𝑑𝑸_1\mathrm{}𝑑𝑸_{N1}\mathrm{exp}\left(\frac{\varphi }{k_\mathrm{B}T}\right)\right]^1$$
(74)
we can show, by adopting a procedure similar to that above that,
$$\underset{\genfrac{}{}{0pt}{}{\stackrel{~}{d}0}{\mathrm{or},\stackrel{~}{d}\mathrm{}}}{lim}𝒩_{\mathrm{eq}}=𝒩_{\mathrm{eq}}^R$$
(75)
As a result, we have established that,
$`\underset{\genfrac{}{}{0pt}{}{\stackrel{~}{d}0}{\mathrm{or},\stackrel{~}{d}\mathrm{}}}{lim}P_{\mathrm{eq}}(𝒓_{\nu \mu })`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle 𝑑𝒌e^{i𝒓_{\nu \mu }𝒌}\left\{𝑑𝑸_1\mathrm{}𝑑𝑸_{N1}\psi _{\mathrm{eq}}^Re^{i\left[_{j=\mu }^{\nu 1}𝑸_j\right]𝒌}\right\}}`$ (76)
$`=`$ $`P_{\mathrm{eq}}^R(𝒓_{\nu \mu })`$
## Appendix B Codeformational Memory-Integral Expansion
Upon expanding the tensors $`𝝈_{jk}`$ in the manner displayed in eq 46, substituting the expansion into the second moment equation, eq 43, and separating the resultant equation into equations for each order in the velocity gradient, the following two equations are obtained, Equilibrium:
$$\frac{d}{dt}f_{jk}=\frac{2k_\mathrm{B}T}{\zeta }A_{jk}\left(\frac{H}{\zeta }\right)\underset{m=1}{\overset{N1}{}}\left[f_{jm}(A_{mk}z^{}\mathrm{\Delta }_{km}^{(0)})+(A_{jm}z^{}\mathrm{\Delta }_{jm}^{(0)})f_{mk}\right]$$
(77)
where,
$$\mathrm{\Delta }_{jm}^{(0)}=\underset{\mu =1}{\overset{N}{}}\left[\frac{(B_{j+1,m}B_{\mu m})}{\left(d_{}^{}{}_{}{}^{2}+\widehat{f}_{j+1,\mu }\right)^{5/2}}\frac{(B_{jm}B_{\mu m})}{\left(d_{}^{}{}_{}{}^{2}+\widehat{f}_{j\mu }\right)^{5/2}}\right]$$
(78)
with the quantities $`\widehat{f}_{\nu \mu }`$ given by,
$$\widehat{f}_{\nu \mu }=\left(\frac{H}{k_BT}\right)\underset{j,k=\mathrm{min}(\mu ,\nu )}{\overset{\mathrm{max}(\mu ,\nu )1}{}}f_{jk}$$
(79)
First Order:
$$\frac{d}{dt}\mathit{ϵ}_{jk}=(𝜿+𝜿^T)f_{jk}\left(\frac{H}{\zeta }\right)\underset{m,n=1}{\overset{N1}{}}\overline{A}_{jk,mn}\mathit{ϵ}_{mn}$$
(80)
where,
$`\overline{A}_{jk,mn}=(A_{jm}\delta _{kn}+\delta _{jm}A_{kn})z^{}(\mathrm{\Delta }_{jm}^{(0)}\delta _{kn}+\delta _{jm}\mathrm{\Delta }_{kn}^{(0)})`$
$`+z^{}\left({\displaystyle \frac{H}{k_BT}}\right){\displaystyle \underset{p=1}{\overset{N1}{}}}\left[f_{jp}\mathrm{\Delta }_{kp,mn}^{(1)}+\mathrm{\Delta }_{jp,mn}^{(1)}f_{pk}\right]`$ (81)
with the quantities $`\mathrm{\Delta }_{jp,mn}^{(1)}`$ given by,
$$\mathrm{\Delta }_{jp,mn}^{(1)}=\underset{\mu =1}{\overset{N}{}}\left[\frac{(B_{j+1,p}B_{\mu p})\theta (\mu ,m,n,j+1)}{\left(d_{}^{}{}_{}{}^{2}+\widehat{f}_{j+1,\mu }\right)^{7/2}}\frac{(B_{jp}B_{\mu p})\theta (\mu ,m,n,j)}{\left(d_{}^{}{}_{}{}^{2}+\widehat{f}_{j\mu }\right)^{7/2}}\right]$$
(82)
The function $`\theta (\mu ,m,n,\nu )`$ has been introduced previously in the treatment of hydrodynamic interaction . It is unity if $`m`$ and $`n`$ lie between $`\mu `$ and $`\nu `$, and zero otherwise,
$$\theta (\mu ,m,n,\nu )=\{\begin{array}{cc}1\hfill & \text{if }\mu m,n<\nu \text{ or }\nu m,n<\mu \hfill \\ \multicolumn{2}{c}{}\\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(83)
Introducing new indices for the pairs $`(j,k)`$ and $`(m,n)`$, the quantity $`\overline{A}_{jk,mn}`$ may be considered an $`(N1)^2\times (N1)^2`$ matrix. The inverse can then be defined in the following manner,
$$\underset{r,s=1}{\overset{N1}{}}\overline{A}_{jk,rs}^1\overline{A}_{rs,mn}=1\mathrm{I}_{jk,mn}$$
(84)
where, $`1\mathrm{I}`$ is the $`(N1)^2\times (N1)^2`$ unit matrix $`1\mathrm{I}_{jk,mn}=\delta _{jm}\delta _{kn}`$.
In the equilibrium second moment equation, eq 77, the term $`(df_{jk}/dt)`$ on the left hand side, is identically zero. It is retained here, however, to indicate that the equation is solved for $`f_{jk}`$ by numerical integration of the ODE’s until steady state is reached.
Upon integrating eq 80 with respect to time, and substituting the solution into eq 38, we finally obtain the expression, eq 47, for the first order codeformational memory-integral expansion, where, the memory function $`G(t)`$ is given by,
$$G(t)=\underset{j,k=1}{\overset{N1}{}}\underset{m,n=1}{\overset{N1}{}}f_{jm}_{jk}\left[\frac{t}{\lambda _H}\stackrel{~}{A}\right]_{mk,nn}$$
(85)
In eq 85, $`\lambda _H=(\zeta /4H)`$ is the familiar time constant, $`_{jk}`$ is defined by,
$$_{jk}=n_\mathrm{p}H\left[\delta _{jk}\frac{1}{2}z^{}\underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}\frac{d_{}^{}{}_{}{}^{2}}{\left(d_{}^{}{}_{}{}^{2}+\widehat{f}_{\nu \mu }\right)^{7/2}}\theta (\nu ,j,k,\mu )\right]$$
(86)
and the quantity $`\stackrel{~}{A}_{jk,mn}`$ is given by,
$$\stackrel{~}{A}_{jk,mn}=\underset{r,s=1}{\overset{N1}{}}\frac{1}{4}f_{jr}^1\overline{A}_{rk,sn}f_{sm}$$
(87)
The exponential operator $`[M]`$ maps one matrix into another according to:
$$[M]_{jk,mn}=1\mathrm{I}_{jk,mn}+M_{jk,mn}+\frac{1}{2!}\underset{r,s=1}{\overset{N1}{}}M_{jk,rs}M_{rs,mn}+\mathrm{}$$
and has the useful properties,
$$\frac{d}{dt}[Mt]_{jk,mn}=\underset{r,s=1}{\overset{N1}{}}M_{jk,rs}[Mt]_{rs,mn}=\underset{r,s=1}{\overset{N1}{}}[Mt]_{jk,rs}M_{rs,mn}$$
$$\underset{r,s=1}{\overset{N1}{}}[aM]_{jk,rs}[b\mathrm{\hspace{0.17em}1}\mathrm{I}]_{rs,mn}=[aM+b\mathrm{\hspace{0.17em}1}\mathrm{I}]_{jk,mn}$$
for arbitrary constants $`a`$ and $`b`$.
As mentioned in section 5, once the memory function $`G(s)`$ is obtained, one can obtain the material functions in small amplitude oscillatory shear flow, and the zero shear rate viscometric functions. Following the procedure outlined in section 5, we can show that,
$`\eta ^{}(\omega )`$ $`=`$ $`\lambda _H{\displaystyle \underset{j,k=1}{\overset{N1}{}}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}{\displaystyle \underset{r,s=1}{\overset{N1}{}}}f_{jk}_{jm}\mathrm{\Phi }_{km,rs}^1\stackrel{~}{A}_{rs,nn}`$
$`\eta ^{\prime \prime }(\omega )`$ $`=`$ $`\lambda _H^2\omega {\displaystyle \underset{j,k=1}{\overset{N1}{}}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}f_{jk}_{jm}\mathrm{\Phi }_{km,nn}^1`$ (88)
where,
$$\mathrm{\Phi }_{jk,mn}=\underset{r,s=1}{\overset{N1}{}}\left[\stackrel{~}{A}_{jk,rs}\stackrel{~}{A}_{rs,mn}+\lambda _H^2\omega ^21\mathrm{I}_{jk,mn}\right]$$
(89)
Using the relations between the zero shear rate viscometric functions and $`\eta ^{}`$ and $`\eta ^{\prime \prime }`$ (eqs 51), one can show that,
$`\eta _{p,0}`$ $`=`$ $`4\lambda _H{\displaystyle \underset{j,k=1}{\overset{N1}{}}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}_{jk}\overline{A}_{jk,mn}^1f_{mn}`$ (90)
$`\mathrm{\Psi }_{1,0}`$ $`=`$ $`32\lambda _H^2{\displaystyle \underset{j,k=1}{\overset{N1}{}}}{\displaystyle \underset{m,n=1}{\overset{N1}{}}}{\displaystyle \underset{r,s=1}{\overset{N1}{}}}_{jk}\overline{A}_{jk,mn}^1\overline{A}_{mn,rs}^1f_{rs}`$ (91)
## Appendix C Viscometric functions correct to first order in $`z^{}`$
The first step in calculating the first order excluded volume corrections to the Rouse viscometric functions, as mentioned earlier, is to evaluate the time integrals in eqs 52 and 57. These time integrals can be carried out by using the forms of the tensors $`𝜸_{[1]}`$ and $`𝑬`$ in steady shear flow, tabulated in reference 12. One can show that the expression for the second moment $`𝝈_{jk}^R`$, which is required to explicitly evaluate all the averages carried out with the Rouse distribution function $`\psi _R`$, is given by,
$$𝝈_{jk}^R=\frac{k_\mathrm{B}T}{H}\left\{\delta _{jk}\mathrm{𝟏}+2\lambda _HC_{jk}\left(𝜿+𝜿^T\right)+8\lambda _H^2C_{jk}^2\left(𝜿𝜿^T\right)\right\}$$
(92)
while the stress tensor in steady shear flow has the form,
$`𝝉^p`$ $`=`$ $`n_\mathrm{p}k_BT{\displaystyle \underset{j=1}{\overset{N1}{}}}\left[2\lambda _HC_{jj}\left(𝜿+𝜿^T\right)+8\lambda _H^2C_{jj}^2\left(𝜿𝜿^T\right)\right]`$ (93)
$``$ $`n_\mathrm{p}H{\displaystyle \underset{j,k=1}{\overset{N1}{}}}\{2\lambda _HC_{kj}𝒀_{jk}^R+4\lambda _H^2C_{kj}^2[𝜿𝒀_{jk}^R+𝒀_{jk}^R𝜿^T]`$
$`+`$ $`16\lambda _H^3C_{kj}^3[𝜿𝒀_{jk}^R𝜿^T]\}+𝒁^R`$
A similar expression, without the $`𝒁^R`$ term, has been derived by Öttinger in his renormalisation group treatment of excluded volume effects—within the framework of polymer kinetic theory—using a $`\delta `$-function excluded volume potential.
The next step is to explicitly evaluate the tensors $`𝒀_{jk}^R`$ and $`𝒁^R`$, in terms of the velocity gradient $`𝜿`$ and the Kramers matrix $`C_{kj}`$, by using eq 92 for $`𝝈_{jk}^R`$. The resultant expressions are then substituted into eq 93, and after a lengthy calculation, the following explicit expression for the excluded volume contribution to the stress tensor, correct to first order in $`z^{}`$, is obtained,
$`𝝉_{(E)}^p`$ $`=`$ $`{\displaystyle \frac{1}{8\lambda _H}}n_\mathrm{p}k_BTz^{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{1}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{5/2}e_{\mu \nu }(\dot{\gamma })^{3/2}}}\{(\alpha _{\mu \nu }^{(0)}\beta _{\mu \nu }^{(0)})\mathrm{𝟏}`$ (94)
$`+`$ $`\left(\alpha _{\mu \nu }^{(1)}\beta _{\mu \nu }^{(1)}\right)\left(𝜿+𝜿^T\right)+\left(\alpha _{\mu \nu }^{(2)}\beta _{\mu \nu }^{(2)}\right)\left(𝜿𝜿^T\right)+\left(\alpha _{\mu \nu }^{(3)}\beta _{\mu \nu }^{(3)}\right)\left(𝜿^T𝜿\right)`$
$`+`$ $`\left(\alpha _{\mu \nu }^{(4)}\beta _{\mu \nu }^{(4)}\right)\left(𝜿𝜿^T𝜿\right)+\left(\alpha _{\mu \nu }^{(5)}\beta _{\mu \nu }^{(5)}\right)\left(𝜿^T𝜿𝜿^T\right)`$
$`+`$ $`(\alpha _{\mu \nu }^{(6)}\beta _{\mu \nu }^{(6)})(𝜿𝜿^T𝜿𝜿^T)\}`$
where, the functions $`\alpha _{\mu \nu }^{(j)}`$ and $`\beta _{\mu \nu }^{(j)};(j=0,1,\mathrm{},6)`$, which represent the indirect and direct contributions respectively, are given in Tables 1 and 2, and the function $`e_{\mu \nu }(\dot{\gamma })`$ is defined by,
$$e_{\mu \nu }(\dot{\gamma })=1+\frac{\lambda _H^2\dot{\gamma }^2}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^2}\left[2\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)S_{\mu \nu }^{(2)}S_{\mu \nu }^{(1)}S_{\mu \nu }^{(1)}\right]$$
(95)
Equation 94 for the stress tensor can then be used to find the excluded volume contributions to the viscometric functions, correct to first order in $`z^{}`$, by using the definitions in eqs 22,
$`\eta _p^{(E)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\lambda _H^2z^{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{1}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{7/2}e_{\mu \nu }(\dot{\gamma })^{3/2}}}[S_{\mu \nu }^{(0)}S_{\mu \nu }^{(1)}+\lambda _H^2\dot{\gamma }^2S_{\mu \nu }^{(1)}S_{\mu \nu }^{(2)}`$ (96)
$`+`$ $`d_{}^{}{}_{}{}^{2}S_{\mu \nu }^{(1)}]`$
$`\mathrm{\Psi }_1^{(E)}`$ $`=`$ $`\lambda _H^2z^{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{1}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{7/2}e_{\mu \nu }(\dot{\gamma })^{3/2}}}`$ (97)
$`\times `$ $`\left[\mathrm{\hspace{0.17em}2}S_{\mu \nu }^{(2)}\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)S_{\mu \nu }^{(1)}S_{\mu \nu }^{(1)}+\lambda _H^2\dot{\gamma }^2\left(3S_{\mu \nu }^{(2)}S_{\mu \nu }^{(2)}2S_{\mu \nu }^{(3)}S_{\mu \nu }^{(1)}\right)\right]`$
$`\mathrm{\Psi }_2^{(E)}`$ $`=`$ $`{\displaystyle \frac{1}{8\lambda _H}}n_\mathrm{p}k_BTz^{}{\displaystyle \underset{\genfrac{}{}{0pt}{}{\mu ,\nu =1}{\mu \nu }}{\overset{N}{}}}{\displaystyle \frac{\alpha _{\mu \nu }^{(3)}\beta _{\mu \nu }^{(3)}}{\left(d_{}^{}{}_{}{}^{2}+S_{\mu \nu }^{(0)}\right)^{5/2}e_{\mu \nu }(\dot{\gamma })^{3/2}}}=0`$ (98)
These expressions have been derived earlier by Öttinger , in an arbitrary number of space dimensions, in the limit $`d^{}0`$. It is clear from eq 98 that the second normal stress difference coefficient is zero because the indirect and direct excluded volume contributions cancel each other out. When $`d^{}0`$, however, both the quantities $`\alpha _{\mu \nu }^{(3)}`$ and $`\beta _{\mu \nu }^{(3)}`$ are identically zero. |
warning/0002/quant-ph0002064.html | ar5iv | text | # Robust unravelings for resonance fluorescence
## I Introduction
Some states of open quantum systems are more robust than others. That is, they are less perturbed by the system dynamics. This fact has been the subject of a long-running and active research program . Recently, one of us and Vaccaro have introduced into this program a formalism with a number of distinctive features . This formalism consists of finding the maximally robust unraveling (MRU) for the open quantum system. Its introduction was motivated by a desire to better understand the rich dynamics of open quantum systems in general , and that of the atom laser in particular .
The use of the term “robust unraveling” rather than “robust state” encapsulates two of the distinctive features of the work of Refs. . An unraveling is a way of measuring the environment of an open quantum system such that the system state can be described by a pure state undergoing stochastic evolution. This is always possible in principle if the unmonitored system obeys a Markovian master equation, which we will assume to have a unique stationary state. In the long time limit, the “unraveled” system will be in a pure state drawn at random from a particular ensemble of pure states defined by the unraveling. It is from the consideration of such an ensemble of pure states that the two distinctive features of our approach are met. The first is that it is not individual pure states whose robustness are to be calculated, but rather a whole ensemble of pure states, the average robustness of which is calculated. The second is that the pure states in this ensemble are physically realizable in the sense that they are the states of the system known to an experimenter using the appropriate measurement scheme on the system’s environment.
As noted above, the principle application of the maximally robust unraveling formalism has been to a model for an atom laser (a continuously damped and replenished gaseous Bose-Einstein condensate). This work is a specialized application in two ways. First, the system itself has a high excitation number and so is a quantum system in the classical limit. Second, only a subset of the set of all possible unravelings was considered. This subset (which is still infinite) contains those unravelings that lead to continuous and Markovian evolution of the system state vector . This restriction was necessary to make the problem tractable and was justified by the classicality of the system.
In this work we apply the formalism of MRU to a system with no classical limit (in the usual sense at least), a resonantly-driven fluorescent two-level atom. This system is one of the canonical examples of an open quantum system, and has surprisingly complex dynamics for its size. It is therefore worth investigating in its own right. But, even more importantly, it is simple enough that the maximally robust unraveling can be found analytically. It turns out that this MRU is neither continuous nor Markovian. This enables us to investigate the question of how closely one can approximate the ensemble of this MRU if one is restricted to considering continuous Markovian unravelings. The answer to this question has implications for the general usefulness of the MRU formalism, since for typical systems it would be necessary to impose this restriction in order to make the formalism practical.
The structure of this paper is as follows. In Sec. II we briefly review the MRU formalism. In Sec. III we introduce the two-level atom model and derive an expression for the ensemble average survival probability (which is used to quantify the robustness) in terms of moments of the ensemble of state vectors. In Sec. IV we look at one simple (but not maximally robust) unraveling, that resulting from direct detection of the atom’s fluorescence, for comparison with other, more robust, unravelings. In Sec. V we present the most robust unraveling and its ensemble of (in this case, just two) state vectors. In Sec. VI we find the most robust unraveling from within the set of continuous Markovian unravelings. We compare this ensemble to the MRU of Sec. V in Sec. VII. We conclude with a discussion of the implications of our results in Sec. VIII.
## II Maximally Robust Unravelings
### A The Master Equation
Open quantum systems generally become entangled with their environment, and this causes their state to become mixed. In many cases, the system will reach an equilibrium mixed state in the long time limit. This is the sort of system for which our approach to robustness, of finding the maximally robust unraveling (MRU), can be applied without modification.
If the system is weakly coupled to the environmental reservoir, and many modes of the reservoir are roughly equally affected by the system, then one can make the Born and Markov approximations in describing the effect of the environment on the system . Tracing over (that is, ignoring) the state of the environment leads to a Markovian evolution equation for the state matrix $`\rho `$ of the system, known as a quantum master equation. The most general form of the quantum master equation that is mathematically valid is the Lindblad form
$$\dot{\rho }=i[H,\rho ]+\underset{\mu =1}{\overset{M}{}}𝒟[c_\mu ]\rho \rho ,$$
(1)
where for arbitrary operators $`A`$ and $`B`$,
$$𝒟[A]BABA^{}(A^{}AB+BA^{}A)/2.$$
(2)
If the master equation has a unique stationary state (as we will assume it does), then that is defined by
$$\rho _{\mathrm{ss}}=0.$$
(3)
This assumption requires that $``$ be time-independent. In many quantum optical situations, such as resonance fluorescence, one is only interested in the dynamics in the interaction picture, in which the free evolution at optical frequencies is removed from the state matrix. Indeed, if one treats the driving field as classical, as we will do, it is necessary to move into such an interaction picture in order to obtain a time-independent Liouvillian superoperator $``$.
The stationary state matrix $`\rho _{\mathrm{ss}}`$ can be expressed as an ensemble of pure states as follows:
$$\rho _{\mathrm{ss}}=\underset{k}{}\mathrm{}_k|\psi _k\psi _k|,$$
(4)
where the $`|\psi _k`$ are normalized state vectors and the $`\mathrm{}_k`$ are positive weights summing to unity. The (possibly infinite) set of ordered pairs,
$$E=\{(|\psi _k,\mathrm{}_k):k\},$$
(5)
we will call an ensemble $`E`$ of pure states. Note that there is no restriction that the states be mutually orthogonal. This means that there are continuously infinitely many ensembles $`E`$ that represent $`\rho _{\mathrm{ss}}`$. The aim of finding the MRU is to find the “best” or “most natural” representation for $`\rho _{\mathrm{ss}}`$.
### B Unravelings
As explained in the Introduction above, the first criterion for our most natural ensemble is that it be physically realizable by monitoring the environment of the system. In the situation where a Markovian master equation can be derived, it is possible (in principle) to continually measure the state of the environment on a time scale large compared to the reservoir correlation time but small compared to the response time of the system. This effectively continuous measurement is what we mean by “monitoring”. In such systems, monitoring the environment does not disrupt the system–reservoir coupling and the system will continue to evolve according to the master equation if one ignores the results of the monitoring.
By contrast, if one does take note of the results of monitoring the environment, then the system will no longer obey the master equation. Because the system–reservoir coupling causes the reservoir to become entangled with the system, measuring the former’s state produces information about the latter’s state. This will tend to undo the increase in the mixedness of the system’s state caused by the coupling.
If one is able to make perfect rank-one projective (i.e. von Neumann) measurements of the reservoir state, the system state will usually be collapsed towards a pure state. However this is not a process that itself can be described by projective measurements on the system, because the system is not being directly measured. Rather, the monitoring of the environment leads to a gradual (on average) decrease in the system’s entropy.
If the system is initially in a pure state then, under perfect monitoring of its environment, it will remain in a pure state. Then the effect of the monitoring is to cause the system to change its pure state in a stochastic and (in general) nonlinear way. Such evolution has been called a quantum trajectory , and can be described by a nonlinear stochastic Schrödinger equation for the system state vector . The nonlinearity and stochasticity are present because they are a fundamental part of measurement in quantum mechanics.
On average, the system still obeys the master equation. That is, if the increment in $`|\psi `$ under the SSE is $`|d\psi `$ then
$$\mathrm{E}[|d\psi \psi |+|\psi d\psi |+|d\psi d\psi |]=\mathrm{E}[|\psi \psi |].$$
(6)
Here $`\mathrm{E}`$ denotes the ensemble average with respect to the stochasticity of the SSE. This stochasticity is evidenced by the necessity of retaining the Itô term $`|d\psi d\psi |`$ .
Because the ensemble average of the system still obeys the master equation, the stochastic Schrödinger equation) is said to unravel the master equation . It is now well-known that there are many (in fact continuously many) different unravelings for a given master equation , corresponding to different ways of monitoring the environment.
Each unraveling $`𝒰`$ gives rise to an ensemble of pure states
$$E^𝒰=\{(|\psi _k^𝒰,\mathrm{}_k^𝒰):k\},$$
(7)
where $`|\psi _k^𝒰`$ are the possible pure states of the system at steady state, and $`\mathrm{}_k^𝒰`$ are their weights. For master equations with a unique stationary state $`\rho _{\mathrm{ss}}`$, the SSE is ergodic over $`E^𝒰`$ and $`\mathrm{}_k^𝒰`$ is equal to the proportion of time the system spends in state $`|\psi _k^𝒰`$. The ensemble $`E^𝒰`$ represents $`\rho _{\mathrm{ss}}`$ in that
$$\underset{k}{}\mathrm{}_k^𝒰|\psi _k^𝒰\psi _k^𝒰|=\rho _{\mathrm{ss}},$$
(8)
as guaranteed by Eq. (6).
### C Survival Probability
Imagine that the system has been evolving under a particular unraveling $`𝒰`$ from an initial state at time $`\mathrm{}`$ to the stationary ensemble at the present time $`0`$. It will then be in the state $`|\psi _k^𝒰`$ with probability $`\mathrm{}_k^𝒰`$. If we now cease to monitor the system then the state will no longer remain pure, but rather will relax toward $`\rho _{\mathrm{ss}}`$ under the evolution of Eq. (1).
This relaxation to equilibrium will occur at different rates for different states. For example, some unravelings will tend to collapse the system into a pure state that is very fragile, in that it changes into a very different (and mixed) state as it relaxes to equilibrium. In this case the ensemble would rapidly become a poor representation of the observer’s expected knowledge about the system. Hence we can say that such an ensemble is a “bad” or “unnatural” representation of $`\rho `$. Conversely, an unraveling that produces robust states would remain an accurate description for a relatively long time. We expect such a “good” or “natural” ensemble to give more intuition about the dynamics of the system. The most robust ensemble we interpret as the “best” or “most natural” such ensemble.
We quantify the robustness of a particular state $`|\psi _k^𝒰`$ by its survival probability $`S_k^𝒰(t)`$. This is the probability that the system would be found (by a hypothetical projective measurement) to still be in the state $`|\psi _k^𝒰`$ at time $`t`$. It is given by
$$S_k^𝒰(t)=\psi _k^𝒰|e^t\left[|\psi _k^𝒰\psi _k^𝒰|\right]|\psi _k^𝒰.$$
(9)
Since we are considering an ensemble $`E^𝒰`$ we must define the average survival probability
$$S^𝒰(t)=\underset{k}{}\mathrm{}_k^𝒰S_k^𝒰(t).$$
(10)
In the limit $`t\mathrm{}`$ the ensemble-averaged survival probability will tend towards the stationary value
$$S^𝒰(\mathrm{})=\mathrm{Tr}[\rho _{\mathrm{ss}}^2].$$
(11)
This is independent of the unraveling $`𝒰`$ and is a measure of the mixedness of $`\rho _{\mathrm{ss}}`$.
There are many possible figures-of-merit that may be obtained from the survival probability $`S^𝒰(t)`$, as discussed in Ref. . Here we choose the simplest one, also adopted in : the time it takes for $`S^𝒰(t)`$ to fall half-way to its equilibrium value. That is,
$$\tau ^𝒰=\mathrm{min}\{t:S^𝒰(t)=(1+\mathrm{Tr}[\rho _{\mathrm{ss}}^2])/2\}.$$
(12)
This survival time $`\tau ^𝒰`$ quantifies the robustness of a particular unraveling $`𝒰`$.
Let the set of all unravelings be denoted $`J`$. Then the subset of maximally robust unravelings $`J_M`$ is
$$J_M=\{J:\tau ^{}\tau ^𝒰𝒰J\}.$$
(13)
Even if $`J_M`$ has many elements $`_1,_2,\mathrm{}`$, these different unravelings may give the same ensemble $`E^{}=E^_1=E^_2=\mathrm{}`$. In this case we claim $`E^{}`$ is the most natural ensemble representation of the stationary solution of a given master equation. Different definitions of survival time will obviously lead to different numerical values for $`\tau ^{}`$. We are less concerned with such numerical values than with the robust ensemble $`E^{}`$, which has been found to depend little on the precise definition used.
## III The Two-Level Atom
### A The Resonance Fluorescence Master Equation
Consider an atom with two relevant levels $`\{|g,|e\}`$. Let there be a dipole moment between these levels so that the coupling to the continuum of electromagnetic field modes in the vacuum state will cause the atom to decay at rate $`\gamma `$. So that the atom does not simply decay to the state $`|g`$, add driving by a classical field (such as that produced by a laser) of Rabi frequency $`\mathrm{\Omega }`$. We work in the interaction picture with respect to the free Hamiltonian $`H_0=\mathrm{}\omega _0|ee|`$ so that the classical driving at frequency $`\omega `$ becomes time-independent. The evolution of the atom’s state matrix can then be described by the resonance fluorescence (RF) master equation
$$\dot{\rho }=i\frac{\mathrm{\Omega }}{2}[\sigma _x,\rho ]+\gamma 𝒟[\sigma ]\rho .$$
(14)
In this equation we have used the Pauli matrices
$`\sigma _x`$ $`=`$ $`|eg|+|ge|`$ (15)
$`\sigma _y`$ $`=`$ $`i|eg|+i|ge|`$ (16)
$`\sigma _z`$ $`=`$ $`|ee||gg|`$ (17)
$`\sigma `$ $`=`$ $`|ge|={\displaystyle \frac{1}{2}}(\sigma _xi\sigma _y)`$ (18)
$`\sigma ^{}`$ $`=`$ $`|eg|={\displaystyle \frac{1}{2}}(\sigma _x+i\sigma _y)`$ (19)
In terms of these, any state of the atom can be written as a 3-vector $`(x,y,z)`$ satisfying
$$x^2+y^2+z^21,$$
(20)
with equality only for a pure state. From this Bloch vector the state matrix is defined by
$$\rho =\frac{1}{2}\left(I+x\sigma _x+y\sigma _y+z\sigma _z\right).$$
(21)
The linear equations of motion for the Bloch vector that result from Eq. (14) are known as the Bloch equations. The solution satisfying the initial condition
$$x(0)=u,y(0)=v,z(0)=w$$
(22)
is
$`x(t)`$ $`=`$ $`ue^{(\gamma /2)t},`$ (23)
$`y(t)`$ $`=`$ $`c_+e^{\lambda _+t}+c_{}e^{\lambda _{}t}+y_{\mathrm{ss}},`$ (24)
$`z(t)`$ $`=`$ $`c_+{\displaystyle \frac{\gamma 4i\stackrel{~}{\mathrm{\Omega }}}{4\mathrm{\Omega }}}e^{\lambda _+t}+c_{}{\displaystyle \frac{\gamma +4i\stackrel{~}{\mathrm{\Omega }}}{4\mathrm{\Omega }}}e^{\lambda _{}t}+z_{\mathrm{ss}}.`$ (25)
Here $`c_\pm `$ are constants given by
$$c_\pm =\frac{1}{8i\stackrel{~}{\mathrm{\Omega }}}\left[4\mathrm{\Omega }(wz_{\mathrm{ss}})\pm (\gamma \pm 4i\stackrel{~}{\mathrm{\Omega }})(vy_{\mathrm{ss}})\right].$$
(26)
The eigenvalues $`\lambda _\pm `$ are defined by
$$\lambda _\pm =\frac{3}{4}\gamma \pm i\stackrel{~}{\mathrm{\Omega }}.$$
(27)
Here
$$\stackrel{~}{\mathrm{\Omega }}=\sqrt{\mathrm{\Omega }^2(\gamma /4)^2}$$
(28)
is a real modified Rabi frequency for $`\mathrm{\Omega }>\gamma /4`$, and is imaginary for $`\mathrm{\Omega }<\gamma /4`$. The stationary solutions appearing in the above equations are
$`x_{\mathrm{ss}}`$ $`=`$ $`0,`$ (29)
$`y_{\mathrm{ss}}`$ $`=`$ $`{\displaystyle \frac{2\gamma \mathrm{\Omega }}{\gamma ^2+2\mathrm{\Omega }^2}},`$ (30)
$`z_{\mathrm{ss}}`$ $`=`$ $`{\displaystyle \frac{\gamma ^2}{\gamma ^2+2\mathrm{\Omega }^2}}.`$ (31)
### B The Survival Probability
Using the Bloch vector representation of the atomic state matrix it is easy to show that the survival probability for a pure state $`|\psi _k`$ with projector
$$|\psi _k\psi _k|=\frac{1}{2}\left(1+u_k\sigma _x+v_k\sigma _y+w_k\sigma _z\right)$$
(32)
is
$$S_k(t)=\frac{1}{2}\left(1+x_k(t)u_k+y_k(t)v_k+z_k(t)w_k\right),$$
(33)
where $`(x,y,z)_k(t)`$ is the Bloch vector at time $`t`$ with the initial condition $`(x,y,z)_k(0)=(u,v,w)_k`$ as in Eqs. (23)–(25). From that solution it is evident that $`S_k(t)`$ will contain terms that are constant, linear, and bilinear in the vector components $`(u,v,w)_k`$ of the initial state.
As explained in the preceding section we are interested in the survival probability not for a single state but for an ensemble of states. This is the ensemble average of Eq. (10). After some work, the survival probability in this case is found to be simply
$`S(t)`$ $`=`$ $`{\displaystyle \underset{k}{}}\mathrm{}_kS_k(t)`$ (34)
$`=`$ $`{\displaystyle \frac{1}{2}}(1+y_{\mathrm{ss}}^2+z_{\mathrm{ss}}^2)+{\displaystyle \frac{1}{2}}`$ (36)
$`\times \left[(1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2)e^{(\gamma /2)t}+V_vf_+(t)+V_wf_{}(t)\right],`$
where
$$f_\pm (t)=e^{(\gamma /2)t}+e^{(3\gamma /4)t}\left(\mathrm{cos}\stackrel{~}{\mathrm{\Omega }}t\pm \frac{\gamma }{4\stackrel{~}{\mathrm{\Omega }}}\mathrm{sin}\stackrel{~}{\mathrm{\Omega }}t\right).$$
(37)
In Eq. (36), all the information about the ensemble is contained in the moments
$`V_v`$ $``$ $`\mathrm{E}[v^2]E[v]^2={\displaystyle \underset{k}{}}\mathrm{}_kv_k^2\left({\displaystyle \underset{k}{}}\mathrm{}_kv_k\right)^2,`$ (38)
$`V_w`$ $``$ $`\mathrm{E}[w^2]E[w]^2={\displaystyle \underset{k}{}}\mathrm{}_kw_k^2\left({\displaystyle \underset{k}{}}\mathrm{}_kw_k\right)^2.`$ (39)
This is possible because we have used the following relations:
$`\mathrm{E}[(u,v,w)]`$ $`=`$ $`(x_{\mathrm{ss}},y_{\mathrm{ss}},z_{\mathrm{ss}}),`$ (40)
$`\mathrm{E}[u^2]`$ $`=`$ $`1\mathrm{E}[v^2]\mathrm{E}[w^2].`$ (41)
To find the robustness of any particular unraveling we thus need to find simply the two ensemble averages $`V_v`$ and $`V_w`$.
## IV Unraveling by Direct Detection
The most obvious way to unravel the RF master equation is by direct detection. This requires detecting all of the atom’s fluorescence by unit-efficiency photodetectors. This is beyond current technology, but not by so much that the experiment should be considered unphysical. As we will find, unraveling by direct detection is actually not very robust (by the definition of Sec. II) but it is nevertheless useful to consider as a point of comparison with more robust unravelings.
The stochastic evolution of an atom undergoing RF with direct detection has been considered many times before . It has one feature that enables an enormous simplification over a generic unraveling. This is that immediately following a detection, the atomic state is independent of its state before, and is just the ground state $`|g`$. Between these jumps to the ground state, the conditioned atomic state evolves deterministically. At steady state, when there has certainly been at least one detection, all members of the ensemble are therefore identified simply by the time $`t`$ since the last detection.
If there is a detection at time $`t_0`$, then, until the next detection occurs, the state of the atom at time $`t_0+t`$ evolves according to the equation
$$\frac{d}{dt}|\stackrel{~}{\psi }_0(t)=\left(\frac{\gamma }{2}\sigma ^{}\sigma i\frac{\mathrm{\Omega }}{2}\sigma _x\right)|\stackrel{~}{\psi }_0(t).$$
(42)
The solution, satisfying the initial condition $`|\stackrel{~}{\psi }_0(0)=|g`$ is
$$|\stackrel{~}{\psi }_0=\stackrel{~}{c}_e(t)|e+\stackrel{~}{c}_g(t)|g,$$
(43)
where
$`\stackrel{~}{c}_g(t)`$ $`=`$ $`\left[\mathrm{cos}(\stackrel{ˇ}{\mathrm{\Omega }}t/2)+{\displaystyle \frac{\gamma }{2\stackrel{ˇ}{\mathrm{\Omega }}}}\mathrm{sin}(\stackrel{ˇ}{\mathrm{\Omega }}t/2)\right]e^{(\gamma /4)t},`$ (44)
$`\stackrel{~}{c}_e(t)`$ $`=`$ $`i{\displaystyle \frac{\mathrm{\Omega }}{\stackrel{ˇ}{\mathrm{\Omega }}}}\mathrm{sin}(\stackrel{ˇ}{\mathrm{\Omega }}t/2)e^{(\gamma /4)t}.`$ (45)
Here
$$\stackrel{ˇ}{\mathrm{\Omega }}=\sqrt{\mathrm{\Omega }^2(\gamma /2)^2}$$
(46)
is a real modified Rabi frequency for $`\mathrm{\Omega }>\gamma /2`$ and is imaginary for $`\mathrm{\Omega }<\gamma /2`$. Note that it is different from $`\stackrel{~}{\mathrm{\Omega }}`$ defined in Eq. (28).
The state in Eq. (43) is unnormalized, and the norm $`\stackrel{~}{\psi }_0(t)|\stackrel{~}{\psi }_0(t)`$ represents the probability that there has been no detection since time $`t_0`$, given that there was a detection at that time. Let us write this probability as
$$P_0(t)=|\stackrel{~}{c}_e(t)|^2+|\stackrel{~}{c}_g(t)|^2.$$
(47)
We show in the appendix that this probability is related to $`\mathrm{}(t)`$, the probability that, at steady state, the last detection was a time $`t`$ ago, by
$$\mathrm{}(t)=\frac{P_0(t)}{_0^{\mathrm{}}P_0(s)𝑑s}.$$
(48)
As noted above, in steady state under direct detection the possible atomic states are parametrized by the real variable $`t`$, the time since the last detection. The state at that time has projector
$$\widehat{P}(t)=\frac{|\stackrel{~}{\psi }_0(t)\stackrel{~}{\psi }_0(t)|}{\stackrel{~}{\psi }_0(t)|\stackrel{~}{\psi }_0(t)},$$
(49)
and the weight for each of these members of the ensemble is $`\mathrm{}(t)dt`$. Physically, all members of the ensemble exist on the $`u=0`$ great circle of the Bloch sphere, because that is where the modified Rabi cycling of Eq. (42) takes the ground state. This distribution is shown in Fig. 1(a). For $`\stackrel{ˇ}{\mathrm{\Omega }}`$ imaginary (that is, $`\mathrm{\Omega }<\gamma /2`$) the states never reach the excited state. For $`\stackrel{ˇ}{\mathrm{\Omega }}`$ real (that is, $`\mathrm{\Omega }>\gamma /2`$, the states may undergo an arbitrary number of cycles. For $`\mathrm{\Omega }\gamma `$ the states are likely to undergo many cycles before a spontaneous emission event occurs so that the ensemble consists of all the states on the $`u=0`$ great circle, almost uniformly distributed.
From Eq. (49) it can be verified analytically that
$$_0^{\mathrm{}}\widehat{P}(t)\mathrm{}(t)𝑑t=\frac{1}{𝒩}_0^{\mathrm{}}𝑑t|\stackrel{~}{\psi }_0(t)\stackrel{~}{\psi }_0(t)|=\rho _{\mathrm{ss}},$$
(50)
where
$$𝒩=_0^{\mathrm{}}\left[|\stackrel{~}{c}_e(t)|^2+|\stackrel{~}{c}_g(t)|^2\right]𝑑t.$$
(51)
. Moreover we can easily find numerically the ensemble averages necessary to find the ensemble average survival probability, namely
$`V_v`$ $`=`$ $`{\displaystyle \frac{1}{𝒩}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{[i\stackrel{~}{c}_e(t)\stackrel{~}{c}_g^{}(t)i\stackrel{~}{c}_e^{}(t)\stackrel{~}{c}_g(t)]^2}{|\stackrel{~}{c}_e(t)|^2+|\stackrel{~}{c}_g(t)|^2}}𝑑ty_{\mathrm{ss}}^2,`$ (52)
$`V_w`$ $`=`$ $`{\displaystyle \frac{1}{𝒩}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{[|\stackrel{~}{c}_e(t)|^2|\stackrel{~}{c}_g(t)|^2]^2}{|\stackrel{~}{c}_e(t)|^2+|\stackrel{~}{c}_g(t)|^2}}𝑑tz_{\mathrm{ss}}^2,`$ (53)
In Fig. 2 we plot the survival probability for the direct detection ensemble, for a variety of driving strengths $`\mathrm{\Omega }`$. We see that for $`\mathrm{\Omega }<\gamma `$ the survival probability decays approximately exponentially at rate of order $`\gamma `$. For $`\mathrm{\Omega }\gamma `$ the stationary state matrix is close to the ground state, and most members of the direct detection ensemble are also. For $`\mathrm{\Omega }\gamma `$ the ensemble is equally spread over the $`u=0`$ great circle and consequently the variances $`V_v`$ and $`V_w`$ are approximately equal to $`1/2`$. Using this, the survival probability is found to be approximately
$$S(t)\frac{1}{2}\left(1+e^{(3/4)\gamma t}\mathrm{cos}\stackrel{~}{\mathrm{\Omega }}t\right).$$
(54)
The oscillations in the survival probability are due to the Rabi oscillations. As noted above, the direct detection ensemble consists of states on the $`u=0`$ great circle. Rabi cycling around the $`x`$-axis according to the RF master equation (14) rotates this circle around, rapidly moving the states away from their initial positions and then back close to their initial conditions after one cycle. They do not return exactly to their initial states because of the slow (at rate $`3\gamma /4`$) decay towards the equilibrium state. This behaviour is illustrated for a typical member of the direct detection ensemble in Fig. 1(b).
The change from damped to oscillatory behaviour has a dramatic effect on the survival time in Eq. (12). It is plotted in Fig. 3 as a function of $`\mathrm{\Omega }`$. For $`\mathrm{\Omega }\gamma `$ it is given by
$$\tau 2\mathrm{ln}2\gamma ^1,$$
(55)
as in this limit the survival probability decays as $`e^{\gamma t/2}`$. For $`\mathrm{\Omega }\gamma `$, we can use Eq. (54) to get the approximate expression
$$\tau \frac{\pi }{3}\mathrm{\Omega }^1.$$
(56)
That is, the survival time is here determined by the Hamiltonian evolution only. As shown in Fig. 3, this is quite a good approximation even for moderate $`\mathrm{\Omega }`$.
## V The Most Robust Unraveling
### A The Most Robust Ensemble
It was stated above that the unraveling by direct detection is not the most robust unraveling. In fact, from an examination of the survival probability in Eq. (36) we can see that it is one of the least robust unravelings for $`\mathrm{\Omega }\gamma `$. That is because of the large variances in $`v`$ and $`w`$ in this limit.
It is not difficult to show that the two functions $`f_+(t)`$ and $`f_{}(t)`$, defined in Eq. (37), are non-positive for all $`\mathrm{\Omega }`$ and for $`t>0`$. Since the variances $`V_v`$ and $`V_w`$ are non-negative, it is easy to see that to maximize the survival probability, one would wish to minimize $`V_v`$ and $`V_w`$. The ideal limit would be $`V_v=V_w=0`$. This corresponds to an ensemble in which all members have the same Bloch vector components $`v`$ and $`w`$. Since the ensemble average must equal $`\rho _{\mathrm{ss}}`$ it follows then that for all members
$$v=y_{\mathrm{ss}},w=z_{\mathrm{ss}}.$$
(57)
Furthermore, since $`x_{\mathrm{ss}}=0`$, and since the members of the ensemble must be pure, it follows that for all members
$$u=\pm \sqrt{1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2},$$
(58)
where the two alternatives are equally weighted. These two members of the ensemble are shown in Fig. 1(c).
This ensemble is guaranteed to give the maximum survival probability
$$S(t)=\frac{1}{2}(1+y_{\mathrm{ss}}^2+z_{\mathrm{ss}}^2)+\frac{1}{2}\left[(1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2)e^{(\gamma /2)t}\right],$$
(59)
which is plotted in Fig. 5. It gives the maximum survival time
$$\tau =2\mathrm{ln}2\gamma ^1,$$
(60)
which is independent of $`\mathrm{\Omega }`$. There is no Rabi cycling because the states defined by the Bloch vector $`(u,v,w)`$ have $`v`$ and $`w`$ already equal to their stationary values. Under the master equation evolution $`u`$ simply decays towards its stationary value of zero and $`v`$ and $`w`$ remain constant. This decay of the Bloch vector to equilibrium at rate $`\gamma /2`$ is shown in Fig. 1(d).
### B Adaptive Interferometric Detection
The most robust ensemble defined here would be a mere curiosity if it were not for the fact that there is a detection scheme that realizes it. This scheme, proposed by one of us and Toombes , involves interfering the light from the atom with a resonant local oscillator before detection. This is done using a highly transmitting beam splitter as shown in Fig. 4. In the limit of a large local oscillator, this is known as homodyne detection. However we require a very weak local oscillator, with reflected intensity comparable to the intensity of the light from the atom. Furthermore, we require the local oscillator amplitude to be continually adjusted by a real-time feedback loop. To be specific, the field detected should be proportional to
$$\sigma +\mu (t),$$
(61)
where the field from the atom is proportional to $`\sigma `$, the atomic lowering operator as usual, and the local oscillator field is represented by the complex number $`\mu (t)`$. This complex number is given by
$$\mu (t)=\pm \frac{1}{2},$$
(62)
where the sign is changed every time a detection occurs.
Remarkably, this relatively simple detection scheme has the consequence that, after initial transients have died away, a driven atom jumps between the states with projector $`\widehat{P}_\pm `$, where
$$\widehat{P}_\pm =\frac{1}{2}\left(I\pm \sqrt{1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2}\sigma _x+v_{\mathrm{ss}}\sigma _y+w_{\mathrm{ss}}\sigma _z\right),$$
(63)
every time a detection occurs. The rate of detections moreover is independent of which of these states the atom is in, and is equal to $`\gamma /4`$. Thus in the long time limit the atom will have a probability $`\mathrm{}_\pm =1/2`$ to be in each of them, and the maximally robust ensemble will be physically realized.
## VI Continuous Markovian Unravelings
The most robust unraveling is neither continuous nor Markovian. It is not continuous because the atomic state jumps every time there is a detection. It is not Markovian because the evolution of the atomic state does not depend only on its present state. Rather, it depends on the past history of detections through the local oscillator amplitude $`\mu (t)`$. As noted in the introduction, previous investigations of robust unravelings in other systems have been restricted to unravelings that do have these properties. It is therefore interesting to ask for the present system, how close to the MRU is the most robust unraveling that is continuous and Markovian?
Continuous Markovian unraveling (CMU) of the RF master equation (14) can be represented by a two-parameter family of nonlinear SSEs for the non-normalized state vector $`|\overline{\psi }(t)`$ of the form
$$d|\overline{\psi }(t)=dt\left[iH\frac{\gamma }{2}\sigma ^{}\sigma +J(t)\sigma \right]|\overline{\psi }(t),$$
(64)
which is to be interpreted in the Itô sense . Here $`J(t)`$ is a complex “current” given by
$$J(t)dt=\gamma \upsilon \sigma +\sigma ^{}dt+\sqrt{\gamma }dW(t),$$
(65)
where $`\upsilon `$ is a complex number satisfying
$$\upsilon ^{}\upsilon 1,$$
(66)
and the angle brackets denote a quantum average using the normalized state vector $`|\psi (t)`$. We use $`|\overline{\psi }`$ rather than $`|\stackrel{~}{\psi }`$ because the norm of $`|\overline{\psi }`$ has no interpretation in terms of probability, unlike that of $`|\stackrel{~}{\psi }`$ in Sec. IV. The stochastic term $`dW(t)`$ is a complex Gaussian white noise term satisfying
$`\mathrm{E}[dW]`$ $`=`$ $`0,`$ (67)
$`\mathrm{E}[dW^{}dW]`$ $`=`$ $`dt,`$ (68)
$`\mathrm{E}[(dW)^2]`$ $`=`$ $`\upsilon dt.`$ (69)
The complex parameter $`\upsilon `$ comprises the two parameters for the family of unravelings of the form of Eq. (64). From Eq. (67) and Eq. (68) it can be shown that Eq. (6) is satisfied for all $`\upsilon `$. Thus
$$\rho (t)=\mathrm{E}\left[\frac{|\overline{\psi }(t)\overline{\psi }(t)|}{\overline{\psi }(t)|\overline{\psi }(t)}\right]$$
(70)
obeys the RF master equation (14), while in an individual trajectory the state remains pure. In this case there is no simple way to find the steady state ensemble of pure states. A numerical simulation of Eq. (64), with time averages replacing ensemble averages, is the only way to proceed.
Finding the maximally robust continuous Markovian unraveling (MRCMU) in this case reduces to a search over the ball $`|\upsilon |^21`$ in the complex plane. We find that the MRCMU is for $`\upsilon =1`$. For this value of $`\upsilon `$ the “current” $`J(t)`$ (which is real in this case) has a deterministic part equal to
$$\mathrm{E}\left[J(t)\right]=\gamma \sigma _x.$$
(71)
That is, the measurement yields information about the $`\sigma _x`$ quadrature of the atomic dipole. As a consequence it tends to localize the atom near the $`\sigma _x`$ eigenstates. This localization is relatively stable for high driving, since these are also eigenstates of the Hamiltonian $`\mathrm{\Omega }\sigma _x/2`$. This is shown in Fig. 1(e), which is a stochastically generated sample of 10000 states from the equilibrium ensemble for the MRCMU. In complete contrast to the direct detection ensemble in Fig. 1(a), most of the states lie close to the $`\sigma _x`$ eigenstates. For $`\mathrm{\Omega }`$ large as in this figure, this means close to the two members of the MRU shown in Fig. 1(c).
Because a typical member of the MRCMU ensemble is fairly close to the $`\sigma _x`$ axis, it is relatively little affected by the Hamiltonian, which causes rotation around that axis. Under the RF master equation (14) its evolution consists of decay to the equilibrium, with relatively small Rabi oscillations superimposed. This is as shown in Fig. 1(f). As a consequence, although the survival probability oscillates, it remains above its equilibrium value and decays towards it at a rate proportional to $`\gamma `$. This is shown in Fig. 5 for $`\mathrm{\Omega }=10\gamma `$. Also shown, for comparison, is the survival probability for the minimally robust CMU, which occurs for $`\upsilon =1`$. This is almost identical to the corresponding curve for direct detection in Fig. 2, as the ensemble is confined to the $`u=0`$ great circle in both cases.
The survival time for the MRCMU is shown in Fig. 3. Because of our definition of survival time in Eq. (12), the survival time for the MRCMU is determined by the rapid (for $`\mathrm{\Omega }\gamma `$) oscillations in the survival probability rather than the slow mean decay. Thus it is qualitatively similar to the survival time for the direct detection ensemble, starting at $`2\mathrm{ln}2\gamma ^1`$ and then falling like $`\mathrm{\Omega }^1`$ as $`\mathrm{\Omega }`$ increases.
## VII Comparison of the MRU and the MRCMU
We have seen that no continuous Markovian unraveling (CMU) is as robust as the maximally robust unraveling (MRU), which is neither continuous nor Markovian. Furthermore, the robustness for the maximally robust CMU, as measured by the survival time, scales in the same way as that for the very non-robust direct detection scheme. However, because of the arbitrariness in any definition of survival time, we are interested more in the robust ensembles themselves than in the numerical values of their survival times. As discussed above, the MRCMU realizes an ensemble that has a common feature with that realized by the MRU: the states in the ensemble tend to have well-defined values of $`u=\sigma _x`$. In this section we wish to answer the question: just how close is the MRCMU ensemble to the MRU ensemble?
### A Closeness of Two Ensembles
To answer this question we require a measure (not necessarily transitive) from ensemble $`E^A`$ to ensemble $`E^B`$, where each ensemble represents the same state matrix. It seems best to choose an operationally defined measure, which we do as follows.
By allowing the same projector to reappear with different indices, we can write, to any desired degree of accuracy,
$`E^A`$ $`=`$ $`\left\{(|\varphi _k,N^1):1kN\right\},`$ (72)
$`E^B`$ $`=`$ $`\left\{(|\psi _\mu ,N^1):1\mu N\right\},`$ (73)
where $`|\varphi _k`$ and $`|\psi _\mu `$ are normalized states such that
$$\rho =N^1\underset{k}{}|\varphi _k\varphi _k|=N^1\underset{\mu }{}|\psi _\mu \psi _\mu |,$$
(74)
where $`N`$ is an arbitrarily large integer.
To define the closeness of the ensembles, imagine that there are two people, Alice and Bob. Alice has in her possession a measuring device with $`N`$ settings, corresponding to the $`N`$ projectors $`\{|\varphi _k\varphi _k|:k\}`$. If setting $`k`$ is chosen then the device makes a projective measurement with projector $`|\varphi _k\varphi _k|`$. Bob has in his possession an ensemble of quantum states $`\{|\psi _\mu :\mu \}`$. It is Bob’s aim to try to convince Alice that he is actually in possession of the ensemble $`\{|\varphi _k:k\}`$. He must submit each of his $`N`$ systems $`|\psi _\mu `$ to Alice, telling her which of her states $`|\varphi _k`$ each is supposed to be in. She then makes the appropriate measurement and, unless Bob’s ensemble really is the same as Alice’s, is likely to find errors some of the time. An error is when a state that Bob claims is $`|\psi _k`$ is found to give the answer “no” to Alice’s projective measurement “is the state $`|\varphi _k`$?” Assuming Bob chooses a good strategy, then the higher the probability of error, the larger the distance between the two ensembles.
We can formalize this as follows. Say Bob actually sends state $`|\psi _\mu `$, but claims it is $`|\varphi _{k(\mu )}`$, where the functional dependence here indicates that Bob makes his choice of index $`k`$ based on his actual state. Then the probability of error for this state is
$$ϵ_{k(\mu )|\mu }=1\left|\varphi _{k(\mu )}|\psi _\mu \right|^2.$$
(75)
The ensemble average probability of error, for Bob’s optimum strategy, is
$$ϵ_{\mathrm{opt}}=N^1\underset{M}{\mathrm{min}}\underset{\mu =1}{\overset{N}{}}ϵ_{k(\mu )|\mu },$$
(76)
where the minimum is over all one-to-one mappings $`M`$
$$\mu \stackrel{M}{}k(\mu ),k\stackrel{M^1}{}\mu (k).$$
(77)
Bob’s strategies have to correspond to a mapping of this form because unless Bob names each of Alice’s states once and once only Alice would know that Bob is lying when he claims to be in possession of Alice’s ensemble $`\{|\varphi _k:k\}`$.
We could take the distance between the ensembles to be equal to this minimum average error probability. However, if the state matrix $`\rho `$ is close to being pure (with $`\mathrm{Tr}[\rho ^2]`$ close to one), then the average error probability would be small regardless of what strategy Bob chose. In particular, if Bob had a totally random strategy then the error probability for an individual state $`|\psi _\mu `$ would be, on average,
$$ϵ_{k|\mu }=1\psi _\mu |\overline{|\varphi _k\varphi _k|}|\psi _\mu =1\psi _\mu |\rho |\psi _\mu ,$$
(78)
and the ensemble average error probability would be
$$\overline{ϵ}=N^1\underset{\mu =1}{\overset{N}{}}\left(1\mathrm{Tr}[\rho |\psi _\mu \psi _\mu |]\right)=1\mathrm{Tr}[\rho ^2].$$
(79)
That is, the average error probability would be close to zero even though Bob does not take into account the difference between his states so that his effective ensemble consists of $`N`$ copies of the mixed state $`\rho `$.
For this reason, it seems better to define the distance between the ensembles by the normalized error probability:
$$d(E^A|E^B)=\frac{1\mathrm{min}_M_{\mu =1}^N\frac{1}{N}\left|\varphi _{k(\mu )}|\psi _\mu \right|^2}{1\mathrm{Tr}[\rho ^2]}.$$
(80)
In this case it is easy (at least for a two-level system) to see that
$$0d(E^B|E^A)1,$$
(81)
where the lower bound is attained if and only if the ensembles are identical, and where no tighter upper bound can be found for a given $`\rho `$. Thus the two ensembles could be said to be close if and only if $`d(E^B|E^A)1`$.
### B Closeness of the MRU and MRCMU Ensembles
In the case at hand the reference ensemble $`E^A`$ is the most robust ensemble of Sec. V, while the ensemble $`E^B`$ whose closeness we wish to gauge is that of the most robust continuous Markovian unraveling of Sec. VI. The fact that $`E^A`$ has only two elements, whereas $`E^B`$ has an infinitude of elements causes no problems. We simply allow an arbitrarily large number $`N`$ of elements for each ensemble and let half of those for ensemble $`A`$ be the state with projector $`\widehat{P}_+`$ and half the state $`\widehat{P}_{}`$, as defined in Eq. (63).
Because the two states $`\widehat{P}_\pm `$ differ only by the sign of $`\sigma _x`$, and because both ensembles are symmetric about a reflection in the $`yz`$ plane, Bob’s best strategy is easy to find. For each of his states $`|\psi _\mu `$ he tells Alice that it is state $`\widehat{P}_+`$ if it has a positive mean $`\sigma _x`$, and $`\widehat{P}_{}`$ if it has a negative mean $`\sigma _x`$. If $`|\psi _\mu `$ is the state
$$|\psi _\mu \psi _\mu |=\frac{1}{2}(1+u_\mu \sigma _x+v_\mu \sigma _y+w_\mu \sigma _z),$$
(82)
then the probability of error for this state is
$`ϵ_\mu `$ $`=`$ $`1\psi _\mu |\widehat{P}_{\mathrm{sign}(u_\mu )}|\psi _\mu `$ (83)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}(|u_\mu |\sqrt{1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2}+v_\mu y_{\mathrm{ss}}+w_\mu z_{\mathrm{ss}}).`$ (84)
Averaging over all of Bob’s states we get
$$ϵ_{\mathrm{opt}}=\frac{1}{2}\frac{1}{2}\left(\frac{1}{N}\underset{\mu }{}|u_\mu |\sqrt{1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2}+y_{\mathrm{ss}}^2+z_{\mathrm{ss}}^2\right),$$
(85)
and the distance from the MRU ensemble to the MRCMU ensemble is
$$d(E^A|E^B)=1\frac{\mathrm{E}[|u|]}{\sqrt{1y_{\mathrm{ss}}^2z_{\mathrm{ss}}^2}}.$$
(86)
Thus to find how close the MRCMU ensemble is to the MRU ensemble we simply need to evaluate $`\mathrm{E}[|u|]`$, the ensemble average of $`|\sigma _x|`$ for the former. The result is plotted in Fig. 6. The distance is always less than about $`0.3`$, the value to which it appears to asymptote for large $`\mathrm{\Omega }`$. Since this is moderately small compared to one, we can say that the two ensembles are moderately close. This is in stark contrast to either the direct detection ensemble or the $`\upsilon =1`$ CMU ensemble. For these two ensembles, $`u=0`$ for all members, so the distance to the MRU ensemble is the maximum value of unity.
Also plotted in Fig. 6 is the distance from the CMU with $`\upsilon =0`$ to the MRU ensemble. This CMU has a number of special properties and is sometimes known as quantum state diffusion . We see that it also gives an ensemble whose distance to the MRU ensemble is less than unity. However, the distance is considerably greater than that of the MRCMU, asymptoting to a value greater than $`0.4`$.
## VIII Discussion
The resonance fluorescence master equation generates surprisingly rich dynamics for a two-level atom. Here we have investigated those dynamics using the technique of finding the maximally robust unraveling. That is, finding the scheme for monitoring the fluorescent radiation that collapses the atom into pure states that are, on average, the most robust. By robust we mean that they survive best under the master equation evolution so that, once the monitoring ceases, the probability for the atom to be found at some later time to still be in the state into which it was collapsed by the monitoring, is maximized.
The property of producing robust states may give the maximally robust unraveling potential applications, particularly in quantum information technology where minimizing the effect of environmental decoherence is essential. Quite separately from any application, the maximally robust unraveling is useful for characterizing the dynamics of open quantum systems, such as the resonantly driven atom of this study. It has been suggested before that the ensemble arising from the MRU is the most natural representation of the system’s stationary state matrix in terms of pure states (state vectors).
In this work we have found that the MRU for the RF master equation is an adaptive interferometric monitoring scheme proposed in Ref. . The atom’s radiation is, prior to detection by a photodetector, interfered at a beam splitter with a reflected local oscillator. The measurement is adaptive because the local oscillator amplitude (which is comparable in magnitude to the field radiated by the atom) has its phase changed by $`\pi `$ every time a detection occurs. This detection scheme has the remarkable property that, in steady state, the atom simply jumps between two fixed pure states. In the large driving limit, these two pure states are close to eigenstates of the driving Hamiltonian $`\mathrm{\Omega }\sigma _x/2`$.
The adaptive interferometric monitoring scheme was designed in Ref. specifically to have this property of producing a stationary ensemble containing just two members. In that reference it was found that other detection schemes, such as spectral detection (resolving the three Mollow peaks), and another adaptive scheme, give rise to similar behaviour. That is, the atom jumped between states that were close to $`\sigma _x`$ eigenstates in the large driving limit. In Ref. it was speculated that this behaviour was what the atom “wanted to do”. Here we have confirmed that the resulting two-member ensemble is indeed the most robust and hence arguably the most natural. In the context of the study of decoherence and the classical limit, it appears that jumping between two fixed states is the most classical behaviour for a strongly-driven two level atom.
Another issue we have investigated in this paper is how close to the MRU one can approach if one restricts the unravelings to continuous Markovian ones. In the context of the fluorescent atom, this means unravelings realizable from homodyne measurements. They give rise to evolution on the Bloch sphere that is continuous (but not differentiable) and Markovian. This is an interesting question because the set of continuous Markovian unravelings is easily parameterized by real numbers, unlike the set of all possible unravelings, which is too large to be finitely parameterized in this way. For this reason, previous work in MRU has concentrated on finding the most robust CMU.
In this work we have found that the MRCMU has a robustness, as measured by the survival time, which falls as $`\mathrm{\Omega }^1`$ as the driving $`\mathrm{\Omega }`$ increases. This is similar to the result for direct detection, and contrary to that of the MRU for which the survival time is constant at $`2\mathrm{ln}2\gamma ^1`$. However, as we have shown graphically, the distribution of states on the Bloch sphere for the MRCMU ensemble is qualitatively much closer to that of the MRU than to that of direct detection. Furthermore, we have introduced a quantitative measure for the closeness of two ensembles of pure states, and applied this to the various unravelings. We find that the MRCMU ensemble is reasonably close to the MRU ensemble, while the direct detection ensemble is as distant as is possible from the most robust ensemble.
This result has wider implications in the program of decoherence, robustness, and the classical limit. A two-level atom is an extremely non-classical system. In the limit of strong driving the stationary state matrix is almost fully mixed, and the existence of two discrete levels manifests strongly in the dynamics: the maximally robust unraveling has jumps between two almost orthogonal states. By contrast, under a continuous Markovian unraveling the atomic state does not jump, but rather diffuses around the Bloch sphere.
We have shown that, despite the atom’s nonclassicality, the ensemble generated by a CMU can be reasonably close to that generated by the MRU. This suggests that restricting an investigation to continuous Markovian unravelings is not a serious restriction (provided that one is not interested in the numerical value of the maximum survival time). The basic nature of the maximally robust unraveling should be evident from that of the maximally robust CMU. For the two-level atom it is fortunate that the absolute maximally robust unraveling can be found analytically. For more general systems a numerical search would be necessary, and, to be practical, would have to be confined to finitely parameterizable unravelings such as the continuous Markovian unravelings. Thus our results lends credence to the whole program of finding the (approximately) maximally robust unraveling for open quantum systems.
## Derivation of Eq. (4.7)
Let $`n(t)`$ be the event that there were no detections from time $`t_0`$ to the present time $`t_0+t`$. Let $`d(t)`$ be the event that there was a detection at time $`t`$ before the present. Then
$$P_0(t)=P[n(t)|d(t)],$$
(87)
where $`P_0(t)`$ is as in Eq. (47) and where $`P[A|B]`$ means the probability of $`A`$ given $`B`$.
Now what we want is $`\mathrm{}(t)`$, the probability that the last detection was at a time $`t`$ before the present, at steady state. That is, the probability that there was a detection at time $`t`$ in the past, and that there were no detections from then until now. In other words,
$$\mathrm{}(t)=P[n(t)d(t)],$$
(88)
where $`P[AB]`$ means the joint probability of $`A`$ and $`B`$.
Now from the definition of conditional probability, $`P[AB]=P[A|B]P[B].`$ Therefore
$$\mathrm{}(t)=P_0(t)P[d(t)].$$
(89)
But at steady state (that is, after initial transients have decayed), the probability that there was a detection at time $`t`$ in the past, given no other information, does not depend on $`t`$. That is, $`P[d(t)]`$ is a constant, so we simply have
$$\mathrm{}(t)=𝒩^1P_0(t),$$
(90)
for some constant $`𝒩`$.
To find the constant of proportionality we just note that, since the last detection must have been at some time $`t`$ in the past, $`_0^{\mathrm{}}\mathrm{}(t)𝑑t=1.`$ Here we are actually treating $`\mathrm{}(t)dt`$ as the probability that the last detection was in the interval $`[t,t+dt)`$. From this condition it is easy to see that
$$𝒩=_0^{\mathrm{}}P_0(t)𝑑t,$$
(91)
which gives Eq. (48).
###### Acknowledgements.
This work was partly supported by the Australian Research Council. HMW would like to acknowledge ongoing discussions with John Vaccaro. |
warning/0002/hep-th0002164.html | ar5iv | text | # References
hep-th/0002164
A Comment on Self-Tuning and Vanishing Cosmological Constant in the Brane World
Stefan Förste<sup>1</sup>, Zygmunt Lalak<sup>1,2</sup>, Stéphane Lavignac<sup>1,3</sup> and Hans Peter Nilles<sup>1</sup>
<sup>1</sup>Physikalisches Institut, Universität Bonn
Nussallee 12, D-53115 Bonn, Germany
<sup>2</sup>Institute of Theoretical Physics
Warsaw University, Poland
<sup>3</sup>Service de Physique Théorique, CEA-Saclay
F-91191 Gif-sur-Yvette Cédex, France
Abstract
In this note we elaborate on various five dimensional contributions to the effective 4D cosmological constant in brane systems. In solutions with vanishing 5D cosmological constant we describe a non-local mechanism of cancellation of vacuum energy between the brane and the singularities. We comment on a hidden fine tuning which is implied by this observation.
One of the main issues in the ‘Brane World’ scheme of the superunification of gauge forces with gravity is the question of having the observable sector (containing the Standard Model degrees of freedom) confined to a lower dimensional brane embedded into higher-dimensional bulk without violating experimental bounds on corrections to the precisely measured parameters of the Standard Model and Newton’s law of gravity. In this context the prominent observable which needs to be taken care of is the four dimensional cosmological constant, known to be close to zero at present. In four dimensions the vanishing cosmological constant, or in fact vanishing vacuum expectation value of the energy-momentum tensor, is necessary to obtain Minkowski space as the solution of the 4d Einstein equations. This leads one to consider vacuum configurations in higher dimensions which have four dimensional part of the metric in the form of the Minkowski metric multiplied by the warp factor which depends on transverse coordinates only. The simplest setup corresponds to just a single transverse coordinate, with the line element
$$ds^2=e^{2A(x^5)}\eta _{\mu \nu }dx^\mu dx^\nu +\left(dx^5\right)^2.$$
(1)
The sources for such configurations are assumed to be located on a number of four dimensional branes, one of which represents the observable gauge sector. It is important to note that although the observable gauge interactions are strictly confined to the 3-brane, the gravity and moduli fields permeate the whole space, effectively connecting the walls in a nontrivial way. This implies that every particle localized on the wall feels sources of gravitational forces which are located all over the bulk. In some cases the influence of the remote sources will be suppressed, like in the case of the exponentially falling off graviton wave function in ref. , which would effectively restrict the relevant gravitational sources to the thin layer around the brane, but sometimes the suppression would be so mild, that the influence of the whole bulk contribution will be highly relevant. Thus in any case, even though the gauge forces are restricted to the branes, the gravitational sector has to be completely integrated out when going to the effective four dimensional theory. In particular, this implies that when one computes the effective four dimensional energy density, or four dimensional vacuum pressure, one has to integrate over the whole causally accessible portion of the transverse space. The issue of exact solutions to Einstein equations with Minkowski-space 4d foliations and vanishing 4d vacuum energy and pressure can be studied in detail in the quasi-stringy setup of five dimensional dilaton gravity. Models of this type, with and without singularities located in transverse space, have been recently studied in a series of papers (and references therein). However, the problem of precisely how the cancellation of various contributions to the vacuum energy and pressure occurs in these models has not been explicitly addressed. The details of this mechanism are also important in view of the question about the amount of fine-tuning between contributions coming from spatially disconnected branes required to achieve vanishing 4d parameters. We believe that this sheds more light on the recently advertised miraculous self-tuning mechanism of ,. In fact, the way we see it the self-tuning is a fine-tuning in disguise.
In this note we shall discuss the details of the nonlocal cancellations which occur in the five dimensional dilaton gravity models with various number of branes, with and without bulk cosmological constant. The 5d action we consider is
$$S_5=d^5x\sqrt{g}(R\frac{4}{3}(\varphi )^2\mathrm{\Lambda }e^{a\varphi })+d^5x\sqrt{g_4}(f_i(\varphi )\delta \left(x^5x_i^5\right))$$
(2)
where $`f_i=V_ie^{b_i\varphi }`$ and the index $`i`$ counts the branes. The corresponding Einstein equations are
$$\sqrt{g}(G_{MN}\frac{4}{3}_M\varphi _N\varphi +\frac{2}{3}(\varphi )^2g_{MN}+\frac{1}{2}\mathrm{\Lambda }e^{a\varphi }g_{MN})+\frac{1}{2}\sqrt{g_4}f_i\delta \left(x^5x_i^5\right)g_{\mu \nu }\delta _M^\mu \delta _N^\nu =0$$
(3)
and the dilaton equation of motion is
$$\mathrm{\Lambda }e^{a\varphi }a\sqrt{g}\sqrt{g_4}\delta \left(x^5x_i^5\right)_{x^5}f_i+\frac{8}{3}_M(\sqrt{g}g^{MN}_N\varphi )=0,$$
(4)
where $`M,N=1,\mathrm{},5`$ and $`\mu ,\nu =1,\mathrm{},4`$. The ansatz (1) implies that the 4d cosmological constant should vanish. In the known solutions with non-vanishing 5d constant $`\mathrm{\Lambda }`$ the vacuum energy has contributions from the bulk integration and at the sources. Performing the bulk integral one obtains boundary terms which locally cancel the vacuum energy at the branes. However, for vanishing $`\mathrm{\Lambda }`$ the situation is different as we will explain in what follows.
Our claim ist that the effective four dimensional cosmological constant vanishes only when one conjectures contributions to the vacuum energy being located at the singularities. To show in detail how the vacuum energy located at the singularities cancels the contribution from the brane<sup>1</sup><sup>1</sup>1It was independently observed in that taking only the contribution from the brane at the origin gives a nonvanishing result. we need to resolve the singularity. The easiest way to do this is to put additional sources there such that Einstein’s equation is satisfied everywhere. Therefore we supplement the action of by additional source terms, and (2) takes the form
$`S`$ $`=`$ $`{\displaystyle d^5x\sqrt{G}\left[R\frac{4}{3}\left(\varphi \right)^2\right]}`$ (5)
$`{\displaystyle d^4x\sqrt{g}Ve_{}^{b\varphi }{}_{|x^5=0}{}^{}}`$
$`{\displaystyle d^4x\sqrt{g}V_+e_{}^{b_+\varphi }{}_{|x^5=x_+}{}^{}}`$
$`{\displaystyle d^4x\sqrt{g}V_{}e_{}^{b_{}\varphi }{}_{|x^5=x_{}}{}^{}}`$
where the indices $`\pm `$ refer to the singularities at $`x^5{}_{<}{}^{>}\mathrm{\hspace{0.17em}0}`$. This modifies equations (2.7) and (2.9) of correspondingly. Now let us consider in more detail how this influences solution (I) of . Before discussing the effect of the singularities let us recall formulae which we are going to use in the subsequent considerations. First of all solution (I) is obtained with the ansatz $`A^{}=\frac{1}{3}\varphi ^{}`$ for $`x^5>0`$ ($`x^5<0`$), and the prime denotes derivative with respect to $`x^5`$. The solution for $`\varphi `$ reads
$$\varphi \left(x^5\right)=\{\begin{array}{cc}\frac{3}{4}\mathrm{log}\left|\frac{4}{3}x^5+c_1\right|+d_1,\hfill & x^5<0\hfill \\ \frac{3}{4}\mathrm{log}\left|\frac{4}{3}x^5+c_2\right|+d_2,\hfill & x^5>0,\hfill \end{array}$$
(6)
where the $`c_i,d_i`$ are integration constants. In the following $`c_1`$ is taken to be positive and $`c_2`$ negative. Continuity at $`x^5=0`$ is used to eliminate $`d_2`$. Boundary conditions on the first derivatives at $`x^5=0`$ give equations which can be solved for the $`c_i`$,
$`{\displaystyle \frac{2}{c_2}}`$ $`=`$ $`\left[{\displaystyle \frac{3b}{8}}{\displaystyle \frac{1}{2}}\right]Ve^{bd_1}\left|c_1\right|^{\frac{3}{4}b}`$ (7)
$`{\displaystyle \frac{2}{c_1}}`$ $`=`$ $`\left[{\displaystyle \frac{3b}{8}}+{\displaystyle \frac{1}{2}}\right]Ve^{bd_1}\left|c_1\right|^{\frac{3}{4}b}.`$ (8)
Hence, the solution does exist for any value of $`V`$ and $`b`$ and has one undetermined integration constant.
Now, we want to investigate how putting additional sources at the singularities will modify this observation of self-tuning. By resolving the singularities with additional source terms we get two more boundary conditions at those points ($`x_+=\frac{3}{4}c_2`$, $`x_{}=\frac{3}{4}c_1`$),
$$\frac{8}{3}\left(\varphi ^{}\left(x_\pm +0\right)\varphi ^{}\left(x_\pm 0\right)\right)=b_\pm V_\pm e^{b_\pm \varphi \left(x_\pm \right)}$$
(9)
and
$$3\alpha _{}\left(\varphi ^{}\left(x_\pm +0\right)\varphi ^{}\left(x_\pm 0\right)\right)=\frac{1}{2}V_\pm e^{b_\pm \varphi \left(x_\pm \right)}$$
(10)
with $`\alpha _{}=\frac{1}{3}`$. In order to make sense out of (9),(10) $`\varphi `$ needs to be continued beyond the singularities. There are two (perhaps equivalent) ways of doing that: (a) periodic continuation of the solution or (b) cutting off the fifth direction by defining $`|xx_\pm |=0`$ for $`x_<^>x_\pm `$ . We will follow the second option (b). Technically this means that for $`x^5>0`$ ($`x^5<0`$) we drop the first (second) terms on the left hand sides of (9), (10). We obtain the following conditions
$$b_\pm =\pm \frac{4}{3}$$
(11)
and
$$V_{}e^{\frac{4}{3}d_1}=V_+e^{\frac{4}{3}d_2}=2.$$
(12)
Before plugging in the explicit value of $`d_2`$ we write down the contribution of the singularities to the vacuum energy. There are delta-peaked terms in the bulk Lagrangian and the additional source terms at the singularities. Adding this up the singularities give the following contribution to the four dimensional energy density
$$E_++E_{}=\frac{1}{3}\left(V_+e_{}^{4A+b_+\varphi }{}_{|x^5=x_+}{}^{}+V_{}e_{}^{4A+b_{}\varphi }{}_{|x^5=x_{}}{}^{}\right).$$
(13)
To be specific we choose $`A=\frac{1}{3}\varphi `$ for $`x^5<0`$. Requiring continuity at zero gives $`A=\frac{1}{3}\varphi +\frac{1}{4}\mathrm{log}\left|\frac{c_1}{c_2}\right|+\frac{1}{3}\left(d_1+d_2\right)`$ for $`x_5>0`$. With (11) and (12) and formulae (7) and (8) we get
$$E_++E_{}=\frac{2}{3}e^{\frac{4}{3}d_1}\left(\left|\frac{c_1}{c_2}\right|+1\right)=\frac{2}{3}e^{\frac{4}{3}d_1}\frac{8}{43b}.$$
(14)
The contribution at zero is found to be
$$E_0=\frac{1}{3}Ve_{}^{4A+b\varphi }{}_{|x^5=0}{}^{}.$$
(15)
Using (8) one finds
$$E_0=\frac{2}{3}\frac{8}{43b}e^{\frac{4}{3}d_1}.$$
(16)
Adding up (14) and (16) we obtain for the 4d effective cosmological constant
$$\mathrm{\Lambda }_{4d}=E_++E_{}+E_0=0.$$
(17)
This shows that it is important to take into account contributions from singularities to get vanishing effective cosmological constant.
A remark on self- versus fine-tuning is in order. Our conditions (11), (12) clearly impose fine-tuning on the parameters of the sources at the singularities. The hope of the authors of , is that a different way of resolving the singularities will automatically give the desired contributions to the cosmological constant. One could imagine that at the singularities new light degrees of freedom appear which adjust their vev such that (17) holds. To find a specific example where one can see the detailed dynamics underlying the self tuning would certainly be very interesting.
In what follows we compute the vacuum energy for the solution (II) of ref.
$$\varphi \left(x^5\right)=\{\begin{array}{cc}\pm \frac{3}{4}\mathrm{log}\left|\frac{4}{3}x^5+c\right|+d,\hfill & x^5<0\hfill \\ \pm \frac{3}{4}\mathrm{log}\left|\frac{4}{3}x^5c\right|+d,\hfill & x^5>0.\hfill \end{array}$$
(18)
In the following $`c`$ is taken to be a positive constant. The remaining parameters are $`b=\frac{4}{3}`$ and $`V_0=4e^{\pm \frac{4}{3}d}`$. With this choice the solution has two singularities, one on each side of the brane. If one neglects possible contributions from singularities then one obtains after the calculation analogous to the one in the previous example the vacuum energy
$$E=E_0=\frac{1}{3}V_00.$$
(19)
This is obviously nonzero for a nonzero brane tension $`V_0`$. The trouble is that, again, the solution does not fulfill the Einstein equations at the positions of singularities, so is not a global solution. The point is that the sources supporting the singularities of the Einstein tensor, and that of the gradient energy of the dilaton, are missing. The simplest way to repair the solution is to paste in suitable sources at singularities. Even without worrying about the dynamical origin of these sources, one can easily find out the total contribution from these sources which is needed to repair the solution and to make the vacuum energy vanish. Straightforward calculation shows that the total vacuum energy including contributions from singularities (and corresponding sources) is
$$E=E_0+E_++E_{}=E_0\frac{1}{3}e^{4A(\frac{3}{4}c)}V_+e^{b\varphi (\frac{3}{4}c)}\frac{1}{3}e^{4A(\frac{3}{4}c)}V_{}e^{b\varphi (\frac{3}{4}c)}$$
(20)
where one finds directly from the solution that $`V_+=V_{}=\frac{1}{2}V_0`$. Again, the total vacuum energy vanishes when one includes the total (sources + singularities in the fields) contribution at each singularity. This means that one needs to cut-off the space at singularities by putting a stiff (infinite) potential wall at each singularity, or alternatively one might imagine repeating the whole module consisting of the original brane and two singularities with sources along the fifth dimension.
Next we would like to comment on orbifold examples. First, let us divide a line by $`Z_2`$. Then the acceptable solution must be symmetric around the origin. Since we do not want singularities to appear we take $`\varphi =\frac{3}{4}\mathrm{log}(\frac{4}{3}|x^5|+c)+d`$ with positive $`c`$. When one computes vacuum energy, then one obtains a nonzero, but finite, result. This solution is a valid solution of equations of motion everywhere, hence one would think that one obtains a consistent example of the metric which admits Minkowski-type foliation and gives a nonzero vacuum energy. However, the resolution of the puzzle comes from the observation that the effective four dimensional Planck scale which is proportional to the integral of the warp factor $`\sqrt{\frac{4}{3}|x^5|+c}`$ diverges<sup>2</sup><sup>2</sup>2 The graviton zero mode, which is actually proportional to the warp factor, $`\sqrt{\frac{4}{3}|x^5|+c}`$ is not normalisable.. Thus gravitational degrees of freedom become frozen and effectively gravity decouples from the physics on the brane.
The situation changes when one considers dilaton gravity on the orbifold $`S^1/Z_2`$ extending between $`\pi \rho `$ and $`\pi \rho `$. There the above nonsingular solution can be extended to the solution on this orbifold if one puts on the orbifold plane at $`\pi \rho `$ a system which conspires in such a way as to produce there the brane tension $`V_{\pi \rho }=V_0`$. The vacuum energy of such an orbifold vanishes due to cancellation between contributions from the two branes. However, the correlation between brane tensions on spatially separated branes must be considered to be a fine-tuning, similar to that in the Randall-Sundrum model . We want to point out that the presence of the dynamical dilaton does not decrease the degree of fine-tuning with respect to the Randall-Sundrum model. The important result due to the dilaton is the softening of the dependence of the warp factor, and consequently of the graviton zero mode, on the fifth coordinate. This dependence changes from the exponential fall-off to the mild fractional power-law dependence. The result is that one cannot naturally produce the large hierarchy of mass scales in these quasi-stringy models.
We have shown that in examples where the 5d cosmological constant vanishes there is a non local mechanism of cancellation between the vacuum energy at the brane and at the singularities. (For the low energy four dimensional observer this looks like a vanishing vacuum energy at the brane.) We have argued that this mechanism leads to a hidden fine-tuning even for the self-tuning brane solution. The self-tuning feature can survive only when one finds some dynamical mechanism by which the vacuum energy at the singularity adjusts its value in such a way that it cancels the contribution from the Standard Model brane. Without the knowledge of such a mechanism the self-tuning brane solution seems qualitatively quite similar to the fine tuned orbifold model we discussed above. For this solution it is much simpler to calculate corrections to Newton’s law as there are no subtleties due to a degenerate metric. Work in that direction is in progress.
Acknowledgments
We thank Radosław Matyszkiewicz for useful discussions and Shamit Kachru for correspondence on and .
This work has been supported by TMR programs ERBFMRX–CT96–0045 and CT96–0090. Z.L. is additionaly supported by the Polish Committee for Scientific Research grant 2 P03B 05216(99-2000). |
warning/0002/quant-ph0002072.html | ar5iv | text | # Dynamical Generation of Noiseless Quantum Subsystems
## Abstract
We present control schemes for open quantum systems that combine decoupling and universal control methods with coding procedures. By exploiting a general algebraic approach, we show how appropriate encodings of quantum states result in obtaining universal control over dynamically-generated noise-protected subsystems with limited control resources. In particular, we provide an efficient scheme for performing universal encoded quantum computation in a wide class of systems subjected to linear non-Markovian quantum noise and supporting Heisenberg-type internal Hamiltonians.
03.67.-a,03.67.Lx,03.65.-w,89.70.+c
Quantum bang-bang control has recently emerged as a general strategy for manipulating quantum evolutions by enforcing suitable time scale separations between the controller and the natural dynamics of the system . For open quantum systems, this has lead to establishing quantum error suppression schemes, whereby active decoupling from environmental noise is achieved by continuously undoing system-bath correlations on time scales that are short compared to the typical memory time of the bath . Decoupling techniques were shown to be consistent with efficient quantum information processing , thereby offering an alternative scenario compared to error-correcting and error-avoiding quantum codes . In contrast to the latter methods, no redundant encoding is necessary for preserving or manipulating quantum information provided that the required control operations can be implemented. However, one may ask whether quantum coding could be advantageous or necessary in situations where the available control options are limited.
Answering the above question naturally connects the decoupling formalism with the notion of noiseless subsystem that has been identified as the most general route to noise-free information storage . The basic philosophy is to envision the bang-bang control procedure as a tool for effectively endowing the system dynamics with a nontrivial group of symmetries. Such symmetries generate structures in the system’s state space which are in principle inaccessible to unwanted interactions and are therefore suited for encoding quantum information. Mathematically, the crucial requirement relates to the reducibility properties of operator algebras associated with the action of the decoupling group. Variants of the same basic idea have been argued to lie at the heart of all existing approaches for stabilizing quantum information in a recent work by Zanardi .
In this Letter we examine the implications of the above concept within the decoupling framework, by showing that the action of the control group allows for a complete classification of the choices available for both safe information encoding and universal control over coded states. At variance with the case where noiseless subsystems emerge by virtue of preexisting static symmetries in the overall Hamiltonian, the dynamical origin of the noise-protected structures also precisely constrains the admissible methods for implementing universal control in a way which simultaneously preserves the effect of decoupling as well as the selected coding space.
Using coding methods has several attractive consequences. First, bang-bang operations are needed only for noise suppression. Additional manipulations on encoded subsystems become fully implementable via weak strength controls . Second, for schemes where the relevant Hamiltonians are allowed to be turned on or off slowly, an advantage is that the corresponding pulses can be made more easily frequency-selective. Finally, coded states may be intrinsically more robust against imperfections in the decoupler operations. For a potentially large class of quantum information processors characterized by linear quantum noise, we outline a scheme where noise-decoupling involves a minimal set of two collective bang-bang rotations and universal quantum computation on encoded qubits can be performed entirely through slow tuning of two-body bilinear interactions.
Decoupling.$``$ Let $`S`$ be a finite-dimensional quantum system with self-Hamiltonian $`H_S`$ on $`_S`$, dim$`(_S)=d`$. $`S`$ interacts with the environment $`B`$ via a Hamiltonian $`H_{SB}=_\alpha E_\alpha B_\alpha `$, the $`B_\alpha `$’s being linearly independent environment operators. The error operators $`E_\alpha `$ are assumed to belong to a linear space $``$ that we call the interaction space. We require that tr($`E_\alpha )=0`$, thereby removing from $`H_{SB}`$ the internal evolution of the environment. Let $`𝒜_{}`$ denote the algebra generated by the identity, $`H_S`$, and $``$. $`𝒜_{}`$ is a subalgebra of the full operator algebra End$`(_S)`$ closed under Hermitian transpose (-closed). For $`n`$-qubit systems, $`_S^d`$, End$`(_S)\text{Mat}(d\times d,)`$, with $`d=2^n`$.
In its essence, decoupling via bang-bang (b.b.) control relies on the idea of exploiting full strength/fast switching control actions , meaning that a certain set of Hamiltonians can be (ideally) turned on/off instantaneously with arbitrarily large strength. Let $`𝒢`$ denote a finite group determining the realizable b.b. operations (decoupling group), $`|𝒢|=\text{order}(𝒢)`$. We identify the abstract group $`𝒢`$ with its image under a unitary, faithful representation $`\mu `$ by $`d\times d`$ matrices. A decoupler operates by subjecting the overall system to a cyclical time evolution, the elementary temporal loop (of duration $`T_c`$ = cycle time) being designed to effect a suitable group-theoretical averaging determined by $`𝒢`$. In the ideal limit of arbitrarily fast cycle time, the action of the decoupler is equivalent to a modification of the effective dynamics according to $`𝒜_{}\mathrm{\Pi }_𝒢(𝒜_{})`$, $`\mathrm{\Pi }_𝒢`$ being defined by
$$X\mathrm{\Pi }_𝒢(X)=\frac{1}{|𝒢|}\underset{g𝒢}{}g^{}Xg,X𝒜_{}.$$
(1)
The quantum operation $`\mathrm{\Pi }_𝒢`$ is identical with the projector on the commutant (or centralizer) of $`𝒢`$ in End$`(_S)`$, $`Z(𝒢)=\{𝒪\text{End}(_S)|[𝒪,g]=0g𝒢\}`$. Since $`\mathrm{\Pi }_𝒢(H_S)=H_{eff}Z(𝒢)`$, the decoupler essentially induces a $`𝒢`$-symmetrization of the dynamics due to $`H_S`$.
The commutant $`Z(𝒢)`$ has a natural structure as a subalgebra of End($`_S)`$. A second algebraic structure associated with $`𝒢`$ is the algebra generated by $`𝒢`$, $`𝒢`$, which is the (at most) $`|𝒢|`$-dimensional vector space spanned by complex combinations of elements in $`𝒢`$ . Let $`𝒢^{}`$ denote the set of operators commuting with $`𝒢`$. Clearly, $`𝒢^{}=Z(𝒢)`$. The fact that both $`𝒢`$ and $`𝒢^{}`$ are -closed subalgebras of End($`_S`$) will play an important role. $`𝒢`$ and $`𝒢^{}`$ are linked together by the property of reducibility . $`𝒢`$ is said to be irreducible (and $`𝒢`$ to act irreducibly on $`_S`$) if $`𝒢^{}=\{\lambda 𝟙\}=𝟙`$. A similar definition applies to $`𝒢^{}`$. Since $`𝒢^{\prime \prime }=𝒢`$, the non-triviality of $`𝒢`$ automatically implies that $`𝒢^{}`$ is reducible. Whether or not $`𝒢`$ acts irreducibly on $`_S`$ distinguishes, at the algebraic level, between maximal decoupling, where $`𝒢^{}=𝟙`$, and selective decoupling, in which case $`𝒢^{}𝟙`$ .
The goal of decoupling is to dynamically maintain evolutions of the system so as to have a place where quantum information can safely reside and undergo the required logical manipulations. The possibility to carry out such a program without resorting to redundant encoding was demonstrated in . Is this the only relevant situation?
Encoding.$``$ The basic idea is provided by the notion of a subsystem . Mathematically, subsystems are identified as factors of subspaces by observing that the action of $`𝒢`$ and $`𝒢^{}`$ on $`_S`$ can be represented as
$`𝒢`$ $``$ $`_J𝟙_{𝕟_𝕁}\text{Mat}(𝕕_𝕁\times 𝕕_𝕁,),`$ (2)
$`𝒢^{}`$ $``$ $`_J\text{Mat}(n_J\times n_J,)𝟙_{𝕕_𝕁},`$ (3)
where the index $`J`$ labels the $`J`$-th $`d_J`$-dimensional irreducible component of $`𝒢`$, appearing with multiplicity $`n_J`$. Obviously, $`_Jn_Jd_J=d`$. Such representations are associated with the following decomposition of $`_S`$:
$$_S_J_J_J𝒞_J𝒟_J,$$
(4)
with dim$`(𝒞_J)=n_J`$, dim$`(𝒟_J)=d_J`$. Results (2)-(4) stem from the general decomposition theory of -closed operator algebras. As argued in , they provide the common algebraic ground for discussing noise control strategies. In our setting, the above relationships are linked to the decomposition of $`\mu `$ according to the irreducible representations (irreps) of $`𝒢`$, $`\mu =_Jn_J\mu _J`$ . Eq. (4) reflects the fact that the subspace $`_J`$ of states transforming according to $`\mu _J`$ arises from $`n_J`$ replicas of a $`d_J`$-dimensional irrep. In a suitably chosen orthonormal basis of $`_J`$, $`\{|J,l,m|l=1,\mathrm{},n_J;m=1,\mathrm{},d_J\}`$, such a one-to-one mapping is given by a correspondence of the form $`|l,m|l|m`$. Thus, the $`J`$-th eigenspace factorizes into the tensor product of two factors $`𝒞_J`$ and $`𝒟_J`$, carrying irreps of $`𝒢^{}`$ and $`𝒢`$ respectively. By construction, the dimensions of $`𝒢^{}`$-irreps are found as multiplicities of $`𝒢`$-irreps, and vice versa.
The physical meaning behind the above construction is simple: The overall state space $`_S`$ is decomposed into invariant subspaces $`_J`$, each of which can be regarded as the state space of a bipartite system. For fixed $`J`$, $`𝒞_J`$ is the state space of a subsystem which is only acted on non-trivially by operators in $`𝒢^{}`$, while $`𝒟_J`$ is the state space of a subsystem which is only acted on non-trivially by operators in $`𝒢`$. Clearly, one is left with the freedom of exploiting any of those subsystems for encoding quantum states. Under what conditions is such an encoding noiseless ?
Let us first consider encoding in the left factors $`𝒞_J`$ (“commutant coordinates”), assuming that $`n_J>1`$. When, in situations with underlying static symmetry, the decomposition (2) is applied to the interaction algebra $`𝒜_{}`$, this generalizes the standard case of noiseless subspaces, where coding takes place in the singlet sector of $`𝒜_{}`$, $`d_{J_0}=1`$, $`𝒞_{J_0}_{J_0}`$ . Within the decoupling framework, protection against environmental noise is guaranteed if $`\mathrm{\Pi }_𝒢()=0`$ i.e., $``$ is correctable by $`𝒢`$ . In fact, this condition is no longer necessary and can be replaced by the weaker requirement $`\mathrm{\Pi }_𝒢()𝒢^{}𝒢`$, meaning that the effective error space belongs to the so-called center of $`𝒢`$. Noise suppression is then ensured by the trivial action of the central elements on $`𝒞_J`$, $`𝒢^{}𝒢_Jq_J𝟙_{𝕟_𝕁}𝟙_{𝕕_𝕁}`$, $`q_J`$. Note that $`𝒢^{}𝒢=𝒢`$ for Abelian decouplers.
As a second coding method, we can choose the right factors $`𝒟_J`$ (“group coordinates”). Such an option requires $`d_J>1`$, thereby excluding one-dimensional irreps. As a limiting case, this is the only possibility if $`𝒢`$ acts irreducibly on $`_S`$, in which case the decomposition (2) collapses to a single term $`𝒢\text{Mat}(d\times d,)`$ and the whole space is a noiseless subsystem . In general, since symmetrized noise generators $`\mathrm{\Pi }_𝒢(E_\alpha )𝒢^{}`$ act trivially on factors carrying a $`𝒢`$-irrep, subsystems of the form $`𝒟_J`$ are automatically immune to environmental noise irrespective of the decoupler’s ability to suppress the errors. Although the overall effective dynamics is not unitary in this case, corruption of states in $`𝒟_J`$ is fully prevented due to their symmetry.
In addition to protecting against the environment, encoding may also offer improved stability against faults in the implementations of b.b. control. In particular, while imperfections of operations in $`𝒢`$ directly affect the group component, states that carry $`𝒢^{}`$-coordinates are still unaffected as long as $`𝒢^{}`$ is preserved. Thus, encoding in the commutant degrees of freedom $`𝒞_J`$ is robust against imperfections of the b.b. rotations which stay in $`𝒢`$. Experience from nuclear magnetic resonance suggests that such imperfections do not severely affect the ability of the decoupler to maintain noiselessness of the commutant degrees of freedom. This effect will be analyzed elsewhere.
Universal control.$``$ Since the group-theoretic averaging of the decoupler is intrinsically associated with a minimum time scale $`T_c`$ , it is not surprising that control operations are to be effected according to different timing criteria depending on whether the intended action is on the group or the commutant coordinates. Regardless of the choice of $`𝒞_J`$ or $`𝒟_J`$ as the preferred coding space, transformations over a given subsystem should not be allowed to ever draw states out of the protected factor. This determines the symmetry of the Hamiltonians to be applied for control, $`H𝒢^{}`$ or $`H𝒢`$ for action on $`𝒞_J`$\- or $`𝒟_J`$-subsystems respectively. Since the application of Hamiltonians in $`𝒢^{}`$ does not interfere with the decoupler performances, encoding in $`𝒞_J`$ has the virtue that programming operations can be effected via the weak strength/slow switching scheme introduced in . On the other hand, when encoding in $`𝒟_J`$ is chosen, slow application of arbitrary Hamiltonians produces a trivial action. The least demanding option for applying $`H𝒢𝒢^{}`$ relies then on the ability of fast-modulating $`H`$ according to the weak strength/fast switching scheme of .
Let $`𝒰(𝒞_J)`$ and $`𝒰(𝒟_J)`$ denote the subgroups of unitary transformations over the state space $`𝒞_J`$ and $`𝒟_J`$, respectively. Universality results can be established by observing that, by (2)-(3), $`𝒢^{}|_{𝒞_J}\text{Mat}(n_J\times n_J,)=\text{End}(𝒞_J)`$ and, similarly, $`𝒢|_{𝒟_J}\text{Mat}(d_J\times d_J,)=\text{End}(𝒟_J)`$ i.e., the elements of $`𝒢^{}`$ ($`𝒢`$) restricted to the coding space span the whole operator algebra of the associated subsystem. Thus, by standard universality arguments , almost any pair of Hamiltonians $`H_i𝒢^{}`$ or $`H_i𝒢`$, $`i=1,2`$, is universal over $`𝒞_J`$ or $`𝒟_J`$, respectively. Similar existential results for control over commutant coordinates are formally derived in .
If $`𝒢`$ is irreducible, the possibility to attain complete control over $`_S`$ is directly found as a special case of the above results. When $`𝒢`$ acts reducibly, reachability of arbitrary states in $`_S`$ necessarily occurs through control operations that steer the system through different irreps of $`𝒢`$ and $`𝒢^{}`$. The criteria for universality with no redundant encoding derived in can then be regarded in terms of a symmetry mixing which arises from either combining commutant coordinates associated with different decouplers or from exploiting the action on both group and commutant coordinates of a single group $`𝒢`$.
It is worth stressing that complete controllability of noiseless subsystems does not by itself imply the potential of efficiently implementing a quantum network. This depends on the available physical Hamiltonians as well as on the details of the architecture by which subsystems are actually configured to encode and process information. We focus on quantum computation (QC).
Universal quantum computation.$``$ Let $`S`$ be a quantum computer with $`n`$ qubits, $`_S(^2)^n`$. We consider henceforth a linear interaction Hamiltonian of the form
$$H_{SB}=\underset{a,i}{}\sigma _a^{(i)}B_a^{(i)},$$
(5)
for suitable environment operators $`B_a^{(i)}`$, $`a=x,y,z`$, $`i=1,\mathrm{},n`$. Eq. (5) encompasses various models of interest where the error space is spanned by single-qubit Pauli operators. Notably, the two extreme situations of independent and collective decoherence correspond to error generators of the form $`\{E_\alpha \}=\{\sigma _a^{(i)}\}`$, dim$`()=3n`$, and $`\{E_\alpha \}=\{_i\sigma _a^{(i)}\}`$, dim$`()=3`$, respectively.
Example 1: The collective spin-flips decoupling group. Let us assume that $`n`$ is even and define $`X_j=\sigma _x^{(j)}`$, $`Z_j=\sigma _z^{(j)}`$, with $`Y_j=Z_jX_j=i\sigma _y^{(j)}`$. The group of collective $`\pi `$-rotations is the set $`𝒢=\{𝟙,_{𝕚=\mathrm{𝟙}}^𝕟𝕏_𝕚,_{𝕚=\mathrm{𝟙}}^𝕟𝕐_𝕚,_{𝕚=\mathrm{𝟙}}^𝕟_𝕚\}`$. $`𝒢`$ is an Abelian subgroup of the Pauli group for $`n`$ qubits, with $`k=2`$ generators $`_iX_i`$, $`_iZ_i`$, $`|𝒢|=2^k=4`$. Besides being identical with the stabilizer of distance-two $`[n,n2,2]`$ error-correcting codes , $`𝒢`$ is also a subgroup of the full group of collective rotations that plays the role of a generalized stabilizer for noiseless codes within the collective decoherence model . Decoupling with $`𝒢`$ is effective at suppressing any linear interaction of the form (5) since $`\mathrm{\Pi }_𝒢(\sigma _a^{(i)})=0`$. A single decoupling cycle is specified by a pulse sequence of the form $`[\delta 𝒫_x\delta 𝒫_z]^2`$, $`\delta =T_c/4`$ and $`𝒫_a`$ denoting a time delay and a collective $`\pi `$-pulse along the $`\widehat{a}`$-axis respectively . Since $`𝒢`$ is Abelian, $`𝒢`$ has $`|𝒢|=4`$ one-dimensional irreps and the decomposition of $`_S`$ is identical to the decomposition according to joint eigenspaces $`_J`$, $`J=1,\mathrm{},2^k=4`$, dim($`_J)=n_J=2^{nk}`$. Encoding into commutant factors $`𝒞_J`$ is the only nontrivial option. Accordingly, each of the four (equivalent) joint $`𝒢`$-eigenspaces is able to encode $`n2`$ logical qubits.
Control operations over each $`2^{n2}`$-dim noiseless subspace can be implemented in the weak/slow fashion. Here is an explicit scheme for performing universal QC on encoded qubits. The key point is to look at the available operations in $`𝒢^{}`$, which is easily done by exploiting the isomorphism of $`𝒢`$ with the binary vector space $`𝒵_2^{2n}`$ along with standard results from stabilizer theory . As a group, $`𝒢^{}`$ has a set of $`2n2`$ independent generators, two of which are also generators for $`𝒢`$. The $`2(n2)`$ generators of $`𝒢^{}𝒢`$ can be chosen among interactions of the form $`X_iX_j`$, $`Z_iZ_j`$, $`ij=1,\mathrm{},n`$. These correspond to nontrivial encoded operations. For instance, the choice $`\overline{X}_j=X_1X_{j+1}`$, $`\overline{Z}_j=Z_{j+1}Z_n`$, $`j=1,\mathrm{},n2`$, defines a set of $`n2`$ logical qubits in terms of their encoded $`\sigma _x`$ and $`\sigma _z`$ observables . A universal set of quantum gates is generated by observing that $`𝒢^{}`$ also contains the Heisenberg couplings $`\stackrel{}{\sigma }_i\stackrel{}{\sigma }_j=X_iX_j+Y_iY_j+Z_iZ_j`$ enabling one to implement swapping between any pair of encoded qubits i.e., $`\stackrel{}{\sigma }_{\overline{i}}\stackrel{}{\sigma }_{\overline{j}}=\stackrel{}{\sigma }_{i+1}\stackrel{}{\sigma }_{j+1}`$. Since the square-root-of-swap gate together with one-qubit gates are a universal set , we can noise-tolerantly perform universal QC on $`n2`$ encoded qubits by slowly turning on and off two-body interactions in parallel with the decoupler.
Example 2: The symmetric decoupling group. Let $`𝒢=𝒮_n`$ be the natural representation of the permutation group on the $`n`$-fold tensor product space $`_S`$. In the presence of general linear interactions (5), decoupling according to $`𝒮_n`$ forces effective permutation symmetry, thereby simulating the collective decoherence model . Because $`\mathrm{\Pi }_𝒢()=\{_i\sigma _a^{(i)}\}`$ $`𝒮_n^{}`$, noiseless subsystems are only supported by group factors $`𝒟_J`$ carrying $`𝒮_n`$-irreps. By recalling that $`𝒮_n^{}`$ is identical with the algebra of totally symmetric operators generated by the global $`su(2)`$, the dimensions of such coding spaces can be calculated from the irrep multiplicities of angular momentum theory . Thus, dim($`𝒟_J)=(2J+1)n!/[(n/2+J+1)!(n/2J)!]`$, $`J𝐍/2`$. An explicit scheme has been recently proposed for performing universal QC on logical qubits encoded in clusters of $`n=4,J=0`$ physical qubits . The same construction applies in our setting, with the additional constraint that the exchange Hamiltonians required to implement universal gates should be fast-modulated at the same rate as the b.b. control within a cycle.
Example 3: The collective rotations decoupling group. Let $`𝒢`$ be the continuous group generated by the Lie algebra $`=su(2)`$ of collective spin operators. Decoupling according to $`𝒢`$ can be achieved by performing the quantum operation (1) with respect to a suitable finite-order symmetrizing group of unitaries $``$, whose explicit form is given in . Since $`Z(𝒢)=Z()=𝒮_n`$, $`𝒮_n`$-irreps emerge here as commutant factors, making this example the dual of the previous one. However, being $`\mathrm{\Pi }_𝒢()=\mathrm{\Pi }_{}()=0`$, noiseless subsystems can be supported now by both commutant factors, in which case dim($`𝒞_J)=(2J+1)n!/[(n/2+J+1)!(n/2J)!]`$, or by group factors, for which dim($`𝒟_J)=2J+1`$. In particular, if a $`J=0`$ four-qubits encoding in $`𝒞_J`$ is chosen as above, the scheme for universal QC proposed by can be fully implemented according to weak/slow control.
Discussion.$``$ We presented dynamical procedures for generating and controlling sectors of the state space of a generic open quantum system, which are (ideally) immune to environmental noise. In addition to substantially expanding the range of possibilities for using active decoupling methods, our analysis sheds light on the connections with passive error protection schemes, where the relevant degrees of freedom are decoupled from the noise-inducing interactions by virtue of preexisting symmetries. The presence of nontrivial symmetries is found to be at the root of both active and passive stabilization methods, thereby enabling the identification of common algebraic structures. In spite of the mathematical resemblance, however, the two strategies are physically very different. In particular, the limit of long reservoir correlation length, which underlies passive error prevention in the presence of collective noise , is replaced by the dynamical requirement of long reservoir correlation time in active decoupling, which explicitly relies on the non-Markovian nature of quantum noise . The combination of decoupling and coding procedures results in a scheme for performing universal quantum computation on noise-protected subsystems which is highly appealing in terms of both the attainable encoding efficiency and the overall control resources. Even in the limit where environmental noise is fully tolerated, the scheme is not guaranteed to be robust against arbitrary errors due to imperfect control. The performance of decoupling in the presence of faulty control implementations along with the stability properties of the corresponding dynamically generated subsystems will be discussed in a forthcoming work.
L. V. is grateful to D. P. DiVincenzo for inspiring discussions on stabilizer codes. This work was supported by DARPA/ARO under the QUIC initiative. E. K. received support from the DOE, under contract W-7405-ENG-36, and from the NSA.
vlorenza@mit.edu; knill@lanl.gov; slloyd@mit.edu |
warning/0002/math0002189.html | ar5iv | text | # Geometrically Graded ℎ-𝑝 Quadrature Applied to the Complex Boundary Integral Equation Method for the Dirichlet Problem with Corner Singularities
## 1 Introduction
This report describes a research project carried out from March to October 1992, at the Department of Mathematics, The University of Queensland, Australia. It was carried out under the supervision of Dr Graeme A. Chandler, and was accredited as a #30 project, coded MN882.
Techniques related to the CBIEM have been analysed in , and used in . The CBIEM is also closely related to the ‘Complex Variable Boundary Element Method’ . This report contains an application of it, using $`h`$-$`p`$ quadrature to achieve high rates of convergence, even in the presence of corner singularities. This application owes its conception to my supervisor.
The ideas of graded meshes and $`h`$-$`p`$ quadrature (numerical integration) are presented in §2, and are illustrated by experimental results. The CBIEM itself is described in §3. §4 details numerical implementation of the CBIEM, using the quadrature technique described in §2, and presents error results for some test problems. §5 concludes the report with suggestions for further development. matlab code written for the implementation is listed in Appendix A.
## 2 $`h`$-$`p`$ Quadrature Methods
### 2.1 Introduction
This section describes a high order numerical integration (quadrature) technique, that retains its high order in the case of end point singularities in the integrand. The method uses a *graded mesh*, with integration rules of high order used on larger intervals, and low order on smaller intervals. To achieve the ‘best’ possible convergence rates, whilst including the end points of each interval, the basic quadrature rules used are Gauß–Lobatto. The underlying mesh is graded in a *geometric* manner. As the method of using different quadrature rules on internal intervals is a generalisation of earlier ‘$`h`$’ and ‘$`p`$’ methods, the resulting composite quadrature rule is called a ‘geometrically graded $`h`$-$`p`$’ method .
### 2.2 Quadrature Methods – the Questions
Given an integrand $`f:[a,b]`$, consider the numerical approximation of the definite integral by a rule $`\{x_k,w_k\}`$ on $`n`$ points:
$`{\displaystyle _a^b}f\left(x\right)𝑑x{\displaystyle \underset{k=1}{\overset{n}{}}}f\left(x_k\right)w_k.`$
The interval $`[a,b]`$ is possibly infinite or semi-infinite, but this report considers only finite intervals; and without loss of generality, let $`[a,b]=[0,1]`$. Similarly, the integrand could include a weighting factor $`\omega \left(x\right)`$, but this is not required here.
The *degree* of a quadrature rule is the maximal degree of the polynomial that it can integrate exactly.<sup>1</sup><sup>1</sup>1Comments on errors refer to discretisation, not machine roundoff error unless explicitly stated. That is, if the degree of a rule on $`n`$ points is $`p`$, then:
$`{\displaystyle _0^1}x^j𝑑x={\displaystyle \underset{k=1}{\overset{n}{}}}x_k^jw_k,j=0:p.`$
If $`f`$ is smooth, the rate of convergence for $`n`$ point Gaußian quadrature is $`𝒪\left(\rho ^n\right)`$ (for some $`\rho <1`$), and for the composite Simpson’s rule it is $`𝒪\left(n^4\right)`$. That is, the error decreases more quickly for Gaußian quadrature. This is *not* true in general if $`f`$ has a singularity.<sup>2</sup><sup>2</sup>2‘Singularity’ is intended to always mean ‘end point singularity’. If a particular singularity in the integrand is not at an end point, then the interval can be subdivided so that the singularity is at the end points of the two subintervals. Most quadrature methods perform poorly on integrands with internal singularities. For example, consider the ‘square root’ singularity $`f(x)=\sqrt{x},x[0,1]`$. In this case, the rate of convergence for Gaußian quadrature falls to $`𝒪\left(n^3\right)`$, whilst that of Simpson’s rule is $`𝒪\left(n^{3/2}\right)`$. Even so, using a *composite* Simpson’s rule and a *graded mesh*, a convergence rate of $`𝒪\left(n^4\right)`$ can be recovered.
A *composite* quadrature rule is created by subdividing the interval of integration into $`m`$ subintervals, and evaluating the integral over each subinterval using an appropriate quadrature rule. Choosing $`x_{j1}<x_j,j=1:m`$, and $`x_0=0,x_m=1`$:
$`{\displaystyle _0^1}f\left(x\right)𝑑x={\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle _{x_{j1}}^{x_j}}f\left(x\right)𝑑x.`$
*Grading* the mesh means that the subdivision is organised in some way. A description of a generalised mesh grading to cater for complicated possibilities, such as *adaptive* quadrature is found in . Here, the simplification of nonadaptive meshes is used. Meshes are graded by assigning mesh points according to some simple function. A quadrature rule of degree $`p_j`$ (a function of $`n_j`$, the number of points used by the rule) is used on the interval $`[x_{j1},x_j]`$. This raises two issues:
* How should $`x_j`$ and $`n_j`$ be chosen? That is, how should the mesh be graded, and how should the degree of the quadrature rules on each subinterval vary?
* How is this procedure dependent on the integrand? Consider the generalisation of $`\sqrt{x}`$ to $`|x|^\alpha `$, or more pathological cases. For experimental work, a known integrand is easy to deal with. What of more general cases, where no explicit functional information is known?
The remainder of this section describes some partial answers to these questions, and displays some experiments. The answers deal with the case $`|x|^\alpha ,1<\alpha <1`$, and the experiments demonstrate the case $`\alpha =1/2`$.
### 2.3 Gaußian Quadrature
Gaußian quadrature rules are the *best* possible rules in the sense that they are of maximal degree.<sup>3</sup><sup>3</sup>3This may of course, not be ideal for the particular application, but in the absence of theoretical functional information about the integrand, nothing beyond this can be said about the convergence rates of *any* quadrature method. In practice, with commonly occurring functions, there is a certain amount of implicit theoretical information which can be used in error analysis. This is due to the exploitation of the maximal number of degrees of freedom in the choice of their nodes and weights. Gauß methods divide into categories depending on the associated weight function, and whether there are any prescribed quadrature points. For the applications in this report, it is preferable to use the end points of the interval as quadrature points, and a unit weight function is assumed. The appropriate set of rules are called the Gauß–Lobatto rules \[6, pages 101–104\].
###### Theorem 2.1 (Gauß–Lobatto Quadrature)
Given $`fC^{2n2}[a,b]`$, the $`n`$ point Gauß–Lobatto quadrature rule ($`n2`$), has nodes $`ax_1<x_2<\mathrm{}<x_{n1}<x_nb`$, and positive weights $`w_1,\mathrm{},w_n`$ such that:
$`{\displaystyle _a^b}f\left(x\right)𝑑x={\displaystyle \underset{k=1}{\overset{n}{}}}f\left(x_k\right)w_k+E_n.`$
Here $`E_n`$ is dependent on $`f`$, $`a`$, $`b`$ and $`n`$:
$$E_n=\frac{n\left(n+1\right)\left[\left(n2\right)!\right]^4}{\left(2n1\right)\left[\left(2n2\right)!\right]^3}\left(ba\right)^{2n1}f^{\left(2n2\right)}\left(\xi \right),\xi (a,b).$$
(1)
The rule is of degree $`p=2n3`$ (this is always odd). Observe that for $`n=2`$ the rule is the trapezoidal rule, and for $`n=3`$ it is Simpson’s rule. Only for $`n4`$ do these rules diverge from the series of closed Newton–Cotes rules (see Table 1 on page 1).
An algorithm for finding $`\{x_k,w_k\}`$ using a matrix eigenvalue technique is implemented in Appendix A.9.
### 2.4 Graded Meshes
A number of different methods for grading meshes appear in the literature. Three important methods are described below. In each case, the mesh subdivides the interval $`[0,1]`$, when the integrand has a singularity at $`0`$. The meshes have $`m`$ subintervals, that is $`m+1`$ points, including the ends. The $`j`$th mesh point is at $`x_j`$, and the $`j`$th interval is of width $`h_j`$:
1. *Quasiuniform*. The mesh is essentially uniform; that is for some constant $`\tau <1`$, $`h_j[h\tau ,h],j=1:m1`$, where $`h=\mathrm{max}_j\left\{h_j\right\}`$.
2. *Algebraic*. For some $`\gamma 1`$, $`x_j=\left({\displaystyle \frac{j}{m}}\right)^\gamma `$, $`j=0:m`$.
3. *Geometric*. For some $`0<\sigma <1`$, $`x_j=\sigma ^{mj}`$, $`j=1:m`$; and $`x_0=0`$.
A geometrically graded mesh is illustrated on one segment of a closed contour in Figure 4 on page 4.
### 2.5 $`h`$, $`p`$ and $`h`$-$`p`$ Quadrature Methods
The ‘$`h`$-$`p`$’ nomenclature presented here originated in papers by Babuška et al. , on finite element methods, based on previous work which did not explicitly use this schema. The following discussion of the three methods refers to their use with *graded* meshes.
#### 2.5.1 $`h`$ Methods
An $`h`$ quadrature method is composed using two steps:
1. Choose an underlying mesh of subintervals; possibly a graded mesh determined by the user, from analysis of the singularities of the integrand.
2. Integrate over each of the $`m`$ mesh intervals, applying the same $`n`$ point quadrature rule. The result is a composite quadrature rule on a total of $`N`$ points. These $`N`$ points will be called the *node* points from now on. The functional relationship $`N(m,n)`$ is dependent on whether the basic quadrature rule is open or closed. (If the rule is open, the original mesh points are not included in the final rule.) Observe that the user cannot arbitrarily select $`N`$, only $`m`$ and $`n`$.
$`N(m,n)=\{\begin{array}{cc}m\left(n2\right)+m+1\hfill & (\mathrm{closed}\mathrm{basic}\mathrm{rule})\hfill \\ mn\hfill & (\mathrm{open}\mathrm{basic}\mathrm{rule}).\hfill \end{array}`$
A particular basic rule is decided upon (e.g. Simpson’s rule), and desired accuracy is hopefully attained by simply increasing $`m`$ (that is, $`N`$). Whatever grading is chosen, the separation of the node points ($`h`$) decreases as $`m`$ is increased, hence the name ‘$`h`$ method’. For a uniform mesh (which works well for smooth integrands), $`h=\left(ba\right)/\left(N1\right)`$ is constant. Alternatively, open rules, or Gauß rules can be used, the only important factor is that all the basic rules are of the same type and degree.
#### 2.5.2 $`p`$ Methods
In a $`p`$ method, again a graded mesh is created. Integration is performed over each mesh interval using a basic quadrature rule on $`n`$ points. Here, $`n`$ instead of $`m`$ is varied by the user. That is, for a given number of mesh subintervals, $`m`$, a set of rules of increasing degree (that is $`n`$, the number of points involved) is used, until desired accuracy is obtained. The same functional relationship $`N(m,n)`$ exists. As $`n`$ is increased, the rules used grow in their degree ($`p`$), hence the name ‘$`p`$ method’ (see also Table 1).
To illustrate, consider the family of closed Newton–Cotes rules. Assume that the interval has been subdivided, possibly using an adaptive algorithm that chooses smaller subdivisions where there the integrand has greater derivative. Approximate the integral over each division using the trapezoidal rule ($`p=1`$), and inspect the result. If it is unacceptable, repeat using Simpson’s rule ($`p=3`$). Continue this process until results are acceptable.
#### 2.5.3 $`h`$-$`p`$ Methods
The $`h`$-$`p`$ method is the natural combination of the two previous methods. The user may vary both $`m`$ and $`n`$. The idea behind this is to create a composite rule that minimises errors in the approximation, for a given number of node points $`N`$. (Experiment demonstrates that this is achievable.) The user chooses a family of basic quadrature rules, then decides how to vary $`n`$ with mesh interval. As the singularities considered are *always* at end points, a good choice is to organise small mesh intervals and low degree rules (small $`n`$) near the end points, and larger mesh intervals and high degree rules away from them, where the integrand is expected to be smooth.
A simple choice is to begin with a rule on $`n=2`$ points on the smallest interval, then linearly increase $`n`$ with the number of the mesh interval. Other discrete integer functions $`n_j,j=1:m`$ are easily designed. The only constraint on these functions is that if any error analysis is to be done, there should be some regularity in $`n_j`$. (Choosing basic rules from the same family facilitates this.) This implementation uses the Gauß–Lobatto rule of degree $`1`$ ($`n=2`$) on the first interval, degree $`3`$ ($`n=3`$) on the second, etc.
Creating the composite quadrature rule is quite difficult. Each of the basic quadrature rules must be appropriately scaled and shifted, and then coincident mesh points must be combined. This is further complicated in the cases of closed meshes, closed quadrature rules, and contour integration, where the end points of various segments of the parameterisation must also be combined. (This is exacerbated if the contour is closed.) The CBIEM requires all of these to be implemented. The (closed) contours involved have corners, and the integrand will usually have singularities at these corners. As it will be important to keep the collocation points (see §3.2.3) between, and not on, the corners, the underlying meshes *must* include end points. This means that the basic quadrature rules must be closed, so as to include the end points.
The literature recommends using a geometrically graded mesh, with an $`h`$-$`p`$ quadrature method. (This is implemented in the CBIEM.) For maximum efficacy, the basic quadrature rules chosen must be of maximal degree, which restricts them to Gauß rules. They must also be closed. An $`n`$ point rule already has two of its points fixed, at the ends. The appropriate rule is known as the Gauß–Lobatto rule, which is of degree $`p=2n3`$.
### 2.6 Error Analysis for the $`h`$ and $`h`$-$`p`$ Methods
This section is tedious, and consists mainly of technical arguments. The important parts are Theorem 2.2 on page 2.2; Theorem 2.3 on page 2.3; and the experimental results in §§2.6.3 and 2.6.4. The rest can be skipped without loss of continuity.
#### 2.6.1 Error Analysis for the $`h`$ Method
This section computes an error bound for the $`h`$ method using an algebraic mesh, for the integrand $`|x|^\alpha `$, on $`[0,1]`$. Clearly $`\alpha >1`$ is necessary for the integral to be proper, and thus make its computation sensible. For $`1<\alpha <0`$, $`|x|^\alpha C^0[0,1]`$, so the integrand is unbounded, but the integral is nonetheless defined. If $`0\alpha <1`$, $`|x|^\alpha C^0[0,1]`$, but $`|x|^\alpha C^1[0,1]`$. If $`\alpha `$, then $`|x|^\alpha C^{\mathrm{}}[0,1]`$, so the singularity vanishes and the case is of lesser interest. If $`\alpha >1`$, but $`\alpha `$, then all of the higher derivatives at $`x=0`$ will not exist. Using the notation $`x`$ and $`x`$ for the least integers (respectively) greater than and less than $`x`$, in general, for $`\alpha >1`$, $`|x|^\alpha C^\alpha [0,1]`$, but $`|x|^\alpha C^\alpha [0,1]`$. These cases are not particularly interesting, so the limit $`\alpha <1`$ is made for simplicity. That is, consider $`1<\alpha <1`$, which includes the paradigm example $`x^{1/2}`$. The $`n`$th derivative of $`f\left(x\right)=|x|^\alpha `$, for $`x0`$, is:
$`f^{(n)}\left(x\right)={\displaystyle \frac{\mathrm{\Gamma }\left(\alpha +1\right)}{\mathrm{\Gamma }\left(\alpha +1n\right)}}|x|^{\alpha n}.`$
Consider the interval $`[0,1]`$ partitioned into $`m`$ subintervals, where $`x_j`$ is the $`j`$th mesh point, $`j=0:m`$, and $`h_j=x_jx_{j1}`$ is the width of the $`j`$th interval, for $`j=1:m`$. Recall, for an *algebraic* grading, a real constant $`\gamma 1`$ is chosen,<sup>4</sup><sup>4</sup>4Choosing $`\gamma <1`$ results in some mesh points possibly lying outside the interval $`[0,1]`$. and the mesh is defined by: $`x_j=\left(j/m\right)^\gamma `$, $`j=0:m`$. Differentiating and applying the mean value theorem shows:
$`h_j=\left({\displaystyle \frac{j}{m}}\right)^\gamma \left({\displaystyle \frac{j1}{m}}\right)^\gamma ={\displaystyle \frac{1}{m^\gamma }}\left[j^\gamma \left(j1\right)^\gamma \right]{\displaystyle \frac{\gamma }{m}}\left({\displaystyle \frac{j}{m}}\right)^{\gamma 1}={\displaystyle \frac{dx_j}{dj}}.`$
For a *geometric* grading, for some constant $`0<\sigma <1`$, $`x_j=\sigma ^{mj}`$, $`j=1:m`$; $`x_0=0`$, so:
$`h_j=\sigma ^{mj}\sigma ^{mj+1}=(1\sigma )\sigma ^{mj},j=2:m,h_1=\sigma ^{m1}.`$
Consider an algebraic grading, using closed basic quadrature rules. For an $`h`$ or $`h`$-$`p`$ method, the integral on the $`j`$th interval is computed using a quadrature rule with $`n_j2`$ points. For an $`h`$ method, $`n_j`$ is constant; for instance, $`n_j2`$ means that the integral over each mesh interval is computed using the trapezoidal rule – the rule is a *composite* trapezoidal rule. The degrees of some common quadrature rules are presented in Table 1.
Now consider the global error $`E_m`$ of the difference between the true solution $`I`$ and the approximation $`I_m`$, induced by the quadrature on $`m`$ subintervals, that is $`I=I_m+E_m`$. For the function $`|x|^\alpha `$, bounds on $`E_m`$ are readily found. Define $`e_j`$ as the component of $`E_m`$ due to the $`j`$th mesh interval, that is $`E_m=_{j=1}^me_j`$. To bound $`E_m`$, note $`E_m_{j=1}^m|e_j|`$, and then bound each of the $`|e_j|`$.
The error result for a degree $`p`$ quadrature rule on an interval $`[x_{j1},x_j]`$, of width $`h_j`$, with a function $`f\left(x\right)C^{p+1}[x_{j1},x_j]`$, is:
$$e_j=C\left(p\right)h_j^{p+2}f^{\left(p+1\right)}\left(\xi \right),\xi [x_{j1},x_j].$$
(3)
For $`n`$ point Gauß–Lobatto quadrature, the degree is $`p=2n3`$. (3) is derived from (1) by a scaling argument, and:
$$C\left(p\right)=\frac{\left(p+3\right)\left(p+5\right)\left[\left(\left(p1\right)/2\right)!\right]^4}{2^2\left(p+2\right)\left[\left(p+1\right)!\right]^3}.$$
(4)
As $`|x|^\alpha C^{\mathrm{}}(0,1]`$, (3) holds for every mesh interval except the first, where the error is known exactly, e.g. for the trapezoidal rule:
$$e_1=\frac{h_1^{\alpha +1}}{\alpha +1}\frac{h_1}{2}h_1^\alpha =\left(\frac{1}{\alpha +1}\frac{1}{2}\right)h_1^{\alpha +1}.$$
(5)
As the maximum value of the $`\left(p+1\right)`$th derivative of $`|x|^\alpha `$ on the $`j`$th interval, is at its left hand end, $`f^{\left(p+1\right)}\left(\xi \right)f^{\left(p+1\right)}\left(x_{j1}\right)`$, the total error can be bounded:
$`E_m|e_1|+C{\displaystyle \underset{j=2}{\overset{m}{}}}h_j^{p+2}f^{\left(p+1\right)}\left(x_{j1}\right)|e_1|+C{\displaystyle \underset{j=1}{\overset{m1}{}}}h_{j+1}^{p+2}f^{\left(p+1\right)}\left(x_j\right).`$
Here, $`C`$ refers to a positive constant, independent of $`m`$, that may vary from line to line. Substituting for $`h_{j+1}`$ and $`f^{\left(p+1\right)}\left(x_j\right)`$ using an algebraically graded mesh yields:
$`E_m`$ $``$ $`|e_1|+C{\displaystyle \underset{j=1}{\overset{m1}{}}}\left[{\displaystyle \frac{\gamma }{m}}\left({\displaystyle \frac{j1}{m}}\right)^{\gamma 1}\right]^{p+2}{\displaystyle \frac{\mathrm{\Gamma }\left(\alpha +1\right)}{\mathrm{\Gamma }\left(\alpha p\right)}}\left({\displaystyle \frac{j}{m}}\right)^{\gamma \left(\alpha p1\right)}`$
$``$ $`|e_1|+C{\displaystyle \underset{j=1}{\overset{m1}{}}}\left[{\displaystyle \frac{\gamma }{m}}\left({\displaystyle \frac{j}{m}}\right)^{\gamma 1}\right]^{p+2}\left({\displaystyle \frac{j}{m}}\right)^{\gamma \left(\alpha p1\right)}.`$
As $`|e_1|`$ is less than some constant multiplied by the ‘$`m`$th term’ in the sum:
$$E_m\frac{C}{m^{\gamma \left(\alpha +1\right)}}\underset{j=1}{\overset{m}{}}j^{\gamma \left(\alpha +1\right)\left(p+2\right)}.$$
(6)
Simplification of (6) (see below) leads to the result:
###### Theorem 2.2 (Convergence of the $`h`$ Method with Algebraic Grading)
Consider the approximation of $`_0^1|x|^\alpha 𝑑x`$, $`1<\alpha <1`$, using an $`h`$ method based on a quadrature rule of degree $`p`$ (an odd positive integer), on an algebraic mesh on a total of $`m2`$ intervals, with mesh parameter $`\gamma 1`$. For some constant $`C`$, the error $`E_m`$ satisfies:
$`E_mCm^z.`$
Here $`z`$ is:
$`z=\{\begin{array}{cc}\gamma \left(\alpha +1\right)\hfill & 1\gamma <\left(p+1\right)/\left(\alpha +1\right)\hfill \\ p+1\hfill & \mathrm{else}.\hfill \end{array}`$
When $`\gamma =\left(p+1\right)/\left(\alpha +1\right)`$, $`E_mC\mathrm{ln}\left(m\right)/m^{p+1}`$.
###### Theorem 2.1
Take (6) and write $`E_m`$ as:
$`E_m`$ $``$ $`{\displaystyle \frac{C}{m^{\gamma \left(\alpha +1\right)}}}{\displaystyle \underset{j=1}{\overset{m}{}}}j^{\gamma \left(\alpha +1\right)\left(p+2\right)}{\displaystyle \frac{C}{m^{\gamma \left(\alpha +1\right)}}}{\displaystyle _1^m}x^{\gamma \left(\alpha +1\right)\left(p+2\right)}𝑑x`$
$``$ $`{\displaystyle \frac{C}{m^{\gamma \left(\alpha +1\right)}}}{\displaystyle _1^{\mathrm{}}}x^{\gamma \left(\alpha +1\right)\left(p+2\right)}𝑑x.`$
The integral converges if: $`\gamma \left(\alpha +1\right)\left(p+2\right)<1`$, and in this case, it converges to the constant $`\left[\left(p+1\right)\gamma \left(\alpha +1\right)\right]^1`$, independent of $`m`$. Absorbing this into the main constant, the result for $`1\gamma <\left(p+1\right)/\left(\alpha +1\right)`$ is created.
In the case $`\gamma =\left(p+1\right)/\left(\alpha +1\right)`$, (6) is
$`E_m{\displaystyle \frac{C}{m^{p+1}}}{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{1}{j}}{\displaystyle \frac{C}{m^{p+1}}}{\displaystyle _1^m}{\displaystyle \frac{1}{x}}𝑑x{\displaystyle \frac{C}{m^{p+1}}}\mathrm{ln}\left(m\right).`$
Thus again, the error is bounded by a constant, this time, dependent on $`m`$. Observe that $`p+1`$ is immediately able to be replaced with $`\gamma \left(\alpha +1\right)`$, demonstrating the continuity of the formulae.
Lastly, bound $`E_m`$ as:
$`E_m`$ $``$ $`{\displaystyle \frac{C}{m^{p+2}}}{\displaystyle \underset{j=1}{\overset{m}{}}}\left({\displaystyle \frac{j}{m}}\right)^{\gamma \left(\alpha +1\right)\left(p+2\right)}{\displaystyle \frac{C}{m^{p+2}}}{\displaystyle _{1/m}^1}x^{\gamma \left(\alpha +1\right)\left(p+2\right)}𝑑x`$
$``$ $`{\displaystyle \frac{C}{m^{p+2}}}{\displaystyle _0^1}x^{\gamma \left(\alpha +1\right)\left(p+2\right)}𝑑x.`$
As $`\gamma \left(\alpha +1\right)\left(p+2\right)>1`$, the integral converges to $`\left[\gamma \left(\alpha +1\right)\left(p+1\right)\right]^1`$, again, another constant independent of $`m`$, which is absorbed into $`C`$. The second result is thus achieved, by observing that $`m^{\left(p+2\right)}<m^{\left(p+1\right)}`$.
The moral of this is that for a particular choice of $`\alpha `$ and $`p`$, there is an ideal choice of $`\gamma `$, that is $`\gamma ^{}=\left(p+1\right)/\left(\alpha +1\right)`$, beyond which the order of the error will not decrease. (Choosing $`\gamma `$ greater than this may reduce $`C`$.) Varying $`\gamma `$ makes no difference to computational expense.<sup>5</sup><sup>5</sup>5Choosing $`\gamma ^{}`$ may be cheaper than nonintegral choices of $`\gamma ^{}`$. Observe that $`\gamma ^{}`$ becomes unbounded as $`\alpha 1^+`$. This is not surprising, as the integral itself becomes unbounded. A similar result to Theorem 2.2 exists for a geometric mesh.
In summary, using an algebraic mesh of $`m`$ subintervals, and an $`h`$ method with a degree $`p`$ quadrature rule on each mesh interval, can achieve, for $`|x|^\alpha `$ integrands, an $`𝒪\left(m^{(p+1)}\right)`$ convergence rate. This is a significant improvement on the equivalent result for a uniform mesh, which is $`𝒪\left(m^{\left(\alpha +1\right)}\right)`$. The exponential convergence rate for smooth integrands is not achieved, but can be, using an $`h`$-$`p`$ method. The methods may be extended to integrands of the form $`|x|^\alpha f\left(x\right)`$, for smooth functions $`f`$.
#### 2.6.2 Error Analysis for the $`h`$-$`p`$ Method
This section discusses the expected order of the error for the $`h`$-$`p`$ method, for algebraic or geometric meshes using integrands with end point singularities. Again, consider the integral $`_0^1|x|^\alpha 𝑑x`$. For the $`h`$ method, $`p_j`$, the degree of the quadrature rule on the $`j`$th interval, was constant. For the $`h`$-$`p`$ method, it is a function of $`j`$. A low degree rule is used on the interval adjacent to the end point singularity; and higher degree rules are used on intervals away from it. A simple choice is to use a rule on $`2`$ points on the first interval, and increase the number of points linearly with $`j`$. That is, using Gauß–Lobatto rules, where $`p_j=2n_j3`$; choosing $`n_j=j+1`$ gives $`p_j=2j1`$.
Recall the error result from (3), where $`C_j=C\left(p_j\right)`$ is given by (4):
$`{\displaystyle _{x_{j1}}^{x_j}}f\left(x\right)𝑑xI_{p_j}=e_j=C_jh_j^{p+2}f^{\left(p+1\right)}\left(\xi \right),\xi [x_{j1},x_j].`$
Use Stirling’s formula for large $`x`$ to approximate the factorials in $`C_j`$, as $`\left(x1\right)!=\mathrm{\Gamma }\left(x\right)`$:
$`\mathrm{\Gamma }\left(x\right)\sqrt{{\displaystyle \frac{2\pi }{x}}}\left({\displaystyle \frac{x}{e}}\right)^x\left\{1+{\displaystyle \frac{1}{12x}}+\mathrm{}\right\}.`$
Applying the first term of this to (4) gives an asymptotic bound for large $`j`$:
$`C_j`$ $`=`$ $`{\displaystyle \frac{\left(j+1\right)\left(j+2\right)}{2j+1}}{\displaystyle \frac{\left[\left(j1\right)!\right]^4}{\left[\left(2j\right)!\right]^3}}{\displaystyle \frac{2^7}{e^3}}\sqrt{\pi j}\left({\displaystyle \frac{e}{8j}}\right)^{2j+3}.`$
Thus:
$$|C_j|C\frac{\mu ^j}{j^{5/2+2j}},\mu =e^2/2^6.$$
(8)
$`E_m`$ cannot be bounded directly using this approximation for $`C_j`$, as it is not a proper bound. Heuristically,<sup>6</sup><sup>6</sup>6This is brought out by experiment. for either a geometric or an algebraic mesh, the rapid convergence of $`C_j`$ to $`0`$ means that $`E_m`$ is expected to be dominated by $`e_1`$. Consider a mesh on $`m`$ intervals, and an associated composite quadrature rule on a total of $`N=m\left(m+1\right)/2+1`$ points. For a geometric mesh $`h_1=\sigma ^{m1}`$, and for an algebraic mesh $`h_1=m^\gamma `$. Using (5), this gives:
$$|e_1|Ch_1^{\alpha +1}C\{\begin{array}{cc}m^{\gamma \left(\alpha +1\right)}\hfill & \mathrm{algebraic}\hfill \\ \sigma ^{\left(m1\right)\left(\alpha +1\right)}\hfill & \mathrm{geometric}.\hfill \end{array}$$
(9)
The error for an algebraic mesh is polynomial, whilst the error for a geometric mesh is exponential. If the errors with increasing $`Nm^2/2`$ are plotted for a geometric mesh, an error of the form $`E_mC\rho ^\sqrt{N}`$ is observed (see Figure 2), for some $`\rho <1`$.
These rough results prompt more rigorous examination of $`E_m`$. Recall:
$`|e_j|`$ $``$ $`C_jh_j^{2j+3}f^{\left(2j+2\right)}\left(x_{j1}\right).`$
For a geometric mesh, $`x_j=\sigma ^{mj}`$, so $`h_j<\left(1\sigma \right)\sigma ^{mj}`$. Using $`f\left(x\right)=|x|^\alpha `$, gives:
$`f^{(n)}\left(x\right)={\displaystyle \frac{\mathrm{\Gamma }\left(\alpha +1\right)}{\mathrm{\Gamma }\left(\alpha +1n\right)}}|x|^{\alpha n}.`$
The Stirling asymptotic approximation for $`C_j`$ can be converted to a genuine bound by observing for $`x`$:
$`\sqrt{{\displaystyle \frac{2\pi }{x}}}\left({\displaystyle \frac{x}{e}}\right)^x`$ $``$ $`\mathrm{\Gamma }\left(x\right)C\sqrt{{\displaystyle \frac{2\pi }{x}}}\left({\displaystyle \frac{x}{e}}\right)^x.`$
Recall from (8), that with $`\mu =e^2/2^6`$, for some $`C`$ independent of $`x`$:
$`|C_j|`$ $``$ $`C{\displaystyle \frac{\mu ^j}{j^{2j+5/2}}}.`$
A bound on the error for the $`j`$th interval is now:<sup>7</sup><sup>7</sup>7Strictly speaking, this only applies for $`j=2:m`$, as the result for $`C_j`$ only holds for $`j=2:m`$. Application of the result in (9) allows $`|e_1|`$ to be bounded by a constant multiple of this $`C_1`$.
$`|e_j|`$ $`<`$ $`C{\displaystyle \frac{\mu ^j}{j^{2j+5/2}}}\left[\left(1\sigma \right)\sigma ^{mj}\right]^{2j+3}f^{\left(2j+2\right)}\left(\sigma ^{mj1}\right).`$
Simplification of this leads to a bound on $`E_m`$. Observe that $`\left(1\sigma \right)^{2j+3}<1`$, so:
$`|e_j|`$ $`<`$ $`C{\displaystyle \frac{\mu ^j}{j^{2j+5/2}}}\sigma ^{\left(mj\right)\left(2j+3\right)}{\displaystyle \frac{\mathrm{\Gamma }\left(\alpha +1\right)}{\mathrm{\Gamma }\left(\alpha 2j1\right)}}\sigma ^{\left(mj1\right)\left(\alpha 2j2\right)}.`$
Absorb $`\mathrm{\Gamma }\left(\alpha +1\right)`$ into $`C`$, and use the Stirling approximation for $`\mathrm{\Gamma }\left(\alpha 2j1\right)`$:
$`{\displaystyle \frac{1}{\mathrm{\Gamma }\left(\alpha 2j1\right)}}`$ $``$ $`\sqrt{{\displaystyle \frac{\alpha 2j1}{2\pi }}}\left({\displaystyle \frac{e}{\alpha 2j1}}\right)^{\alpha 2j1}.`$
Rearranging the exponent of $`\sigma `$, and absorbing the term $`e^{\alpha 1}`$ into $`C`$:
$`|e_j|`$ $`<`$ $`C{\displaystyle \frac{\mu ^j}{j^{2j+5/2}}}\sigma ^{\left(mj\right)\left(\alpha +1\right)\alpha +2j+2}\sqrt{{\displaystyle \frac{\alpha 2j1}{2\pi }}}\left({\displaystyle \frac{e}{\alpha 2j1}}\right)^{\alpha 2j1}`$
$`<`$ $`C{\displaystyle \frac{\mu ^j\sigma ^{m\left(\alpha +1\right)}\sigma ^{j\left(1\alpha \right)}e^{2j}}{j^{2j+5/2}\left(\alpha 2j1\right)^{\alpha 2j1/2}}}.`$
Expand $`\mu =e^2/2^6`$, and observe that $`\left(\alpha 2j1\right)^{\alpha 1/2}>1`$:
$`|e_j|`$ $`<`$ $`C{\displaystyle \frac{e^{2j}\sigma ^{m\left(\alpha +1\right)}\sigma ^{j\left(1\alpha \right)}e^{2j}}{j^{2j+5/2}2^{6j}\left(\alpha 2j1\right)^{2j}}}<C{\displaystyle \frac{\sigma ^{m\left(\alpha +1\right)}\sigma ^{j\left(1\alpha \right)}}{2^{6j}j^{5/2}}}\left({\displaystyle \frac{\alpha 2j1}{j}}\right)^{2j}.`$
As $`|\left(\alpha 2j1\right)/j|<2`$, then:
$`|e_j|`$ $`<`$ $`C{\displaystyle \frac{2^{2j}\sigma ^{m\left(\alpha +1\right)}\sigma ^{j\left(1\alpha \right)}}{2^{6j}j^{5/2}}}<C{\displaystyle \frac{\sigma ^{m\left(\alpha +1\right)}\sigma ^{j\left(1\alpha \right)}}{2^{4j}j^{5/2}}}.`$
Combine these, to bound:
$`E_m`$ $``$ $`{\displaystyle \underset{j=1}{\overset{m}{}}}|e_j|<C{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{\sigma ^{m\left(\alpha +1\right)}\sigma ^{j\left(1\alpha \right)}}{2^{4j}j^{5/2}}}<C\sigma ^{m\left(\alpha +1\right)}{\displaystyle \underset{j=1}{\overset{m}{}}}{\displaystyle \frac{\sigma ^{j\left(1\alpha \right)}}{2^{4j}j^{5/2}}}.`$
Observing that $`2^{4j}j^{5/2}>1`$, this can be written as:<sup>8</sup><sup>8</sup>8This throws away a lot of information!
$`E_m<C\sigma ^{m\left(\alpha +1\right)}{\displaystyle \underset{j=1}{\overset{m}{}}}\left[\sigma ^{\left(1\alpha \right)}\right]^j.`$
Using the result for the sum of a geometric progression:
$`{\displaystyle \underset{j=1}{\overset{m}{}}}\left[\sigma ^{\left(1\alpha \right)}\right]^j`$ $`=`$ $`{\displaystyle \frac{\left(\sigma ^{1\alpha }\right)^{m+1}1}{\sigma ^{1\alpha }1}}={\displaystyle \frac{\sigma ^{\left(1\alpha \right)\left(m+1\right)}1}{\sigma ^{1\alpha }1}}<C\sigma ^{\left(1\alpha \right)\left(m+1\right)}.`$
Thus: $`E_m<C\sigma ^{m\left(\alpha +1\right)}\sigma ^{\left(1\alpha \right)\left(m+1\right)}<C\sigma ^{2m}`$. As $`Nm^2/2`$, using $`\rho =\sigma ^{2\sqrt{2}}`$, this simplifies to:
$`E_m<C\rho ^\sqrt{N}.`$
This constant $`C`$ is actually a function of $`\alpha `$ and $`\sigma `$, and the fact that the degrees of the quadrature rules are linearly graded. This proves Theorem 2.3:
###### Theorem 2.3 (Convergence of the $`h`$-$`p`$ Method with Geometric Grading)
Consider the approximation of $`_0^1|x|^\alpha 𝑑x`$, $`1<\alpha <1`$, using an $`h`$-$`p`$ method based on a geometric mesh on $`m`$ intervals with parameter $`\sigma `$, and using a Gauß–Lobatto quadrature rule on $`j+1`$ points on interval $`j`$. For some constants $`C`$, and $`\rho <1`$, the error $`E_m`$ satisfies:
$`E_m<C\sigma ^{2m}.`$
Babuška et al. describe an approximation theory for the $`h`$, $`p`$ and $`h`$-$`p`$ methods for the finite element method. This material may be able to be simplified and adapted (as the theory for integration should be easier than for approximation), and also used in error analysis.
#### 2.6.3 Experimental Results for a Real Integral
The example $`_0^1\sqrt{x}𝑑x`$ is used to demonstrate the above convergence results. matlab code used (hpmeth.m and funchp.m), is contained in Appendix A, and error results are presented in Figures 1 and 2.
Figure 1 compares convergence rates for various choices of the algebraic grading parameter $`\gamma `$, whilst holding constant the number of points in the quadrature rule used on each interval; that is $`p=6`$ is fixed. Figure 2 shows similar data, but varies $`p`$. (Here, $`\gamma `$ is allowed to vary, and is chosen to be equal to $`p`$ for convenience.) Both plots are shown compared with the corresponding $`h`$-$`p`$ result, using a geometric grading, and $`\sigma =0.15`$. The linear increase in number of points used in the quadrature rule means that $`n_j=j+1`$, $`j=1:m`$, so $`N=m\left(m+1\right)/2+1`$ (as rules are closed, but ends are not), and thus $`m\sqrt{N}`$. The $`h`$-$`p`$ result demonstrates $`𝒪\left(\rho ^\sqrt{N}\right)`$ behaviour. The $`h`$-$`p`$ data does not show as a straight line, but this can be seen in Figure 3, where it is plotted versus $`\sqrt{N}`$.
#### 2.6.4 Extension to Complex Contour Integrals
The method is easily extended to complex contour integrals. Good test problems have closed contours, and integrals which can be directly evaluated using Cauchy’s integral formula. To demonstrate this, consider $`f\left(z\right)=e^{i\pi /4}z^1\sqrt{z1}`$, integrated around the unit circle. The integrand has a simple pole at $`z=0`$, and a derivative singularity at $`z=1`$. The resulting integral is:
$`{\displaystyle _\mathrm{\Gamma }}f\left(z\right)𝑑z=e^{i\pi /4}.`$
matlab code used (cint.m and funcci.m) is contained in Appendix A. Error results for an $`h`$-$`p`$ method using Gauß–Lobatto quadrature rules are presented in Figure 3. The mesh is geometrically graded, with parameter $`\sigma =0.15`$. For a segment of a closed contour, with a corner at either end, let $`D`$ be chosen as the number of mesh intervals between each corner and a wide, central interval, so $`m=2D+1`$ is the number of mesh intervals over that segment. Here, as the contour is the unit circle, $`2`$ artificial corners are placed, and $`D`$ is varied from $`8`$ to $`15`$. The grading of the quadrature rules is similar to that used in §2.6.3 – the number of points used in the quadrature rule increases linearly with the number of mesh intervals from the nearest corner, starting at $`2`$ on the interval nearest the corner, and finishing at $`D+2`$ on the central interval.
Convergence is plotted for the logarithm of the error with $`\sqrt{N}`$. Observe that the plot is linear, that is, the error is $`𝒪\left(\rho ^\sqrt{N}\right)`$. These superb results show that the method is excellent for the numerical approximation of closed complex contour integrals. This success motivates the use of the $`h`$-$`p`$ method in the CBIEM, where quadrature rules for complex contour integrals are required in the numerical approximation of the solution to an integral equation.
### 2.7 Summary – Advantages of $`h`$-$`p`$ Methods
This section has discussed three important aspects of the numerical approximation of integrals with end point singularities:
1. An $`h`$-$`p`$ quadrature method is superior to other methods.
2. A geometrically graded mesh is superior to other choices of grading (maybe only marginally better than an algebraic one).
3. The appropriate family of quadrature rules to use is the Gauß–Lobatto, as they are closed, and of maximal degree for the number of quadrature points used.
The quadrature rule used in §3 is chosen in this manner.
## 3 The Complex Boundary Integral Equation Method
### 3.1 Origins and Description
The CBIEM is a technique which numerically approximates the solution of the Dirichlet problem.<sup>9</sup><sup>9</sup>9The space containing the functions approximating the solution of the Dirichlet problem is a Sobolev space, which is a generalisation of the Banach space of continuous functions to include functions with ‘weak derivatives’. It reformulates the solution of the Dirichlet problem as the real part of a function which can be found as the solution of a complex boundary integral equation. The solution of a discretised version of this integral equation is then found using a collocation technique. Finally, a discretisation of Cauchy’s integral formula is used to approximate the solution to the original problem at interior points, based on the approximate boundary data.
It is related to the ‘Complex Variable Boundary Element Method’ , which is a Galerkin version of the same technique, using ‘hat’ functions as a basis. (The collocation method creates an approximation to the boundary data by interpolating from known data, whilst the Galerkin constructs an approximation in terms of a series of basis functions defined on segments of the boundary.) As originally stated, the CVBEM only works on polygonal domains,<sup>10</sup><sup>10</sup>10The CVBEM has also been generalised to doubly connected domains . whilst the CBIEM is more general in that it also works on non-polygonal domains.
The Dirichlet problem considered is on an open, finite, simply connected and non-empty region $`\mathrm{\Omega }^2`$. $`\mathrm{\Omega }`$ is bounded by $`\mathrm{\Gamma }`$, a piecewise continuous, anticlockwise oriented contour. $`\mathrm{\Gamma }`$ has a finite number of corners, at which its derivative is discontinuous. The Dirichlet problem is:
Given boundary data $`f`$, find $`U:\mathrm{\Omega }`$ subject to the conditions:
$`^2U\left(𝐱\right)=0,𝐱\mathrm{\Omega },U\left(𝐱\right)=f\left(𝐱\right),𝐱\mathrm{\Gamma }.`$
Thus, the problem is to find the solution to Laplace’s equation over a region, given functional data around its perimeter. This has many applications in the solution of potential problems, such as electrostatics and fluid flow. The value of $`U`$ at points interior to $`\mathrm{\Omega }`$ is determined by the boundary data being ‘diffused’ from the boundary inwards, according to the Laplacian operator. It turns out that the problem has a unique solution for all cases of $`f`$. In all but the most trivial cases, this solution is not expressible in closed form, and a numerical approximation is required. With sufficient (possibly enormous) computational effort, an approximation to any degree of accuracy can usually be obtained.
Problems with ‘corner singularities’ are of particular interest. In these problems, $`U`$ is differentiable in the interior, but $`U`$ becomes unbounded as the corner is approached. It is known that this behaviour is typical of solutions to the Dirichlet problem on domains with corners. Even if the boundary data is smooth, $`U`$ still becomes singular near the corner. Numerical methods must be able to produce good approximations to $`U`$, in spite of the corner singularities.
‘Interior’ methods, such as finite difference and finite element methods, become computationally expensive when applied to problems with corner singularities, and boundary integral methods are more appropriate. Interior methods require a finely discretised two dimensional mesh in the region of the corner, which greatly increases the size of the associated linear system. In contrast, a boundary integral method has only to discretise its mesh in one dimension, that of arc length on the boundary, and is expected to be much cheaper.
The usual boundary integral methods based on Green’s functions lead to a kernel with a logarithmic singularity, even on a smooth domain. This is tedious to program, and computationally inefficient if high order methods are used. If the CBIEM is used with singularity subtraction, the integrands are smooth and can be done simply and accurately by direct quadrature.
The problems caused by the corners and corner singularities are dealt with using $`h`$-$`p`$ quadrature methods, and would be difficult to implement with other types of integral equations .
### 3.2 Development of the CBIEM
The solution to the Dirichlet problem, $`U`$, is harmonic, as it satisfies Laplace’s equation in $`\mathrm{\Omega }`$. Identify $`𝐱^2`$ with $`z`$. Now $`U`$ can be thought of as the real component of an analytic function $`W\left(z\right)=U\left(z\right)+iV\left(z\right)`$, where $`V`$ is uniquely determined to within a constant. $`V`$ can be made unique by requiring $`V\left(\zeta _0\right)=0`$ for some $`\zeta _0\mathrm{\Gamma }`$ (see §3.2.5). For all $`z\mathrm{\Gamma }`$, $`U\left(z\right)f\left(z\right)`$ is immediately known. The CBIEM first approximates $`V\left(z\right)`$ on $`\mathrm{\Gamma }`$, and then uses Cauchy’s integral formula to approximate $`W`$, and hence $`U`$, at points within $`\mathrm{\Omega }`$.
#### 3.2.1 Cauchy’s Integral Formula
For an analytic function $`W`$ on a bounded domain $`\mathrm{\Omega }`$, Cauchy’s integral formula is:
$$_\mathrm{\Gamma }\frac{W\left(\zeta \right)}{\zeta z}𝑑\zeta =\pi iW\left(z\right)\times \{\begin{array}{ccc}0\hfill & z\mathrm{\Omega }\mathrm{\Gamma }\hfill & \mathrm{exterior}\hfill \\ 1\hfill & z\mathrm{\Gamma }\hfill & \mathrm{boundary}\hfill \\ 2\hfill & z\mathrm{\Omega }\hfill & \mathrm{interior}.\hfill \end{array}$$
(10)
When $`z`$ is on the boundary,<sup>11</sup><sup>11</sup>11This result is a simplification. If $`z`$ is at a corner, replace $`1`$ with $`\alpha /\pi `$, where $`\alpha `$ is the interior angle subtended by the corner (else $`\alpha =\pi `$). This result requires that collocation points are *not* placed at corners, to avoid unwanted complexities in the implementation. Fortunately, this is already overcome by the use of node points at the corners (see §3.2.3). the integral is a Hilbert transform. The kernel is singular, and the result must be interpreted as a Cauchy principal value integral \[2, page 39\]. The CBIEM requires approximation of the Cauchy integrals by quadrature. In §3.2.3, this is used to set up a linear system for the approximation of $`V`$ on $`\mathrm{\Gamma }`$. After this has been done, in §3.3 it is used to compute an approximation to $`W`$ (and hence $`U`$) in the interior of $`\mathrm{\Omega }`$.
#### 3.2.2 The Complex Boundary Integral Equation
To derive the integral equation underlying the CBIEM, observe that letting $`W\left(\zeta \right)1`$ for the case $`z\mathrm{\Gamma }`$ in (10) yields:
$`{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{1}{\zeta z}}𝑑\zeta =\pi i,z\mathrm{\Gamma }.`$
Multiplying this by the constant $`W\left(z\right)`$ gives:
$`W\left(z\right){\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{1}{\zeta z}}𝑑\zeta ={\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{W\left(z\right)}{\zeta z}}𝑑\zeta =\pi iW\left(z\right).`$
Equating the $`\pi iW\left(z\right)`$ with that in (10) gives:
$$_\mathrm{\Gamma }\frac{W\left(\zeta \right)W\left(z\right)}{\zeta z}𝑑\zeta =0,z\mathrm{\Gamma }.$$
(11)
(11) is called the ‘Complex Boundary Integral Equation’. This derivation is parallel to that involved in *singularity subtraction* \[6, page 184\] and . The integrand is analytic, and as $`\zeta z`$, it converges to $`W^{}\left(z\right)`$. An analytic function $`W=U+iV`$, which has as its real component the solution to the Dirichlet problem, will satisfy (11). The converse is also true – a function $`W`$ that satisfies (11) will have a real component $`U`$ that satisfies the Dirichlet problem. It is hoped that a function that satisfies a discretisation of (11) will have as its real part the solution to a discretisation of the Dirichlet problem.
#### 3.2.3 Discretisation of the CBIE
The CBIEM requires the numerical approximation of the (Cauchy) integral in (11), and this is achieved using quadrature. In particular, given the possibly singular nature of $`W`$ at corners, §2 motivates the use of geometrically graded $`h`$-$`p`$ quadrature, because it is of high order for such integrands. The approximation will be referred to as a discretisation. Nomenclature used in the following discussion is shown in Figure 4. (The distinction between *mesh* and *node* (quadrature) points is made in §2.5.3.)
Consider the contour integral of an arbitrary integrand $`g\left(\zeta \right)`$ around $`\mathrm{\Gamma }`$:
$`{\displaystyle _\mathrm{\Gamma }}g\left(\zeta \right)𝑑\zeta .`$
Parameterise $`\mathrm{\Gamma }`$ using $`\gamma :[0,1]`$, such that $`\zeta _0\gamma \left(0\right)=\gamma \left(1\right)\zeta _N`$, with argument $`t`$ increasing in an anticlockwise direction around $`\mathrm{\Gamma }`$. The contour integral is \[6, page 168\]:<sup>12</sup><sup>12</sup>12This requires that $`\gamma `$ is continuous, and that $`\gamma |_{[t_{i1},t_i]}`$ is continuously differentiable for a finite partition $`0=t_0<t_1<\mathrm{}<t_n=1`$.
$`{\displaystyle _\mathrm{\Gamma }}g\left(\zeta \right)𝑑\zeta ={\displaystyle _0^1}g\left(\gamma \left(t\right)\right){\displaystyle \frac{\gamma }{t}}\left(t\right)𝑑t.`$
Approximate the integral using an $`h`$-$`p`$ quadrature rule $`\{t_j,w_j\}`$ with $`N`$ node points, defining $`\zeta _j=\gamma \left(t_j\right)`$ and $`\dot{\gamma }_j={\displaystyle \frac{\gamma }{t}}\left(t_j\right)`$:
$`{\displaystyle _\mathrm{\Gamma }}g\left(\zeta \right)𝑑\zeta ={\displaystyle _0^1}g\left(\gamma \left(t\right)\right){\displaystyle \frac{\gamma }{t}}\left(t\right)𝑑t{\displaystyle \underset{j=1}{\overset{N}{}}}g\left(\zeta _j\right)\dot{\gamma }_jw_j.`$
This formula can be applied to the Cauchy integrals. For some fixed $`z\mathrm{\Gamma }`$, consider the integrand $`g\left(\zeta \right)=\left(W\left(\zeta \right)W\left(z\right)\right)/\left(\zeta z\right)`$. Let $`W_j=W\left(\zeta _j\right)`$, and redefine $`w_j\dot{\gamma }_jw_j`$ to absorb $`\dot{\gamma }_j`$. The Cauchy integral is:
$$_\mathrm{\Gamma }\frac{W\left(\zeta \right)W\left(z\right)}{\zeta z}𝑑\zeta \underset{j=1}{\overset{N}{}}\frac{W_jW\left(z\right)}{\zeta _jz}w_j,z\mathrm{\Gamma }.$$
(12)
Approximation of the solution to the CBIE requires approximation of the Cauchy integrals without using $`W\left(z\right)`$. Instead, $`\widehat{W}\left(z\right)W\left(z\right)`$ is constructed from values of $`W`$ at the quadrature points. (12) is discretised into a linear system of order $`N`$.
Begin by choosing a set of $`N`$ different values of $`z`$ from around the boundary. These points are called the *collocation* points. It is known from analysis in the case of uniform meshes that collocation points must not lie on node points . The natural choice is to take as the collocation points the $`N`$ midpoints (in the sense of arc length) between the $`N`$ node points.<sup>13</sup><sup>13</sup>13The choice $`\zeta _{k1/2}=\left(\zeta _{k1}+\zeta _k\right)/2`$ is explicitly *not* used, as this assumes the contour is linear between points parameterised by $`t_{k1}`$ and $`t_k`$. Let $`t_{k1/2}=\left(t_{k1}+t_k\right)/2`$, $`\zeta _{k1/2}=\gamma \left(t_{k1/2}\right)`$ and $`W_{k1/2}=W\left(\zeta _{k1/2}\right)`$, for $`k=1:N`$. Interpolation from known values of $`U`$ at the node and collocation points is used with the CBIE to approximate the $`W_{k1/2}`$.
Define two complex $`N`$-vectors of $`W`$ at the node and collocation points:
$`𝐖=\left[\begin{array}{c}W_1\\ \mathrm{}\\ W_N\end{array}\right],𝐖^{}=\left[\begin{array}{c}W_{1/2}\\ \mathrm{}\\ W_{N1/2}\end{array}\right].`$
Also define the real $`N`$-vectors $`𝐔=\mathrm{}\left(𝐖\right)`$, $`𝐔^{}=\mathrm{}\left(𝐖^{}\right)`$ and $`𝐕=\mathrm{}\left(𝐖\right)`$. From (12), an order $`N`$ linear system for the components of $`𝐖`$ and $`𝐖^{}`$ is determined:
$$\underset{j=1}{\overset{N}{}}\frac{W_jW_{k1/2}}{\zeta _j\zeta _{k1/2}}w_j=0,k=1:N.$$
(14)
In summary, discretisation of Cauchy’s integral formula leads to a linear system, the solution to which is an approximation to $`W`$ at the $`N`$ node points on $`\mathrm{\Gamma }`$. This approximation can be used to approximate $`W`$, and hence $`U`$, at points within $`\mathrm{\Omega }`$.
#### 3.2.4 Linear Interpolation of $`W`$ at the Collocation Points
If the $`W_j`$ were known, (14) could be directly used to interpolate the $`W_{k1/2}`$. However, although $`U_j`$ and $`U_{k1/2}`$ are known explicitly, $`V_j`$ and $`V_{k1/2}`$ are not. The CBIEM implicitly approximates the $`V_{k1/2}`$ by interpolation from the as yet undetermined $`V_j`$ at points near $`\zeta _{k1/2}`$. That is, if a rule on $`O`$ points is being used, choose $`O`$ terms from the sequence:
$`\mathrm{},\zeta _{k3},\zeta _{k2},\zeta _{k1};\zeta _k,\zeta _{k+1},\zeta _{k+2},\mathrm{}.`$
The discretisation of the CBIE in (14), coupled with the $`2N`$ knowns $`U_j`$ and $`U_{k1/2}`$, and the interpolation for the $`V_{k1/2}`$, allows the approximation of the $`N`$ unknowns $`V_j`$. (The $`V_{k1/2}`$ are *not* explicitly required to be calculated.)
The simplest interpolation for the $`V_{k1/2}`$ is a linear one, between $`\zeta _{k1}`$ and $`\zeta _k`$, that is use:
$$V_{k1/2}\left(V_{k1}+V_k\right)/2.$$
(15)
More sophisticated interpolations to the $`V_{k1/2}`$ could involve higher degree polynomials, splines, or trigonometric polynomials. Initially, the linear choice will be used to illustrate the process. In §3.5.1, the method is extended to higher degree polynomials. For a particular problem ($`\mathrm{\Omega },f`$), there is an optimal choice of degree for the interpolation, as errors incurred by interpolation increase with the degree, and eventually become of greater magnitude than those due to discretisation.
The $`W_j`$ are found by solving for their imaginary parts $`V_j`$, and combining these with the knowns $`U_j`$. The linear system is set up as follows. Substituting $`W_j=U_j+iV_j`$ and $`W_{k1/2}=U_{k1/2}+iV_{k1/2}`$ into (14), using (15), and collecting knowns and unknowns:
$$\underset{j=1}{\overset{N}{}}\frac{\frac{1}{2}V_{k1}+\frac{1}{2}V_kV_j}{\zeta _j\zeta _{k1/2}}w_j=i\underset{j=1}{\overset{N}{}}\frac{U_{k1/2}U_j}{\zeta _j\zeta _{k1/2}}w_j,k=1:N.$$
(16)
This is a system of $`N`$ equations for the $`N`$ unknowns $`V_j`$, with RHS determined by the knowns $`U_j`$ and $`U_{k1/2}`$ (and of course the associated $`\zeta _j`$ and $`\zeta _{k1/2}`$).
#### 3.2.5 Solution of the Collocation Equations
In order to write (16) as a linear system, consider the LHS of its $`k`$th equation:
$$\frac{1}{2}V_{k1}\underset{j=1}{\overset{N}{}}\frac{w_j}{\zeta _j\zeta _{k1/2}}+\frac{1}{2}V_k\underset{j=1}{\overset{N}{}}\frac{w_j}{\zeta _j\zeta _{k1/2}}\underset{j=1}{\overset{N}{}}\frac{V_jw_j}{\zeta _j\zeta _{k1/2}}.$$
(17)
Define a matrix $`A^{N\times N}`$:
$`A=\left[\begin{array}{cccc}{\displaystyle \frac{w_1}{\zeta _1\zeta _{1/2}}}& {\displaystyle \frac{w_2}{\zeta _2\zeta _{1/2}}}& \mathrm{}& {\displaystyle \frac{w_N}{\zeta _N\zeta _{1/2}}}\\ {\displaystyle \frac{w_1}{\zeta _1\zeta _{3/2}}}& {\displaystyle \frac{w_2}{\zeta _2\zeta _{3/2}}}& \mathrm{}& {\displaystyle \frac{w_N}{\zeta _N\zeta _{3/2}}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ {\displaystyle \frac{w_1}{\zeta _1\zeta _{N1/2}}}& {\displaystyle \frac{w_2}{\zeta _2\zeta _{N1/2}}}& \mathrm{}& {\displaystyle \frac{w_N}{\zeta _N\zeta _{N1/2}}}\end{array}\right].`$
Also define a set of $`N`$ scalars $`H_k`$, for $`k=1:N`$ (the row sums of $`A`$):
$`H_k={\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{w_j}{\zeta _j\zeta _{k1/2}}}.`$
The first two terms of (17) are then $`\frac{1}{2}V_{k1}H_k+\frac{1}{2}V_kH_k`$. As $`\mathrm{\Gamma }`$ is closed, $`\zeta _0\zeta _N`$ and hence $`V_0=V_N`$. Define $`B^{N\times N}`$:
$$B=\frac{1}{2}\left[\begin{array}{cccccc}H_1& & & & H_1& \\ H_2& H_2& & & & \\ & \mathrm{}& \mathrm{}& & & \\ & & H_{N1}& H_{N1}& & \\ & & & H_N& H_N& \end{array}\right].$$
(19)
Let $`C=BA`$, then the LHS of (16) is $`C𝐕`$. Let $`\mathrm{𝟏}_N^N`$ represent the real column vector with all components unity, and the operation $`\mathrm{diag}\left(𝐱\right)`$ on $`N`$-vector $`𝐱`$ create the diagonal matrix of order $`N`$ with the diagonal entries being the respective components of $`𝐱`$. Defining $`𝐝=\mathrm{diag}\left(A\mathrm{𝟏}_N\right)𝐔^{}A𝐔`$, the system for $`𝐕`$ is:
$$C𝐕=i𝐝.$$
(20)
Attempting to directly solve the complex linear system in (20) fails, as $`𝐕`$ is overdetermined in two separate ways. Firstly, $`𝐕`$ is purely real, that is $`\mathrm{}\left(𝐕\right)=\mathrm{𝟎}`$. Partitioning (20) into real and imaginary components yields two purely real linear systems, either one of which can be solved for what should be the same solution $`\widehat{𝐕}`$. The system to be solved is: $`\mathrm{}\left(C\right)𝐕=\mathrm{}\left(𝐝\right)`$ or $`\mathrm{}\left(C\right)𝐕=\mathrm{}\left(𝐝\right)`$. By redefining $`C\mathrm{}\left(C\right)`$ and $`𝐝\mathrm{}\left(𝐝\right)`$, the first choice for the solution of $`𝐕`$ is made. The linear system is thus:
$$C𝐕=𝐝.$$
(21)
The second way that (20) is overdetermined is that $`𝐕`$ is known only to within a constant.<sup>14</sup><sup>14</sup>14A scalar multiple of $`\mathrm{𝟏}_N`$. So, if direct solution of (21) is attempted, singularity problems will occur, and the resultant $`\widehat{𝐕}`$ will be infinite.<sup>15</sup><sup>15</sup>15Well, a numerical approximation to $`\mathrm{}`$! To accommodate this, arbitrarily<sup>16</sup><sup>16</sup>16For test problems, $`\widehat{V}_N=V\left(\zeta _N\right)`$ is actually used. set $`V_N=0`$, and compute the rest of the components of $`𝐕`$ by subtracting rows in (21). The result is a fully determined order $`N1`$ linear system. The $`N`$ rows of (21) are:
$`\left(C𝐕\right)_j=𝐝_j,j=1:N.`$
Subtracting rows gives a system of $`N1`$ equations:
$`\left(C^{}𝐕\right)_{j1}=\left(C𝐕\right)_j\left(C𝐕\right)_{j1}=𝐝_j𝐝_{j1}=𝐝_{j1}^{},j=2:N.`$
Defining $`𝐕^{}^{N1}`$ as the first $`N1`$ entries of $`𝐕`$ (where the last entry is zero):
$$C^{}𝐕^{}=𝐝^{}.$$
(22)
Here:
$`C^{}`$ $`=`$ $`\left[\begin{array}{ccc}C_{2,1}C_{1,1}& \mathrm{}& C_{2,N1}C_{1,N1}\\ C_{3,1}C_{2,1}& \mathrm{}& C_{3,N1}C_{2,N1}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ C_{N,1}C_{N1,1}& \mathrm{}& C_{N,N1}C_{N1,N1}\end{array}\right]^{(N1)\times (N1)}`$
$`𝐝^{}`$ $`=`$ $`\left[\begin{array}{c}d_2d_1\\ d_3d_2\\ \mathrm{}\\ d_Nd_{N1}\end{array}\right]^{N1}.`$
Solution of the order $`N1`$ linear system in (22) yields the approximation to $`𝐕`$, and hence $`𝐖`$ ($`W`$ at the $`N`$ node points).
### 3.3 Approximation of $`U`$ at Interior Points
The approximation $`𝐖=𝐔+i𝐕`$ on $`\mathrm{\Gamma }`$ is used to approximate $`W`$ (and hence $`U`$) at interior points of $`\mathrm{\Omega }`$, using Cauchy’s integral formula for points within $`\mathrm{\Omega }`$:
$`W\left(z\right)={\displaystyle \frac{1}{2\pi i}}{\displaystyle _\mathrm{\Gamma }}{\displaystyle \frac{W\left(\zeta \right)}{\zeta z}}𝑑\zeta ,z\mathrm{\Omega }.`$
Discretising this gives the approximation, for $`z\mathrm{\Omega }`$:
$$W\left(z\right)\frac{1}{2\pi i}\underset{j=1}{\overset{N}{}}\frac{W\left(\zeta _j\right)}{\zeta _jz}w_j.$$
(25)
Numerical problems occur using this simple approximation, as points $`z\mathrm{\Omega }`$ near the boundary (where $`z\zeta _j`$ is small, for some $`\zeta _j`$), generate very large terms in the sum. In fact, the integrand is nearly singular, so only poor accuracy is expected. Instead, the technique of singularity subtraction (referenced in §3.2.2) uses the result from (12):
$`{\displaystyle \frac{1}{2\pi i}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{W\left(\zeta _j\right)W\left(z\right)}{\zeta _jz}}w_j=0.`$
The integrand is now smooth, and good results can be expected from quadrature. This yields $`W\left(z\right)`$ as a ‘corrected’ (25):
$`W\left(z\right)=\left[{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{W\left(\zeta _j\right)}{\zeta _jz}}w_j\right]/\left[{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{1}{\zeta _jz}}w_j\right].`$
Implementation of the CBIEM using this result is successful. The code supplied (see §4) does not go beyond the stage of the computation of $`\widehat{𝐕}`$ on $`\mathrm{\Gamma }`$, as it is known that the approximation to $`U`$ in the interior is actually *more* accurate than the approximations to $`V`$ on the boundary . Experiment demonstrates this, and thus computation of $`U`$ at interior points need not be further described.
### 3.4 Performance of the CBIEM on Model Problems
Application of the CBIEM yields different quality results depending on the continuity of the model problem and the contour. In the simplest case, both are smooth, there is no singularity, and standard quadrature gives good results. In fact, because of periodicity, even the trapezoidal rule on a uniform mesh gives very good results. There is no need to use $`h`$-$`p`$ quadrature, but it will still work well.
Now consider the case where $`\mathrm{\Gamma }`$ is smooth, but $`W`$ has singularities. For example, if $`W\left(\gamma \left(s\right)\right)s^{1/2}`$, the integrand of the CBIE (11) can have behaviour $`s^{1/2}`$ near the corner. However, using singularity subtraction and $`h`$-$`p`$ quadrature, experiment demonstrates that both the discretisation error and the final error in $`\widehat{𝐕}`$ on the boundary are superb. (See also the example of complex contour integration in §2.6.4.)
If $`W`$ is smooth, but $`\mathrm{\Gamma }`$ is not (has corners), in general, the errors will be expected to increase with the sharpness of the corner. The most difficult cases are cusps or reentrant corners (e.g. the corner in a cardioid). Even for a model problem with a smooth solution $`U`$, a corner singularity in $`V`$ (and hence $`W`$) will occur.
Let $`r`$ represent radial distance from a corner. It is known \[23, pages 257–259\], that at a corner with interior angle $`\left(1\chi \right)\pi `$, a singularity of the form $`r^{1/\left(1\chi \right)}`$ will be found. At worst, for a reentrant corner, $`\chi =1`$, so the form is $`r^{1/2}`$. A good model problem is thus a contour with a corner where the true solution has local behaviour $`Ur^{1/2}`$. This is obtained, for example, using $`W\left(z\right)=z^{1/2}`$. If a uniform mesh were used, the greatest component of the error in $`V`$ will come from the intervals adjacent to the corner. The use of a geometrically graded mesh reduces this component to a level comparable with that of other mesh intervals. Errors will not be of the very high order that is expected for smooth contours, but should still be acceptable.
### 3.5 Higher Degree Interpolatory Polynomials
#### 3.5.1 Introduction
To extend the technique described in §3.2.4, the linear interpolation in (15) is replaced by a higher degree interpolation. The net result of this is to change the definition of the matrix $`B`$ in (19). Other possible techniques of improving the accuracy of the interpolation, such as splines, are not considered here, as they are difficult to implement. The principle involved is that increasing the order of the interpolatory polynomial should reduce the discretisation error, which is expected to be greatest on the largest intervals.
The ‘nearest’ $`O`$ node points on either side of $`\zeta _{k1/2}`$ are used to yield an interpolatory polynomial of degree $`O1`$. In general, for points far from the nearest corner, this means to take the first $`O`$ terms of the sequence $`\zeta _k,\zeta _{k1},\zeta _{k+1},\zeta _{k2},\mathrm{}`$. Otherwise, the term ‘nearest’ is used loosely, as interpolation cannot continue around a corner. Where there are less than $`O/2`$ node points between the collocation point and the nearest corner (including the node on the corner); instead $`O`$ points from and including the corner are used. This results in an interpolatory polynomial that is expected to be least accurate at the collocation point adjacent to the corner. (A possible improvement in this schema is to organise the interpolation rules such that their order increases say, linearly, with node index away from the corner.) The constraint on $`O`$ due to the mesh parameter $`D`$ is described in §3.5.4.
§§3.5.2 and 3.5.3 deal with technical implementation issues, and can be skipped without loss of continuity.
#### 3.5.2 Lagrange Form of the Interpolatory Polynomial
The Lagrange form of the interpolatory polynomial is used. Given a set of values for the $`V_j`$, at positions $`\zeta _j=\gamma \left(t_j\right)`$, the approximation to $`V`$ at point $`\zeta _{k1/2}`$ is:
$`V\left(\zeta _{k1/2}\right)\widehat{V}_{k1/2}={\displaystyle \underset{jF_k}{}}\lambda _j\left(t_{k1/2}\right)V_j.`$
Here:<sup>17</sup><sup>17</sup>17Warning: Replacing $`t_\nu `$ with $`\zeta _\nu `$ and $`t_j`$ with $`\zeta _j`$ in this formula *cannot* be done, as the contour segments are *not* necessarily straight.
$`\lambda _j\left(t\right)={\displaystyle \underset{\nu F_k,\nu j}{}}{\displaystyle \frac{tt_\nu }{t_jt_\nu }}.`$
$`F_k`$ is a set of indices of the nearest node points, specific to the collocation point $`\zeta _{k1/2}`$. Specifically, where a degree $`O1`$ interpolatory polynomial is used on $`O`$ points, and $`C`$ is the index of the nearest corner to $`\zeta _{k1/2}`$:
$`F_k=\{\begin{array}{cc}\{kO/2,\mathrm{},k1,k,\mathrm{},k+O/21\}\hfill & \mathrm{in}\mathrm{general}\hfill \\ \{C,\mathrm{},C+O1\}\hfill & \mathrm{if}kO/2<C\hfill \\ \{CO+1,\mathrm{},C\}\hfill & \mathrm{if}k+O/2>C.\hfill \end{array}`$
Let $`F_k`$ be the $`k`$th row of a table $`F`$, and let the above $`\lambda _j`$ be the $`j`$th element of the $`k`$th row of another table, $`L`$, of the associated weights. The notations $`F(k,j)`$ and $`L(k,j)`$ are used to describe the set of $`O`$ nodal indices and weights associated with the interpolation at point $`\zeta _{k1/2}`$, where $`k=1:N`$ and $`j=1:O`$. The structure of $`F`$ becomes more complicated with increasing $`O`$ and with increasing number of corners. Details of its construction are not provided here, but can be read from the program cbiem.m in Appendix A.1.
For some integer $`D`$, on each segment of $`\mathrm{\Gamma }`$ (with a corner at each end), a mesh is constructed that has $`D1`$ internal points between each corner and a (wide) interval which spans the centre of the segment. This results in a total of $`2D`$ mesh points, and $`2D1`$ mesh intervals (see Figure 4).
For example, if $`D=3`$, then there are $`4`$ interior and $`2`$ corner mesh points on each segment, together with $`4`$ extra interior points, for a total of $`S=\left(D+1\right)^2+1=10`$. Further, if there are $`NC=3`$ corners, then the linear system has order $`N=NC(S1)=27`$. If $`O=6`$ (quintic interpolation about the ‘nearest’ $`6`$ points), then $`F^{N\times O}`$. Where divisions in the structure of $`F`$ due to the corners are reflected by partitions, $`F`$ is:<sup>18</sup><sup>18</sup>18Liberal use of matlab notation is made, and there is a confusion between computer array element and mathematical matrix entry notation: $`F(j,k)F_{j,k}`$.
$`\text{ }F=\text{ }\left[\begin{array}{cccccc}27& 1& 2& 3& 4& 5\\ 27& 1& 2& 3& 4& 5\\ 27& 1& 2& 3& 4& 5\\ 1& 2& 3& 4& 5& 6\\ 2& 3& 4& 5& 6& 7\\ 3& 4& 5& 6& 7& 8\\ 4& 5& 6& 7& 8& 9\\ 4& 5& 6& 7& 8& 9\\ 4& 5& 6& 7& 8& 9\\ & & & & & \\ 9& 10& 11& 12& 13& 14\\ 9& 10& 11& 12& 13& 14\\ 9& 10& 11& 12& 13& 14\\ 10& 11& 12& 13& 14& 15\\ 11& 12& 13& 14& 15& 16\\ 12& 13& 14& 15& 16& 17\\ 13& 14& 15& 16& 17& 18\\ 13& 14& 15& 16& 17& 18\\ 13& 14& 15& 16& 17& 18\\ & & & & & \\ 18& 19& 20& 21& 22& 23\\ 18& 19& 20& 21& 22& 23\\ 18& 19& 20& 21& 22& 23\\ 19& 20& 21& 22& 23& 24\\ 20& 21& 22& 23& 24& 25\\ 21& 22& 23& 24& 25& 26\\ 22& 23& 24& 25& 26& 27\\ 22& 23& 24& 25& 26& 27\\ 22& 23& 24& 25& 26& 27\end{array}\right]\text{ }=\text{ }\left[\begin{array}{c}Ft(1)\\ Ft(2)+S\text{ones}(Ft)\\ \mathrm{}\\ Ft(NC)+(NC1)S\text{ones}(Ft)\end{array}\right].`$
For $`j=1:S1`$ and $`k=1:O`$, in general $`Ft^{\left(S1\right)\times O}`$ is:
$`Ft(j,k)=\left[\begin{array}{cc}k1\hfill & jO/2\hfill \\ k1+jO/2\hfill & O/2<j<SO/2\hfill \\ k1+SO\hfill & SO/2j\hfill \end{array}\right].`$
(Exception: $`F(1:O/2,1)N\times \mathrm{𝟏}_{O/2}`$.)
#### 3.5.3 Construction of $`B`$
The matrix $`B`$ is required in the construction of the linear system in (20), and is the only thing that changes when $`O`$ is varied. For the case of linear interpolatory polynomials ($`O=2`$), a formula involving the terms $`H_k`$ is used to approximate the value of $`V`$ at the collocation points. For larger $`O`$, this is replaced with a considerably more sophisticated formula. As $`O`$ increases, the bandwidth of $`B`$ increases. Naturally, this new $`B`$ simplifies to the earlier definition if $`O=2`$ is used, but is obtained at greater computational expense.
The crucial change is in the approximation to the value of $`V_{k1/2}`$, which in (16) is the term $`\frac{1}{2}V_{k1}+\frac{1}{2}V_k`$. This is now replaced with:
$`{\displaystyle \underset{jF_k}{}}\lambda _j\left(\zeta _{k1/2}\right)V_j={\displaystyle \underset{j=1}{\overset{O}{}}}L(k,j)V_{F(k,j)}.`$
The first $`O`$ terms of the LHS of the $`k`$th line of (16) (there is only one other term) are now:
$`{\displaystyle \underset{j=1}{\overset{O}{}}}L(k,j)V_{F(k,j)}H_k.`$
Construction of the real order $`N`$ matrix with $`(k,j)`$th entry $`L(k,j)V_{F(k,j)}`$ is required. Multiplying each row by the corresponding $`H_k`$ converts this to $`B`$.
Details of the construction of $`B`$ are not provided here, but the illustrative example used in §3.5.2 is continued. Recall that $`D=3`$, $`NC=3`$ and $`O=6`$, so $`S=10`$, and $`N=27`$. First construct $`Bt`$ (a ‘skewed’ version of $`L`$), and then calculate $`B`$ by multiplying $`Bt`$ through by the $`H_k`$. $`Bt`$ is constructed using a ‘shift vector’ $`G`$, where $`G_k`$ is equal to the number of zeros to be put in front of row $`k`$ of $`L`$ to make it row $`k`$ of $`Bt`$. This $`G`$ has a structure formed from a temporary $`Gt`$:
$`Gt`$ $`=`$ $`\left[\mathrm{zeros}(1,O/2)[1:SO1](SO)\mathrm{𝟏}_{O/2}^{}\right]^{}`$
$`G`$ $`=`$ $`\left[GtGt+S1\mathrm{}Gt+NC\left(S1\right)\right]^{}`$
$`Bt(k,:)`$ $`=`$ $`\left[\mathrm{zeros}(1,G\left(k\right))L(k,:)\mathrm{zeros}(1,NG\left(k\right)O/2)\right]^{}k=1:N`$
$`B`$ $``$ $`Bt(:,2:N)B(1:O/2,N)Bt(1:O/2,1).`$
The overall structure of $`B`$ is depicted in Figure 5, where $`\times `$ and $``$ represent nonzero and zero entries respectively.
The $`NC\times NC`$ submatrices within the structure are each of order $`S1`$. The $`k`$th row generally consists of $`O=6`$ contiguous nonzero entries, starting at column $`G_k`$. These entries are the $`k`$th row of $`L`$, multiplied by $`H_k`$. That is, the string of $`O=6`$ nonzero elements in row $`k`$ represents:
$`\left[H_kL_{k,1}H_kL_{k,2}\mathrm{}H_kL_{k,O1}H_kL_{k,O}\right].`$
Additionally, in the first $`O/2`$ rows, the first of these $`O`$ entries is shifted to the $`N`$th column, due to a ‘wraparound’ effect. In view of the banded structure of $`B`$, it appears that a method designed to exploit this structure would be appropriate in the solution of (22). Unfortunately $`A`$ is neither banded, nor of a particularly simple structure, so this approach is nontrivial, and is a direction for further work.
#### 3.5.4 Limitations on the Degree of the Interpolatory Polynomial
The discussion and results presented in §2 prompt the use of a geometrically graded mesh, with the number of points in the (closed) quadrature rule on each mesh interval linearly increasing with interval number from the corner, beginning with $`2`$ adjacent to the corner, and becoming $`D`$ on the central (widest) interval. The use of a linear grading is found in the literature of the finite element method and the usual boundary element method . Changing from linear to quadratic or higher degree may reduce the errors, but its implementation is beyond the scope of this report.
Consider a segment of $`\mathrm{\Gamma }`$ divided into a mesh on $`2\left(D+1\right)`$ points, including corners, with $`m=2D+1`$ intervals. Basic quadrature rules on $`n_j=2,3,\mathrm{},D1,D,D1,\mathrm{},3,2`$ points, are used over intervals $`j=1:m`$. After all of the common end points are considered, the total number of points in the final composite quadrature rule for that segment is $`S=\left(D+1\right)^2+1`$. Slightly different results would apply if the quadrature rules were open. Recall that this choice is rejected, as it leads to node points that avoid the singularity.
$`S`$ imposes a limit on $`O`$, the number of points used in the interpolation rule. As $`O`$ is usually small (for reasons of computational efficiency, typically $`O6`$), this limitation is usually not significant. For example, if $`D=2`$, $`O8`$, and if $`D=3`$, $`O16`$.
## 4 Implementation and Results
### 4.1 Implementation
The CBIEM is implemented as a set of functions in matlab<sup>19</sup><sup>19</sup>19matlab is an (interpreted) matrix computation package, and is a trademark of The Mathworks, Inc. code, presented in Appendix A. (These functions appear in alphabetic order, interspersed with several functions referenced in §2.) The main routine is cbiem.m. The parameterisation of the contour and its derivative are computed within cbiem.m, and it calls an auxiliary function (funccb.m) to compute the true solution for a test problem. The quadrature points for basic rules are obtained by calling the function gettw.m, which provides either (closed) Newton–Cotes points (using a routine internal to gettw.m), or calls another function, lobatto.m, which computes the points for a closed Gauß–Lobatto rule. Two further functions are used to create an $`h`$-$`p`$ composite quadrature rule out of a set of basic quadrature rules (hprmesh.m), and compose a quadrature rule over a contour with several corners (rmesh.m). For testing cbiem.m over a large set of parameters (e.g. generating the data for Tables 2 to 10), a driver function, testcb.m, is used.
Within cbiem.m is a description of its input parameters. For test problems, where the true solution is known, it plots and calculates norms of $`𝐕\widehat{𝐕}`$, and also calculates some discretisation errors. Explicit computation of the approximate solution within $`\mathrm{\Omega }`$ is *not* performed. Experiments with doing this demonstrate that the error results obtained are of the same order as those returned.
### 4.2 Experimental Results using a Teardrop Contour
#### 4.2.1 Description
Although the CBIEM code is generalised to the situation of multiple corners, good experimental contours have only one corner, to facilitate isolation of the sources of error. This section describes numerical results for the CBIEM, using a teardrop contour,<sup>20</sup><sup>20</sup>20Another important test contour is a cardioid with a reentrant corner. depicted in Figure 6. The contour is parametrically given by:
$`\gamma \left(t\right)=2\mathrm{sin}\left(\pi t\right)+i\mathrm{sin}\left(2\pi t\right)t[0,1].`$
This is the same contour as that used in . It has a right angle corner at the origin, which facilitates the use of test problems $`W\left(z\right)=z^\alpha `$. For $`\alpha (0,1)`$, there is a discontinuity in the derivative of the true solution at the origin, which becomes more pathological as $`\alpha 0`$.
Error results are presented using an unweighted vector $`2`$ norm:
$`𝐕\widehat{𝐕}_2=\left[_{i=1}^N\left(V_i\widehat{V}_i\right)^2\right]^{1/2}.`$
An appropriately weighted discretisation of the $`L_2`$ norm might seem more appropriate, but would effectively only present the norm over the central interval, as the widths of the end point intervals are very small. The use of an infinity norm is also appealing, but the $`2`$ norm allows the user to experiment with interpolation formula gradings, to independently reduce the error over different regions of the contour (see §4.2.6). Also, experimental data shows the behaviour of the infinity norm is very similar to that of the $`2`$ norm.
Tables 2 to 10 present error results for the teardrop contour, for three different model problems: $`W\left(z\right)=z^2`$, $`z^{1/2}`$ and $`z^{1/4}`$; various choices of two mesh grading parameters ($`\sigma `$ and $`D`$); and choices of $`O`$, the number of points used by the interpolatory polynomial in the collocation process. In each table $`N=\left(D+1\right)^2`$ is the size of the linear system being solved, such that there are $`2D+1`$ mesh intervals between one corner and the next (see §3.5.2). The nine tables cover three illustrative choices of the mesh parameter $`\sigma `$ for each of the three model problems. In each case, the results presented are for a choice of $`\sigma `$ close to the optimal $`\sigma `$, and two nearby values of $`\sigma `$ that demonstrate the increase in the error in each direction. Results have been selected from a much larger data set. Within each table, the minimum error result is emboldened.
#### 4.2.2 Observations of $`z^2`$
This test function does *not* have a singularity, and the results are good. Despite the corner, the error reduces with increasing either $`D`$ or $`O`$, until a point is reached where roundoff error, caused by excessive order in the interpolatory polynomial, begins to encroach.
#### 4.2.3 Observations of $`z^{1/2}`$
This case has a corner singularity, and represents the ‘worst’ that singularities get in practice (that is, for $`z^\alpha `$ singularities, in practice $`\alpha 1/2`$). Although the error decreases with increasing $`D`$ or $`O`$, it does so more slowly than for $`z^2`$, and is orders of magnitude larger. As for $`z^2`$, there comes a point where increasing $`O`$ causes the error to increase, and indeed grow exponentially. The results for $`𝐕\widehat{𝐕}`$ for the case $`\sigma =0.10`$, $`D=9`$, $`O=6`$ are plotted in Figure 7. The abscissae are plotted uniformly, for if they were plotted versus parameter $`t`$, the geometric grading would bunch up most of the results at the ends (corner of teardrop). Observe that these results are, as expected, antisymmetric.
#### 4.2.4 Observations of $`z^{1/4}`$
A $`z^{1/4}`$ singularity is beyond the range of singularities expected for smooth test functions. The errors are worse again than for $`z^{1/2}`$, and they do not decrease as fast with increasing $`D`$ or $`O`$. In fact, when the test problem is this pathological, the Dirichlet problem is fast becoming a boundary layer problem, which should be dealt with using more specialised methods.
#### 4.2.5 The Black Art of Choosing $`\sigma `$
The minimum error results obtained for each test problem are plotted versus $`\sigma `$ in Figure 8, and demonstrate that there is an optimal choice of $`\sigma `$, which varies significantly with the test problem. The use of the CBIEM in applications, where the true solution is not known in advance, could falter on the setting of $`\sigma `$. If the computational cost is to be minimised, then it is important to find the optimal $`\sigma `$, however, it may be expensive to try many $`\sigma `$ until the optimal one is found. The literature does not justify a choice of $`\sigma `$, but merely states it, e.g. uses $`\sigma =0.15`$ for a particular (finite element) application. The optimal choice of $`\sigma `$ for the paradigm test problem $`z^{1/2}`$ is $`\sigma =0.10`$. As $`z^{1/2}`$ is the worst singularity expected in practice (see §3.4), this should be a good guide as a starting guess for any problem with an unknown solution.
#### 4.2.6 Improvements in the Technique
Consider Table 6, where the best error result is obtained using $`O=6`$. The error may be able to be reduced by grading the order of the interpolatory polynomial over the mesh intervals. Near the corner, the use of high order interpolation may actually *increase* the component of the error, although this may be appropriate far away from the corner. It may be ideal to grade the order of the interpolatory polynomial from $`O=2`$ near the corner, to $`O=6`$ (or greater) farthest from the corner.
Direct implementation of this result requires extensive modification to the matrix $`B`$ used by the CBIEM<sup>21</sup><sup>21</sup>21These comments also refer to the associated $`F`$ and $`L`$ matrices. (see §3.5.1), and is beyond the scope of this report. Another way of achieving the same effect is to calculate $`B`$ matrices $`B_2,B_4,\mathrm{},B_{16}`$ for $`O=2,4,\mathrm{},16`$, then construct a new $`B`$ from appropriate rows of them, and insert this new $`B`$ at the relevant point in the CBIEM.
However, intuition is misleading here. The minimum error in Table 6 is $`2.4158\times 10^5`$, using $`\sigma =0.10`$ and $`D=9`$. This corresponds to $`O=6`$ on each mesh interval. Many experiments in variation of the order of the interpolation rule, holding fixed $`\sigma `$ and $`D`$, find that the very best error result that can be obtained is $`1.7444\times 10^5`$ (a $`28\%`$ reduction), using a grading with interpolation rules of $`O=12,2,2,6,6,6,10,10,10`$ over the $`9`$ mesh intervals from the corner to the centre. Surprisingly, the component of the error over the first interval *decreases* with increasing $`O`$. This is depicted in Figure 9, which shows that the error is uniformly distributed around the contour, except for the largest component, at the corner.
It appears that what is happening is that the method has come up against a discretisation error barrier. For this problem, the discretisation error does not decrease particularly quickly, and is a maximum at the corner.
## 5 Further Directions for Research
This section enumerates various possibilities for future work on the CBIEM.
1. There is the potential for error reduction using graded interpolation rules (see §4.2.6). Similarly, other choices for the grading of the quadrature rules may assist in error reduction, e.g. quadratic increase in degree of quadrature rule with node number from the corner, rather than linear as is presently used.
2. A proper examination of the computational efficiency of the CBIEM is required. This would involve setting up, say, a finite difference solution for the Dirichlet problem, and comparing flop counts required to obtain comparable accuracies.
3. Analysis of the choice of optimal $`\sigma `$ is desirable. Currently, the method is hampered by this not being known in advance.
4. Application of the technique to conformal mapping may be worthwhile.
5. It would be computationally efficient if solution of the linear system involving the matrix $`C=BA`$ could exploit the banded structure of $`B`$ (see §3.5.3).
6. An alternative collocation technique is possible . Given nodes $`\zeta _j`$ with weights $`w_j`$, and collocation points $`\zeta _{j+1/2}`$ with weights $`w_{j+1/2}`$:
$`{\displaystyle _\mathrm{\Gamma }}F\left(\zeta \right)𝑑\zeta {\displaystyle \underset{j=1}{\overset{N}{}}}F\left(\zeta _j\right)w_j{\displaystyle \underset{k=1}{\overset{N}{}}}F\left(\zeta _{k1/2}\right)w_{k1/2}.`$
Approximate the unknowns $`V_j=V\left(\zeta _j\right)`$ by collocating at $`\zeta _{k1/2}`$, and the unknowns $`V_{k1/2}=V\left(\zeta _{k1/2}\right)`$ by collocating at $`\zeta _j`$. This gives an order $`2N`$ system for the $`2N`$ unknowns $`V_j`$ and $`V_{k1/2}`$, but avoids interpolation.
$`0`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{W_jW_{k1/2}}{\zeta _j\zeta _{k1/2}}}w_j`$
$`0`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{N}{}}}{\displaystyle \frac{W_{k1/2}W_j}{\zeta _{k1/2}\zeta _j}}w_{k1/2}.`$
The technique appears to be computationally wasteful, but may be worth investigating, as it would be simpler to implement.
7. The CVBEM was developed to solve $`2`$D fluid flow problems ,<sup>22</sup><sup>22</sup>22Other references to the CVBEM include . where components of the complex potential (the fluid potential $`\mathrm{\Phi }`$ or the streamline function $`\mathrm{\Psi }`$) are known at different points around the contour, typically from physical measurements.<sup>23</sup><sup>23</sup>23Warning: the notation used here $`W=U+iV`$ is equivalent to the fluid flow notation $`W=\mathrm{\Phi }+i\mathrm{\Psi }`$, so that $`U`$ and $`V`$ here have a different meaning from the fluid flow case, where they are commonly the components of the velocity $`q=U\widehat{𝐢}+V\widehat{𝐣}`$, and $`U={\displaystyle \frac{\mathrm{\Phi }}{x}}={\displaystyle \frac{\mathrm{\Psi }}{y}}`$, $`V={\displaystyle \frac{\mathrm{\Phi }}{y}}={\displaystyle \frac{\mathrm{\Phi }}{x}}`$. A modification of the CBIEM can convert it to become a solver for Neumann (and thence mixed) boundary value problems. In the Neumann boundary value problem, $`U_\nu `$, the derivative of $`U`$ across $`\mathrm{\Gamma }`$, is known instead of $`U`$. Use the Cauchy–Riemann equations to observe that $`U_\nu =\pm V_\tau `$ (the tangential component of $`V`$). The boundary information can be used to construct an approximation to $`V`$, by integration of $`V_\tau `$ around $`\mathrm{\Gamma }`$, using a suitable zero point (adding in a constant):
$`V\left(\gamma \left(t\right)\right)={\displaystyle _0^t}V_\tau \left(\gamma \left(t^{}\right)\right)𝑑t^{}.`$
The same collocation process previously used to approximate $`V`$ can in this case be used to approximate $`U`$.
Beyond this, the technique is particularly applicable to free boundary problems , and may be able to be generalised to other elliptic (and possibly other second order) operators.
8. The method would easily parallelise. The establishment of the linear system is computationally expensive, more so for high order interpolatory polynomials. This, as well as solution of the linear system, would efficiently (geometrically) parallelise.
## Appendix A Listing of matlab “.m” files
### A.1 cbiem.m
```
function Vnnorm = cbiem(CCase, D, sigma, O, alpha);
%function Vnnorm = cbiem(CCase, D, sigma, O, alpha);
%
% Perform the CBIEM on the Dirichlet problem.
%
% David De Wit March 30 1992 - January 14 1993
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 0.0: Input and other parameters, togther with definitions.
if ~exist(’alpha’), alpha = 2/3; end
if ~exist(’O’), O = 12; end
if ~exist(’sigma’), sigma = 0.32; end
if ~exist(’D’), D = 7; end
if ~exist(’CCase’), CCase = 4; end
format short e; format compact;
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
%
% CCase:Index number to the contour being used:
% 1: Unit circle, with 4 equally spaced artificial corners.
% 2: Chandler’s Teardrop, one corner, right angled, at the origin.
% (This has been reversed to make it ACW. Now different
% from both GAC and DDW thesis.)
% 3: Kress’ ACW Teardrop, one corner, 2 pi/3 - angled.
% 4: Kress’ Reentrant contour, 3 pi/2 - angled. Reversed to
% avoid the branch cut on the negative real axis, and
% make it ACW in orientation.
% 5: ACW Cardioid. Reentrant contour with 2 pi interior angle.
% 6: ACW Heart. Reentrant contour with 2 pi and 0 interior angles.
% 7: ACW Controlled Cardioid, using trig parameterisation.
% 8: ACW Controlled Cardioid, using polynomial parameterisation.
% 9: Modified Boomerang, with a 5 degree external angle.
%
% D: Density of the geometric mesh, a positive integer. Choice
% of D forces N, the size of the linear system being solved,
% to be N = NC.(D+1)^2.
%
% sigma:Mesh parameter. Ratio of distances of consecutive mesh points
% from the nearest corner. 0 < sigma < 0.5. Try sigma = 0.25
% as a starting guess.
%
% O: Order of the interpolatory polynomial used to aproximate
% V at the collocation points. Actually the (even) number of
% nearest points used, thus 2 is linear, 4 is cubic. Must be
% kept 2 <= O <= NS+1; in practice keep O <= 20. This limits
% O <= 8 for D = 2, 16 for D = 3, etc.
%
% alpha:Exponent of the true solution of test problem, W = z^{alpha}.
%
% NC: Numbers of corners/segments of the contour.
%
% NS: Number of points on each side of the contour. N node points
% create NS = (N+1)^2 + 1 mesh points on each segment.
%
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 1.1: Generate quadrature rule for a geometrically graded mesh, on
% one segment in the t domain using hprmesh, then insert this
% into the grid with corners, using rmesh.
numcorn = [4 1 1 1 1 1 1 1 1 1]; NC = numcorn(CCase);
% graded mesh
G = [0 sigma.^(D:-1:1)]’; G = [G; 1-G(D+1:-1:1)];
S = [2:D+2, D+1:-1:2]’; c = [0:NC]’/NC;
[t, w] = hprmesh(G, S, 0); [tn, wn] = rmesh(c, t, w, 1);
N = length(tn);
% Uniform mesh
%%% N = (D+1)^2; tn = [1:N]’/N; wn = ones(size(tn))/N;
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 1.2: Compute the (complex) values of zn and zc (z at the node and
% collocation points specified by tn and tc, respectively).
% zc are found as the midpoints of zn in an arc-length sense, by
% mapping the midpoints of tc to the contour. Also calculate
% gdot, which is used to modify wn.
tc = tn - diff([0; tn])/2;
ptn = pi*tn; ptc = pi*tc;
if (CCase == 1)
zn = exp(2*i*ptn); zc = exp(2*i*ptc);
gdot = 2*pi*i*exp(2*i*ptn);
elseif (CCase == 2) % Chandler’s Teardrop
zn = 2*sin(ptn) - i*sin(2*ptn);
zc = 2*sin(ptc) - i*sin(2*ptc);
gdot = 2*pi*(cos(ptn) - i*cos(2*ptn));
gdot(N) = -2*pi*i;
elseif (CCase == 3) % Kress’s Teardrop
zn = sin(ptn)*2/sqrt(3) - i * sin(2*ptn);
zc = sin(ptc)*2/sqrt(3) - i * sin(2*ptc);
gdot = 2*pi*( cos(ptn)/sqrt(3) - i*cos(2*ptn) );
gdot(N) = -2*pi*i;
elseif (CCase == 4) % Kress’s Boomerang
a = 2/3;
zn = - a * sin(3*ptn) - i * sin(2*ptn);
zc = - a * sin(3*ptc) - i * sin(2*ptc);
gdot = - 2*pi*( (3*a/2)*cos(3*ptn) + i*cos(2*ptn) );
gdot(N) = -2*pi*i;
elseif (CCase == 5) % Plain Cardioid
zn = (-1 + cos(2*ptn)).*exp(i*2*ptn);
zc = (-1 + cos(2*ptc)).*exp(i*2*ptc);
gdot = -2*pi*exp(i*2*ptn).*(sin(2*ptn) + i*(1-cos(2*ptn)));
gdot(N) = 0;
elseif (CCase == 6) % Pointed Heart - Silly
zn = - sin(3*ptn) - i*(sin(2*ptn)).^3;
zc = - sin(3*ptc) - i*(sin(2*ptc)).^3;
gdot = - 3*pi*( cos(3*ptn) + i*2*cos(2*ptn).*(sin(2*ptn)).^2 );
gdot(N/2) = 0; gdot(N) = 0;
elseif (CCase == 7) % My Cardioid
zn = - sin(3*ptn) - 5 * i* tn .* (1 - tn) .* sin(2*ptn);
zc = - sin(3*ptc) - 5 * i* tc .* (1 - tc) .* sin(2*ptc);
gdot = - 3 * pi * cos(3*ptn) - 5 * i * ...
( (1-2*tn).*sin(2*ptn) + 2*ptn.*(1-tn).*cos(2*ptn) );
gdot(N) = 0;
elseif (CCase == 8) % My Stupid Polynomial Cardioid
a = 7/24;
zn = tn .* (tn - a) .* (tn - 1 + a) .* (tn - 1) + ...
i * tn.^2 .* (tn - 1/2) .* (tn - 1).^2;
zc = tc .* (tc - a) .* (tc - 1 + a) .* (tc - 1) + ...
i * tc.^2 .* (tc - 1/2) .* (tc - 1).^2;
gdot = (tn - a) .* (tn - 1 + a) .* (tn - 1) + ...
tn .* (tn-1+a) .* (tn-1) + tn .* (tn-a) .* (tn-1) + ...
tn .* (tn-a) .* (tn-1+a) + ...
i * ( 2 * tn .* (tn - 1/2) .* (tn - 1).^2 + ...
2 * tn.^2 .* (tn - 1/2) .* (tn - 1) + ...
tn.^2 .* (tn - 1).^2 );
gdot(N) = 0;
zn = 100*zn; zc = 100*zc; gdot = 100*gdot;
elseif (CCase == 9) % 5 degree external angle.
deg = 5;
a = 2/( 3 * tan(deg*pi/360));
zn = - a * sin(3*ptn) - i * sin(2*ptn);
zc = - a * sin(3*ptc) - i * sin(2*ptc);
gdot = - 2*pi*( (3*a/2)*cos(3*ptn) + i*cos(2*ptn) );
gdot(N) = -2*pi*i;
elseif (CCase == 10) % 20 degree external angle.
deg = 20;
a = 2/( 3 * tan(deg*pi/360));
zn = - a * sin(3*ptn) - i * sin(2*ptn);
zc = - a * sin(3*ptc) - i * sin(2*ptc);
gdot = - 2*pi*( (3*a/2)*cos(3*ptn) + i*cos(2*ptn) );
gdot(N) = -2*pi*i;
end
%rzn = real(zn); izn = imag(zn);
%plot(rzn,izn,’-’,rzn,izn,’+’); grid;
%plot(tn(1:10),izn(1:10),’-’,tn(1:10),izn(1:10),’+’); grid;
wn = wn.*gdot;
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 2.1: Set up the complex order N matrices A and B. Establish F, a
% matrix of indices to be used in calculating B, using Ft, a
% submatrix of the pattern of indices for one edge. Also compute
% L, the matrix of the coefficients of the interpolatory
% polynomial, using the indices contained in F.
oN = ones(N,1); j = O/2;
A = oN*wn.’ ./ ( oN*zn.’ - zc*oN’); NS = N/NC;
Ft = zeros(NS, O); Ft(1,:) = 0:O-1;
Ft(NS,:) = NS-O+1:NS;
Ft(2:j,:) = ones(j-1,1)*Ft(1,:);
Ft(NS-j+1:NS-1,:) = ones(j-1,1)*Ft(NS,:);
for k = j+1:NS-j, Ft(k,:) = Ft(k-1,:) + 1; end
for k = 1:NC, F((k-1)*NS+1:k*NS,:) = Ft + (k-1)*NS; end
F(1:j,1) = N*ones(j,1);
L = ones(size(F)); o1 = 1:j; o2 = j+1:N;
for k = 1:O, for v = 1:O
if (v ~= k)
tk = tn(F(o2,k)); tv = tn(F(o2,v));
L(o2, k) = L(o2, k).* (tc(o2) - tv) ./ (tk - tv);
tk = tn(F(1:j,k)); tv = tn(F(1:j,v));
if (k == 1), tk = 0; end; if (v==1), tv = 0; end;
L(o1, k) = L(o1, k) .* ( tc(o1) - tv ) ./ ( tk - tv );
end
end, end
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 2.2: Compute the matrix B. J is a shift vector, used to create rows
% of B from rows of L. Multiply the temporary result through by
% H, then write out C.
H = sum(A.’); t1 = ones(j,1);
t2 = [-t1; [0:NS-O-1]’; (NS-O)*t1]; J = zeros(size(H));
for k = 1:NC, J((k-1)*NS+1:k*NS) = t2 + (k-1)*NS; end
B(1:j,1:O-1) = L(1:j,2:O); B(1:j,N) = L(1:j,1);
for k = j+1:N, B(k,J(k)+1:J(k)+O) = L(k,:); end
C = B.*(H’*ones(size(H))) - A;
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 2.3: Set up and solve for Vn (computed approximation to V at the
% node points), the linear system:
% C Vn = A * Utn - sum(A).*Uc
% Given C, establish the real, order N-1 matrix, Cstar, then
% calculate Vn, using Vn(N) = 0.
%sprintf(’Solving the linear system of size %g’, N)
Un = real(funccb(zn, alpha)); Uc = real(funccb(zc, alpha));
d = diag(H)*Uc - A*Un; d = -imag(d(2:N) - d(1:N-1));
C = real(C(2:N,1:N-1) - C(1:N-1,1:N-1));
Vn = zeros(size(Un)); Vn(1:N-1) = C \ d;
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% 2.4: Error norms. Vne and Vce are the differences between the true
% and the computed values of V at the node and collocation points
% respectively. r is the discretisation error in the CBIE. p is
% the residual in the (above) computation for Vtn, when the true
% soln at the node points is substituted into the equations.
Vne = imag(funccb(zn, alpha)) - Vn;
Vnnorm = sqrt(abs(wn)’*(Vne.^2));
% Compute the approximation at the 4 points of Kress:
Wn = Un + i * Vn;
z1 = [ 0.1+0*i 0.2+0*i 0.3+0*i 0+0.2*i ];
for j = 1:4
WK(j) = sum(Wn .* wn ./ (zn - z1(j))) / sum(wn./(zn - z1(j)));
end
WKt = funccb(z1, alpha); Ek = [N O abs(real(WK - WKt))]
```
### A.2 cint.m
```
function [rho, C] = cint(sigma);
% function [rho, C] = cint(sigma);
%
% Contour integration with a geometric h-p grid.
%
% David De Wit July 14 1992 - July 17 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if ~exist(’QCase’), QCase = 2; end
if ~exist(’sigma’), sigma = 0.15; end
DMin = 8; DMax = 15;
format short e; format compact;
QCase = 2;
for D = DMin:DMax
G = [0 sigma.^(D:-1:1)]’; G = [G; 1-G(D+1:-1:1)];
S = [2:D+2, D+1:-1:2]’; NC = 2;
[t, w] = hprmesh(G, S, QCase, 0); c = [0:NC]’/NC;
[tn, wn] = rmesh(c, t, w, 1); zn = exp(2*pi*i*tn);
iN(D-DMin+1) = length(zn);
gdot = 2*pi*i*zn; wn = wn.*gdot;
ier(D-DMin+1) = 1-wn.’*(((zn-1)/i).^(1/2)./zn)/(2*pi*i*sqrt(i));
end
lier = log10(abs(ier))’; sqiN = sqrt(iN)’;
plot(sqiN,lier,’-g’,sqiN,lier,’+r’); grid;
title(’Contour integration on a h-p geometric grid’);
xlabel(’sqrt(N)’); ylabel(’log10(Error)’);
```
### A.3 funccb.m
```
function W = funccb(z, alpha);
% function W = funccb(z, alpha);
%
% True solution to the Dirichlet problem solved by cbiem.
%
% David De Wit April 13 1992 - September 27 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
W = z.^(alpha);
```
### A.4 funcci.m
```
function W = funcci(z);
% function W = funcci(z);
%
% Integrand of the problem solved by cint.
%
% David De Wit July 9 1992 - July 17 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
W = ((z-1)/i).^(1/2)./z;
```
### A.5 funchp.m
```
function f = funchp(x)
% function f = funchp(x);
%
% Function being integrated by hpmeth.
%
% David De Wit July 9 1992 - July 17 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
f = 1 - 3/2*sqrt(x);
```
### A.6 gettw.m
```
function [QRt, QRw] = gettw(R, QCase);
% function [QRt, QRw] = gettw(R, QCase);
%
% Get tables of nodes and weights for quadrature rules. User
% inputs maximum number of points required, and the type required.
% The default type is Gauss--Lobatto (QCase = 2), as it is of
% higher order than Newton--Cotes.
%
% David De Wit July 13 1992 - July 17 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if ~exist(’R’), R = 20; end
if ~exist(’QCase’), QCase = 2; end
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if (QCase == 1)
if (R > 10)
R = 10;
sprintf(’R too large for Newton--Cotes. Now R = 10.\n’);
end
QRw = [
1 1 1 7 19 41 751 989 2857 16067;
1 4 3 32 75 216 3577 5888 15741 106300;
0 1 3 12 50 27 1323 -928 1080 -48525;
0 0 1 32 50 272 2989 10496 19344 272400;
0 0 0 7 75 27 2989 -4540 5778 -260550;
0 0 0 0 19 216 1323 10496 5778 427368;
0 0 0 0 0 41 3577 -928 19344 -260550;
0 0 0 0 0 0 751 5888 1080 272400;
0 0 0 0 0 0 0 989 15741 -48525;
0 0 0 0 0 0 0 0 2857 106300;
0 0 0 0 0 0 0 0 0 16067
]
QRw = QRw(1:R+1,1:R); QRt = zeros(QRw);
for j = 1:R
QRw(1:j+1,j) = QRw(1:j+1,j)/sum(QRw(:,j));
QRt(1:j+1,j) = [0:j]’/j;
end
elseif (QCase == 2)
QRt = zeros(R+1,R); QRw = QRt;
for j = 2:R+1
[QRt(1:j,j-1), QRw(1:j,j-1)] = lobatto(j,0,1);
end
end
```
### A.7 hpmeth.m
```
function [lEip] = hpmeth(DMax, GMax, p, sigma, QCase);
% function [lEip] = hpmeth(DMax, GMax, p, sigma, QCase);
%
% Experiment with h-p integration methods.
%
% David De Wit July 9 1992 - July 17 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if ~exist(’QCase’), QCase = 2; end
if ~exist(’sigma’), sigma = 0.15; end
if ~exist(’p’), p = 6; end
if ~exist(’GMax’), GMax = 6; end
if ~exist(’DMax’), DMax = 19; end
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% Geometrically-graded h-p method.
Ehp = zeros(DMax,1); vN = Ehp;
for D = 1:DMax
G = [0 sigma.^(D:-1:0)]’; S = [2:D+2]’;
[t, w] = hprmesh(G, S, QCase, 0);
Ehp(D) = funchp(t)’*w; vNhp(D) = length(t);
end
lvNhp = log10(vNhp); lEhp = log10(Ehp);
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% Various linear h and p methods. Variable g, constant p.
Eip = zeros(DMax,GMax);
for D = 1:DMax
N = 2*D;
for g = 1:GMax
G = ([0:N]’/N).^g;
[t, w] = hprmesh(G, p, QCase, 0);
vN(D) = length(t); Eip(D,g) = funchp(t)’*w;
end
end
lvN = log10(vN); lEip = log10(Eip);
sprintf(’Order %f, slopes of the blue lines are approximately’, p)
(lEip(DMax-1,:) - lEip(DMax,:))./(lvN(DMax-1)-lvN(DMax))
plot(lvNhp,lEhp,’+r’,lvNhp,lEhp,’-g’,lvN,lEip,’+r’,lvN,lEip,’-b’);
xlabel(’log10(points)’); ylabel(’log10(error)’);
text(0.9,0.7,’g = 1’,’sc’); text(0.9,0.57,’g = 2’,’sc’);
text(0.9,0.45,’g = 3’,’sc’); text(0.9,0.35,’g = 4’,’sc’);
text(0.9,0.27,’g = 5’,’sc’); text(0.9,0.2,’g = 6’,’sc’);
text(0.7,0.15,’h-p method’,’sc’); grid;
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
% Compare h methods for quadrature rules of various p (# points).
OMax = 12;
for p = 2:2:OMax
g = p; j = p/2 + 1;
for D = 1:DMax
N = 2*D; G = ([0:N]’/N).^g;
[t, w] = hprmesh(G, j, QCase, 0);
Ehhp(D,j-1) = funchp(t)’*w; hhN(D) = length(t);
end
end
lhhN = log10(hhN); lEhhp = log10(Ehhp);
sprintf(’Slopes of the blue lines are approximately’)
(lEhhp(DMax-1,:) - lEhhp(DMax,:))./(lhhN(DMax-1)-lhhN(DMax))
plot(lvNhp,lEhp,’+r’,lvNhp,lEhp,’-g’,lhhN,lEhhp,’+r’,lhhN,lEhhp,’-b’);
xlabel(’log10(points)’); ylabel(’log10(error)’);
text(0.9,0.76,’p = 1’,’sc’); text(0.9,0.62,’p = 3’,’sc’);
text(0.9,0.5,’p = 5’,’sc’); text(0.9,0.4,’p = 7’,’sc’);
text(0.9,0.32,’p = 9’,’sc’); text(0.85,0.2,’p = 11’,’sc’);
text(0.7,0.15,’h-p method’,’sc’); grid;
```
### A.8 hprmesh.m
```
function [t, w] = hprmesh(G, S, IClosed)
% function [t, w] = hprmesh(G, S, IClosed)
%
% Create a new quadrature rule, based on a mesh G, where between
% points G(i) and G(i+1) is a (closed) quadrature rule of
% Gauss--Lobatto type on S(i) points, including the 2 end points
% G(i) and G(i+1); for i = 1:length(G)-1. If IClosed is 1, then
% the contour is closed, and the ends are tied together. This
% function is a generalisation of rmesh.
%
% David De Wit July 13 1992 - December 2 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if ~exist(’G’), sigma = 0.1; G = [0 sigma.^(3:-1:1) 1]’; end
if ~exist(’S’), S = [2:5]’; end
if ~exist(’IClosed’), IClosed = 1; end
lG = length(G); lS = length(S);
if ((lG ~= lS+1) & (lS ~= 1))
sprintf(’hprmesh: Danger l(G) = %f, l(S) = %f’, lG, lS)
end
% Set up S for rules with a constant integration rule.
if (lS == 1), S = ones(lG-1,1)*S; end
dG = diff(G); S = S - 1;
N = max(S); R = length(dG);
% Obtain the nodes and weights in a table
QRt = zeros(N+1,N); QRw = QRt;
for j = 2:N+1
[QRt(1:j,j-1), QRw(1:j,j-1)] = lobatto(j,0,1);
end
% Play with the table
QRw(N+1,:) = diag(QRw(2:N+1,1:N))’; QRt(N+1,:) = ones(1,N);
for i = 2:N, for j = 1:i-1, QRt(i,j) = NaN; QRw(i,j) = NaN; end, end
tt = QRt(:,S); tw = QRw(:,S);
j = ones(N+1,1); tw = tw.*(j*dG’);
tt = tt.*(j*dG’) + j*G(1:R)’;
tw(1,2:R) = tw(1,2:R) + tw(N+1,1:R-1);
tw1 = tw(1:N,:); tt1 = tt(1:N,:);
t = [tt1(:); tt(N+1,R)]; t(isnan(t)) = [];
w = [tw1(:); tw(N+1,R)]; w(isnan(w)) = [];
if (IClosed == 1)
N = length(t);
w = [w(2:N-1); w(1)+w(N)]; t = t(2:N);
end
```
### A.9 lobatto.m
```
function [x, w] = lobatto(n, a, b)
% function [x, w] = lobatto(n, a, b)
%
% Return the weights w and points x of the n-point Gauss--Lobatto
% quadrature rule on the interval [a, b].
% See G. H. Golub, SIAM Review 1973 p 318.
%
% Graeme Chandler July 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
n = round(n);
if (n == 2)
x = [a; b]; w = [1; 1]*(b-a)/2;
elseif (n == 3)
x = [a; a+(b-a)/2; b]; w = [1; 4; 1]*(b-a)/6;
elseif (n >= 4)
nn = n-1; m = 1:2:2*nn-1;
m = (1:nn-1) ./ sqrt(m(1:nn-1) .* m(2:nn));
J = (diag(m,-1)+diag(m,1));
I = eye(nn); en = (1:nn)’ == nn;
gam = (J + I)\en; mu = (J - I)\en;
sol = [1 -gam(nn); 1 -mu(nn)]\[-1; 1];
alpha = sol(1); beta = sqrt(sol(2));
[ww,xx] = eig([J beta*en; beta*en’ alpha]);
[xx, i] = sort(diag(xx));
w = ww(1,i)’.^2 * (b-a);
x = [a; (a+b)/2+(b-a)*xx(2:nn)/2; b];
end
```
### A.10 rmesh.m
```
function [t, w] = rmesh(G, t1, w1, IClosed);
% function [t, w] = rmesh(G, t1, w1, IClosed);
%
% Create a new quadrature rule, based on a mesh G, where a
% (closed) quadrature rule (t1, w1) is inserted over each
% interval of G. If IClosed is 1, then the contour is closed,
% and the ends are tied together. This function is generalised
% into hprmesh. Originally conceived by Graeme Chandler.
%
% David De Wit July 13 1992 - September 6 1992
%
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if (length(G) < 2), return; end
if (length(G) ~= 2)
h = diff(G); tw = w1*h’;
n = length(h); m = length(w1);
tw(1,2:n) = tw(1,2:n) + tw(m,1:n-1); tw1 = tw(1:m-1,:);
w = [tw1(:); tw(m,n)];
tt = t1(1:m-1)*h’ + ones(m-1,1)*G(1:n)’;
t = [tt(:); G(n+1)];
else
t = t1; w = w1;
end
if (IClosed == 1)
N = length(t);
w = [w(2:N-1); w(1)+w(N)]; t = t(2:N);
end
```
### A.11 testcb.m
```
function N = testcb(C, sigma, Dmin, Dmax, Omin, Omax, alpha)
% function N = testcb(C, sigma, Dmin, Dmax, Omin, Omax, alpha)
%
% Run cbiem for various parameters, and tabulate results.
%
% David De Wit May 12 1992 - December 21 1992
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
if ~exist(’alpha’), alpha = 1/2; end
if ~exist(’Omax’), Omax = 16; end
if ~exist(’Omin’), Omin = 8; end
if ~exist(’Dmax’), Dmax = 19; end
if ~exist(’Dmin’), Dmin = 16; end
if ~exist(’sigma’), sigma = 0.32; end
if ~exist(’C’), C = 7; end
format short e; format compact
%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%%
for D = Dmin:Dmax
for O = Omin:2:Omax
P = cbiem(C, D, sigma, O, alpha);
if ((P ~= Inf) & (P ~= NaN))
N(D-Dmin+1, (O-Omin)/2+1) = P;
else
N(D-Dmin+1:Dmax-Dmin+1,:) = ...
Inf*ones(Dmax-D+1,(Omax-O)/2 + 1);
return
end
end
if (min(N(D-Dmin+1,:)) > 1)
N(D-Dmin+2:Dmax-Dmin+1,:) = ...
Inf*ones(Dmax-D,(Omax-Omin)/2+1);
return
end
end
``` |
warning/0002/cond-mat0002319.html | ar5iv | text | # On depolarisation level shift in spherical QD
## I Introduction
A complete quantum dot (QD) theory, taking into account all the sophisticated physics of the object, is still a challenge for a theorist. The main reason is that the scale of the calculation is much larger than atomic one (that complicates ab initio techniques). The same time the number of particles is small to use solid state approximations in full extent. For example, an one–electron picture of a quantum confinement potential, arising from the conduction band discontinuity on the QD boundary, does not always yield accurate electron levels.
In the paper we put forward a model to inspect the electrodynamical correction to the one–electron energy in a spherical QD. It was shown that the similar correction occurs to be significant for a ”natural quantum dot” C<sub>60</sub>. A depolarisation level shift due to the interaction with an electromagnetic field is not negligible, as it might be thought, when taking into account localized electromagnetic modes. We present the scaling analysis of some different mechanisms for the level shift (LS) and propose a possible experimental manifestation of the depolarisation effect.
In order to appraise the LS a simple spherical QD model in frame of an effective mass approximation was applied. How is our result sensitive to the model used? The size scaling of the depolarisation shift preserves, being dependent mainly on a corresponding density of states of the field, while the prefactor might be smaller within other approach, though it is not easy to evaluate explicitly.
The group of full rotations, SO(3), was chosen to label the one–electron states. It is possible to perform an analytic quantum–mechanical calculation of the RPA response within the spherical model. The massive peak of a collective excitation is known to show up in the spectrum, resulting from the fast coherent oscillation of the total electron density of the valence states. Thus, within our model the electron–electron interaction is dealt with selfconsistently. Of course, the number of valence electrons involved in the collective motion has not to be small. It is believed to fulfill for the typical QD possessing some hundreds of atoms and even more.
A surface charge density oscillation can be thought as a confined electric field mode or a multipole surface plasmon. We will reflect on the shift of the electron level in the field of zero–point oscillations of the modes connected with the QD, which depolarisation effect is billion times stronger than of the free field zero–point oscillations, so a name ”giant LS” is admitted.
The classical description of the electromagnetic surface modes, via the dielectric function of the matrix and QD material, gives the true plasmon state frequencies and will be used below. Once more, the final result does not depend too much on the computation approach. Our model sketches out the (many–body) depolarisation semiclassically, avoiding a lot of the routine computational intricacy.
The paper proceeds as follows: the brief model description is next to the introduction. Then, the model will be applied to a 3D–plasmon as well as a free field mode, that will explain the calculation technique. However, the LSs from these modes are too small to have an experimental importance. Section IV deals with the confined modes those result in much larger depolarisation shift. The numerical estimations and the scaling properties of the LS will be given in respect to a possible experiment. A brief summary will follow. The calculation of the QD confined mode frequencies is allocated in Appendix.
## II Semiclassical theory for energy level shift
We have considered semiclassically the LS for an arbitrary shell object in . The method follows the one proposed by Migdal to calculate the Lamb shift for a hydrogen–like atom. The frequency of the zero–point oscillations of the external field is much higher than the inverse period of the electron orbit $`\omega _p\pi /\tau `$. Therefore, the adiabatic approximation has to be used and one divides the fast (field) and slow (electron) variables. An electron is subjected to short fast deflections from its original orbit in the high–frequency field of the electromagnetic wave of the zero–point oscillation. Then the energy shift is given by the second order perturbation theory as
$$\delta E=H(r+\delta )H(r)=H\stackrel{}{\delta }+\frac{1}{2}^2H\stackrel{}{\delta }\stackrel{}{\delta }+\mathrm{},$$
(1)
where $`H(r)`$ is the unperturbed Hamiltonian and $`H(r+\delta )`$ is the Hamiltonian with account for the random electron deflection $`\delta `$. The angle brackets represent the quantum mechanical average over the fast variables of the field (or, the same, over the random electron deflections). The perturbed Hamiltonian is expanded in series on the $`\delta `$ and a first nonzero contribution is taken.
The simplest QD Hamiltonian is considered to have only the rotational correction which is given by:
$$\delta H=\frac{\widehat{L}^2}{2mR^2}\left(2\frac{\delta }{R}+3\frac{\delta ^2}{R^2}+\mathrm{}\right),$$
(2)
where $`R`$ is about the spherical QD radius; $`m`$ is the electron mass which is supposed to be constant within the dot; $`\widehat{L}`$ is the angular momentum operator. On the averaging, the first–order term disappears. So far the LS dependence on the QD size includes $`R^4`$ factor besides some power hidden in the mean square deflection $`\overline{\delta ^2}`$. The strength of the electrodynamical interaction changes with this quantity almost exactly. We will show that the dependence of $`\overline{\delta ^2}`$ in $`R`$ is different for different electric modes (confined and free field). The giant deflection is representative for the giant LS and, therefore, the function $`\overline{\delta ^2}(R)`$ will be studied specifically.
## III Bulk plasmon contribution to LS
First we consider the bulk 3D–plasmon modes that could shift the electron level. Nearly self–evidently the bulk plasmon shift is negligible. The mean square deflection, caused by the 3D mode (which is not confined at all), decreases with the QD size too rapidly. The small factor, contained in the 3D LS, comes essentially from the expression for $`\overline{\delta ^2}`$ which scales as $`1/N`$, where $`N`$ is the number of atoms in the QD. It will be explained in this section.
Within the semiclassical approach, the deflection of the electron can be computed with the use of the Newton law:
$$m_t^2\delta =e,$$
(3)
here $`e`$ is the electron charge, $``$ is the field strength due to the zero–point oscillation of some mode, $`m`$ is the electron effective mass. The square of the deflection $`\overline{\delta ^2}=e^2/(2m)^2d^Dk\overline{_k^2}/\omega _k^4`$ is proportional to the mean square of the electric field strength. The dimension of the field, $`D`$, equals 3. The field strength, in turn, can be rewritten as the zero–point oscillation frequency $`\overline{_k^2}=2\pi \mathrm{}\omega _k`$ through the quantized field normalisation.
The scale of the energy is given by the 3D plasmon frequency $`\omega _p=\sqrt{4\pi e^2n_{3\mathrm{D}}/m}`$. Note that the 3D plasmon frequency does not depend on the quantum number $`𝐤`$ and passes through the integral. Hence, the mean square deflection contains the total number of states effecting on the electron level in the QD. The integral is limited above by $`k_{max}1/R`$. In 3D–case it brings the factor $`R^3N^1`$ claimed in the beginning of the section.
This result will change for other confined electric modes because of their different densities of states. This produces the different $`N`$scaling factor for the LS from these modes.
The prefactor of the deflection, for any mode considered here, depends equally on the square root of the density of electrons, which is useful to be converted to $`r_s`$, a characteristic length via the following definition: $`2\pi r_s^3n_{3\mathrm{D}}/3=1`$. Then, for 3D plasmon the deflection reads as:
$$\overline{\delta ^2}=a_B^2\frac{\sqrt{6}}{6^4\pi }\left(\frac{r_s}{a_B}\right)^{3/2}\left(\frac{r_s}{R}\right)^3N^1,$$
(4)
where the atomic length unit, $`a_B=\mathrm{}^2/me^20.53`$Å, or the Bohr radius, gives the scale of the deflection (note that this definition does not include any permittivity unlike an exciton Bohr radius in semiconductors).
The depolarisation (the ratio of the level shift, $`\delta E`$, to the bare energy, $`E^{(o)}`$) due to the 3D modes is as follows:
$$\mathrm{\Delta }_{3\mathrm{D}}=\frac{\delta E}{E^{(o)}}=\frac{1}{72\sqrt{6}\pi }\sqrt{\frac{a_B}{r_s}}\left(\frac{r_s}{R}\right)^5N^{5/3}.$$
(5)
The rude estimation of the prefactor shows that even for the small QD with $`N=100`$ the shift is $`10^6`$ of the bare energy and will not be resolved because of a number of other different factors effecting the level position.
To give a complete picture we note that the standard LS due to the zero–point oscillations of the free electromagnetic modes of the vacuum can be written as:
$$\mathrm{\Delta }_{\mathrm{vac}}=\frac{6\alpha ^3}{\pi }\left(\frac{a_B}{r_s}\right)^2\left[\mathrm{ln}\frac{r_s}{\alpha a_B}+\mathrm{ln}\frac{R}{r_s}\right]\left(\frac{r_s}{R}\right)^2N^{2/3},$$
(6)
where $`\alpha 1/137`$ is the fine structure constant, and the simple check shows that the logarithmic dependence of the last term in square brackets on $`N`$ does not add any extra to the result and has to be dropped in our case. Though the slope of the LS in $`N`$ is much slower than in Eq.(5) the prefactor is tiny ($`10^7`$) because of $`\alpha ^3`$.
## IV Depolarisation: confined modes
Let us consider the specific behavior of the LS materialized by the zero–point oscillations of the confined plasmon modes. The depolarisation in carbon shell cluster was shown to be independent of the cluster size. The mean square deflection scales also as a zero power of the size $`\overline{\delta ^2}N^0`$. While it is interesting by itself, the carbon cluster matter will not be considered in the paper. However, there are confined modes in our QD problem those enhance the electrodynamical correction to the electron energy.
With the decrease of the dimension of the field the plasmon density of states increases. Hence, the mean interaction of the electron with the plasmon field increases that will be evident from the scaling of $`\overline{\delta ^2}`$.
Two possible candidates for the confined plasmon modes in the QD system, those have different densities of states, are the 2D plasmon and the 0D spherical mode. The former mode can arise because of some interface possibly grown within the structure (see inset in the Fig.1). It might be a conducting wetting layer, if it is thick enough to confine the electromagnetic field. The 2D plasmon naturally originates at the interface between the semiconductor structure and a metal . At the boundary of two dielectrics a surface plasmon is known to propagate. Its contribution will be discussed elsewhere as being smaller than 2D–plasmon one by a factor $`10^2`$ at least owing to the fast space decay.
The 0D mode is the inherent property of the spherical inclusion of the foreign material in any matrix. The calculation of the frequency of this mode is slightly cumbersome (see Appendix for details). The surface QD mode has the quantum numbers $`L,M`$, the angular momentum and its projection on an axis, instead of 2D wave vector, $`𝐤`$, for the standard 2D plasmon modes. The depolarisation is anomalous large in the 0D case. It will be seen in this section from the scaling of $`\overline{\delta ^2}`$ and $`\mathrm{\Delta }`$.
### A 2D plasmon
The frequency of 2D plasmon is well known to depend on its 2D wave vector as: $`\omega _k=\sqrt{2\pi e^2n_{2\mathrm{D}}k/m}`$. We will rewrite the 2D electron density, as before, in terms of the characteristic length: $`\pi r_s^2n_{2\mathrm{D}}/2=1`$ and perform the integration over the plasmon states. Then the mean square deflection can be expressed as:
$$\overline{\delta ^2}=a_B^2\frac{1}{32}\left(\frac{r_s}{a_B}\right)^{3/2}\left(\frac{r_s}{R}\right)^{3/2}N^{1/2}.$$
(7)
The scaling in $`N`$ has a lower exponent that reflects the different density of the confined field (plasmon) states. Substituting $`\overline{\delta ^2}`$ into the Hamiltonian given by Eq.(2), one gets the depolarisation as follows:
$$\mathrm{\Delta }_{2\mathrm{D}}=\frac{3}{32}\sqrt{\frac{a_B}{r_s}}\left(\frac{r_s}{R}\right)^{7/2}N^{7/6}.$$
(8)
The shift depends on the inverse size nearly linearly. However, the prefactor dominates at some moderate size of the QD and lessens the LS to $`10^3`$ for $`N=100`$. The depolarisation is still to be too small to expect experimental consequences. To be precise the result also depends on $`w`$, the distance between the 2D electrons and the QD. It is simply included in the consideration by multipling Eq.(8) by a factor $`\sqrt{\pi }\mathrm{Erf}(\sqrt{w/R})/(2\sqrt{w/R})`$, and the depolarisation declines 4 times at $`w/R10`$.
### B QD confined plasmon: Mode of cavity
The $`\overline{\delta ^2}`$ considered above the less, the larger the QD size, that is not the case for the giant deflection due to the completely localized modes. The localized modes are the surface plasmons of the spherical inclusion (with the dielectric function $`ϵ_1`$) in the matrix (with the different dielectric function $`ϵ_2`$). The frequency of the mode, $`\omega _L`$, that we consider, is nearly the frequency of the bulk plasmon in the matrix, $`\omega _{p2}`$, with the weak dependence on the mode angular momentum (see Appendix). The electric field of the zero–point oscillation is given by the formula $`_L^2=\pi (L+1/2)\mathrm{}\omega _L/R^3`$. The summation over all states below some critical value $`L_c`$ gives the mean square deflection:
$$\overline{\delta ^2}=a_B^2\frac{\pi }{9\sqrt{6}}\left(\frac{r_s}{a_B}\right)^{3/2}\left(\frac{r_s}{R}\right)^3\left(L_c+\frac{1}{2}\right)^3,$$
(9)
where it is natural to limit the summation above the excitation which wavelength is about the lattice constant $`d`$. We found that the $`\overline{\delta ^2}`$ does not depend on the QD size:
$$\overline{\delta ^2}=a_B^2\frac{\pi ^4}{9\sqrt{6}}\left(\frac{r_s}{a_B}\right)^{3/2}\left(\frac{r_s}{d}\right)^3N^0.$$
(10)
Sequently, the level shift depends on the size as $`R^2`$ (which comes from Eq.(2)):
$$\mathrm{\Delta }=\frac{\pi ^4}{3\sqrt{6}}\sqrt{\frac{a_B}{r_s}}\left(\frac{r_s}{d}\right)^3\left(\frac{r_s}{R}\right)^2N^{2/3}.$$
(11)
Our estimation shows that the level correction, becoming of the order of 50%, plays the important role for the QD of 100 atoms and smaller. We collected all studied contributions to the depolarisation LS and plot them in the log–log scale versus the QD size in Figure 1.
The depolarisation because of the localized surface QD modes is large enough to propose an experiment supporting our model. It is easy to see that $`\overline{\delta ^2}\omega _L^3`$, whence the LS depends on the mode frequency as well. Therefore, changing the optical properties of the matrix surrounding the QD, one shifts the levels. If the bare energy level, $`E^{(o)}`$, lies deep in the potential well, its position is nearly independent of the well depth which changes along with the matrix parameters. The deep bare level energy depends only on the well width $`R`$. Hence, keeping the same size of the QD and covering it with the different materials, one will derive solely the depolarisation LS, since it is distinguishable from the standard space quantization LS.
## V Summary
The effect of the zero–point oscillations of the free and confined electromagnetic field on the level of the confined electron in the spherical QD is reviewed. The depolarisation due to an interaction with the zero–point oscillations of the field (produced by all other valence electrons) shifts up the bare one-electron state that seems to be a counterpart for the vertex correction (electron–hole interaction, for example) which lowers the transition frequency down. It indicates that the studied effect should be taken into account for a many–body computation of a QD spectrum.
To the best of our knowledge, the scaling dependence of the depolarisation level shift for the QDs is calculated in the first time. The size dependence of the LS is different for 4 cases considered in the paper. This scaling reflects that the different densities of states work in different mechanisms of the depolarisation due to the different 3D, 2D and 0D–dimensional modes of the electric field are involved. Our model allows a theorist to skip a tedious quantum electrodynamical calculation but obtain the analytical selfconsistent estimation for the (many–body) level shift in a nanoscale system with the strong quantization. The result has not only a theoretical importance.
Although, the depolarisation decreases with the QD size in general, the localized surface electromagnetic mode (which is specific to the QD as a void in the matrix material) results in the giant level shift and is to be possibly resolved experimentally for the QD made from some hundred atoms. Another method to detect the effect could be the measurement of a deep level position in the similar QDs buried by the substances with the distinct optical characteristics. Then the localized plasmon frequency changes along with the prefactor of the depolarisation shift, which could be observed by the optical spectroscopy of the QD system.
## A Surface QD plasmon modes
The sought–for modes are given in complete spherical harmonics $`P_L(r)Y_{L,M}(\mathrm{\Omega })`$, where $`P_L(r)`$ is the Legendre polynomial and $`Y_{L,M}`$ is the spherical harmonic. The electrodynamic solution for the modes of the system consisting of the spherical particle with the dielectric function $`ϵ_1(\omega )`$ and the surrounding matrix with the different dielectric function $`ϵ_2(\omega )`$ is one of the roots of the secular equation:
$$\frac{ϵ_2(\omega )}{ϵ_1(\omega )}=\frac{L}{L+1}.$$
(A1)
The similar equation gives a mode of an empty void in the matrix in the limit $`ϵ_1=1`$ (in the limit $`ϵ_2=1`$ it gives a mode of a sphere in the vacuum). The right hand side varies from $`1/2`$ to $`1`$.
Let us suppose the simplest form for the dielectric function defined by one pole at some frequency, which is about the optical gap $`E_g`$, and one zero, which is given by the bulk plasma frequency of the material (both parameters differ for the materials 1 and 2). Then the solutions of the Eq.(A1) are easy to find graphically (see. Fig. 2). The function of $`\omega `$ standing in the left hand of the expression is plotted in solid. It has two poles at $`E_{g2}`$ and $`\omega _{p1}`$ and two zeroes at $`E_{g1}`$ and $`\omega _{p2}`$. All roots with $`L=1,2..L_c`$ lie in between these two pairs as shown in the figure. The dependence of the mode frequency on the angular momentum is very weak.
The lower–frequency modes, shown as black circles, correspond to the sphere–like plasmons. As lying above the QD gap, the lower mode effectively damps inside the QD. Therefore we will consider only the upper branch, shown as open circles, which is similar to the modes of the cavity. The frequency of the upper mode is close to the frequency of the bulk plasmon $`\omega _{p2}`$ as it is seen from Figure (2). The mode frequency dependence on its angular momentum is negligible.
The change of any of the 4 parameters, defining $`ϵ_1`$ and $`ϵ_2`$, results in the QD mode frequency shift. However, for the sought–for cavity–like mode, the most important are the plasma frequencies of the materials. This provides the mechanism for an experimental observation of the described effect in the matrix materials with the different $`\omega _p`$. It influences on the depolarisation LS, while the space quantization of the one–electron level, which is deep enough, depends solely on the QD width which has to be kept. |
warning/0002/math0002043.html | ar5iv | text | # The toric cobordisms
## The toric cobordisms of $`3`$-manifolds
The quotient $`M_\phi =T^2\times I/(x,0)(\phi (x),1)`$ where $`\phi `$ is a diffeomorphism of $`T^2`$ is a torus bundle over the circle with monodromy $`\phi `$. If the diffeomorphisms $`\phi _1,\phi _2`$ are isotopic, then $`M_{\phi _1}=M_{\phi _2}`$. A diffeomorphism of $`T^2`$ is determined up to isotopy by its induced map on the first integral homology group, and hence the diffeotopy group is isomorphic to $`GL(2,)`$.
Let $`G`$ denote the group $`GL(2,`$); let $`G^{}`$ be its commutator subgroup and $`G^2`$ be the subgroup of $`G`$ generated by the squares of the elements of $`G`$.
###### Claim 1.
Let $`M_\phi `$ be a 3-dimensional torus bundle defined as above. There exists a $`4`$-dimensional manifold $`W`$ fibered by tori over an orientable surface with $`W=M_\phi `$ if and only if $`\phi G^{}`$. Such a $`W`$ is orientable if and only if $`\phi `$ can be written as $`\phi =_{j=1}^g[\phi _{2j1},\phi _{2j}]`$ where $`\mathrm{det}\phi _i=1`$ for $`i=1,\mathrm{},2g`$.
###### Proof.
The $`T^2`$-bundles over a finite cell complex $`X`$ are classified by homotopy class of maps from $`X`$ to $`BDiff(T^2)`$. Such a class determines a conjugacy class of homomorphisms
$$\pi _1(X)\pi _1(BDiff(T^2))\pi _0(Diff(T^2))G.$$
In our case $`X`$ is a surface with non-empty boundary, so $`\pi _1(X)`$ has cohomological dimension 1 (being free). Thus $`T^2`$-bundles over a surface with non-empty boundary are in bijection with the homomorphisms $`\pi _1(X)G`$.
The manifold $`W`$ we are looking for exists if and only if there is a commutative diagram
where $`i_{}`$ is induced by inclusion of the boundary, i.e. $`Imi_{}(F_r)^{}`$, hence we have shown the first part of the claim.
To see when $`W`$ is orientable, use the standard construction of an oriented surface from the disk $`D^2`$, by identifying some $`1`$-disks $`I_i`$ on its boundary. Then we can explicitly construct $`W`$ starting from $`D^2\times T^2`$ by gluing $`(I_i\times T^2)`$’s on its boundary. As $`D^2\times T^2`$ is oriented and the gluing must create no orientation-reversing loop, it is not hard to see that the condition $`\mathrm{det}\phi _i=1`$ assures the orientability of $`W`$. ∎
Reasoning similarly, we get for fiber bundles over a non-orientable surface
###### Claim 2.
Let $`M_\phi `$ be a 3-dimensional torus bundle. There exists a $`4`$-dimensional manifold $`W`$ fibered by tori over a non-orientable surface with $`W=M_\phi `$ if and only if $`\phi G^2`$. Such a $`W`$ is orientable if and only if $`\phi `$ can be written as $`\phi =_{j=1}^k\phi _j^2`$ where $`det\phi _i=1`$ for $`j=1,\mathrm{},k`$.
###### Claim 3.
Let $`\phi _1\mathrm{}\phi _nG=GL(2,)`$ and $`M_{\phi _1},\mathrm{},M_{\phi _n}`$ be the corresponding 3-dimensional torus bundles. There exists a $`4`$-dimensional manifold $`W`$ fibered by tori over an orientable surface with $`W=M_{\phi _1}\mathrm{}M_{\phi _n}`$ if and only if $`_{i=1}^n\phi _iG^{}.`$
###### Proof.
The proof is immediately obtained from Claim 1 and the fact that the disjoint union $`M_{\phi _1}\mathrm{}M_{\phi _n}`$ is cobordant to $`M_\psi `$ with $`\psi =_{i=1}^n\phi _i,`$ by a toric cobordism with base the sphere $`S^2`$ with $`n+1`$ holes (it can be constructed similarly to the proof of Claim 1). ∎
###### Claim 4.
For $`\phi _1\mathrm{}\phi _nG`$ and $`M_{\phi _1},\mathrm{},M_{\phi _k}`$ as above, there exists a $`4`$-dimensional manifold $`W`$ fibered by tori over a non-orientable surface with boundary $`W=M_{\phi _1}\mathrm{}M_{\phi _n}`$ if and only if $`_{i=1}^n\phi _iG^2.`$
###### Lemma 1.
There exists an oriented toric cobordism with an orientable base between $`M_\phi `$ and $`M_\psi `$ if and only if there also exists an oriented toric cobordism with a non-orientable base between them.
###### Proof.
By $`G_{}^2`$ we denote the subgroup of $`GL(2,)`$ generated by the squares of matrices with negative determinant
$$G_{}^2=\{a_1^2a_2^2\mathrm{}a_k^2deta_i=1\}.$$
It is evident that $`G_{}^2SL(2,)`$ and is normal in it. We show that $`G_{}^2=(SL(2,))^{}`$ and this implies the claim. We use the following presentations of $`GL(2,)`$ and $`SL(2,)`$; see (, 2.23). For $`A=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$, $`B=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 1\end{array}\right)`$ and $`R=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$,
$$GL(2,)=A,B,R|A^2=B^3,A^4=R^2=(RA)^2=(RB)^2=1,$$
$$SL(2,)=A,B|A^2=B^3,A^4=1.$$
The commutator subgroup $`(SL(2,))^{}`$ of $`SL(2,)`$ is a free group of rank 2 generated by $`[A,B]=(ARB^1)^2G_{}^2`$ and $`[A,B^1]=(B^1RB)^2G_{}^2.`$ Thus $`(SL(2,))^{}G_{}^2`$, and as $`(SL(2,))^{}\mathrm{}SL(2,)`$, we have $`(SL(2,))^{}\mathrm{}G_{}^2`$.
By using the relations $`RA=A^1R`$ and $`RB=B^1R`$, each element $`aGL(2,)`$ can be written in the normal form $`R^\epsilon A^{k_1}B^{l_1}\mathrm{}A^{k_n}B^{l_n}`$ where $`\epsilon \{0;1\}`$. If $`deta=1`$, the element $`a`$ can be written in the form $`RA^{k_1}B^{l_1}\mathrm{}A^{k_n}B^{l_n}`$.
Now $`G_{}^2/(SL(2,))^{}=\{(RA^iB^j)^2|A^2=B^3,(RA)^2=(RB)^2=A^4=R^2=1,AB=BA\}=1.`$ Thus, $`G_{}^2=(SL(2,))^{}`$. ∎
###### Proof of Theorem 1.
It follows from Lemma 1 that for $`\phi ,\psi SL(2,)`$ there exists an oriented toric cobordism between the torus bundles $`M_\phi `$ and $`M_\psi `$ if and only if $`\phi \psi ^1(SL(2,))^{}.`$ Thus
$$\mathrm{\Omega }_3^{T^2,or}SL(2,)/(SL(2,))^{}=A,B|A^2=B^3,A^4=1,AB=BA_{12}.$$
The generator of $`\mathrm{\Omega }_3^{T^2,or}`$ is the conjugacy class of the element $`B^1A^1=`$ $`\left(\begin{array}{cc}1& 1\\ 0& 1\end{array}\right)`$.
The subgroup $`G^{}`$ lies in the subgroup $`G^2`$. Thus, if there exists an unoriented toric cobordism with an orientable base between $`M_\phi `$ and $`M_\psi `$, then there also exists an unoriented toric cobordism with a non-orientable base between them. Hence, Claim 4 lets us calculate the third group of unoriented toric cobordims, $`\mathrm{\Omega }_3^{T^2,unor}`$, as well. Thus,
$$\mathrm{\Omega }_3^{T^2,unor}G/G^2=A,B,R|A^2=B^3,A^4=R^2=(RA)^2=(RB)^2=1,$$
$$AB=BA,AR=RA,BR=RB,A^2=B^2=1_2_2.$$
The generators here are the conjugacy classes of $`A=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$ and $`R=`$ $`\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right).`$
###### Remark 1.
The toric cobordism is a cobordism of manifolds with some fixed torus bundle structures. In some cases, for $`M_\phi ,M_\psi `$ that are not oriented toric cobordant, we can choose other torus bundle structures on their total spaces in such a way that they become oriented toric cobordant. For example, take a $`\phi SL(2,)`$ such that $`\phi ^2(SL(2,))^{}`$. Then, the corresponding $`M_\phi `$ is not oriented toric cobordant to $`M_\phi ^{}`$, but is oriented toric cobordant to $`M_{\gamma \phi \gamma ^1}`$ for $`\gamma `$ with $`det\gamma =1`$. The reason for this is that the total spaces of $`M_\phi `$ and $`M_{\gamma \phi \gamma ^1}`$ are homeomorphic by a fiber preserving homeomorphism inducing the orientation reversing map of the basis $`S^1`$ (see or ).
###### Remark 2.
Let $`\phi SL(2,)`$ and $`W^4`$ be a toric cobordism between $`M_\phi `$ and $`\mathrm{}`$ that has an orientable base of genus $`g`$. Then $`\phi `$ is a product of $`g`$ commutators. By taking $`\psi ^{},\psi ^{\prime \prime }`$ such that $`[\psi ^{},\psi ^{\prime \prime }]=1`$, one can write $`\phi `$ as a product of $`g+1`$ commutators and so one can construct another toric cobordism between $`M_\phi `$ and $`\mathrm{}`$ with orientable base of genus $`g+1`$. Thus, quite naturally we come to the following question: what is the minimal genus of the orientable base of $`W^4`$?
For this, we can utilize Culler’s algorithm () which determines, for finite groups $`A`$ and $`B`$ and $`a(AB)^{}`$, the minimal number of elements of $`AB`$ required to represent $`a`$ as a product of their commutators. In order to extend Culler’s algorithm from free products to $`SL(2,)=_4__2_6`$, consider the projection homomorphism $`\alpha :_4__2_6_2_3.`$ As Culler’s algorithm can be applied to the group $`_2_3`$, it remains to note that the restriction of the homomorphism $`\alpha `$ to the commutator subgroup
$$\alpha ^{}=\alpha |_{(_4__2_6)^{}}:(_4__2_6)^{}(_2_3)^{}$$
is an isomorphism.
## Acknowledgment
The author is grateful to Gilbert Levitt for pointing out Culler’s paper. It is also a pleasure to thank the referee for constructive comments. |
warning/0002/hep-ex0002004.html | ar5iv | text | # Three generation neutrino mixing is compatible with all experiments
(Universidad San Francisco de Quito
2 February 2000
)
## Abstract
We consider the minimal extension of the Standard Model with three generations of massive neutrinos that mix. We then determine the parameters of the model that satisfy all experimental constraints.
PACS 14.60.Pq, 12.15.Ff
Three observables in disagreement with the Standard Model of Quarks and Leptons are: i) A deficit of electron-type solar neutrinos; ii) A deficit of muon-type atmospheric neutrinos; and, possibly, iii) The observation of the apearance of $`\overline{\nu }_e`$ in a beam of $`\overline{\nu }_\mu `$ by the LSND Collaboration. The invisible width of the $`Z`$ implies that the number of massless, or light Dirac, or light Majorana neutrino species is $`N_\nu =2.993\pm 0.011`$. To account for these observations we consider the minimal extension of the Standard Model with three massive neutrinos that mix. The neutrino interaction eigenstates $`\nu _l`$ are superpositions of the neutrino mass eigenstates $`\nu _m`$:
$$|\nu _l=\underset{m}{}U_{lm}|\nu _m$$
(1)
We consider the “standard” parametrization of the unitary matrix $`U_{lm}`$:
$$\left(\begin{array}{c}\nu _e\\ \nu _\mu \\ \nu _\tau \end{array}\right)=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta }\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }& c_{12}s_{23}s_{12}c_{23}s_{13}e^{i\delta }& c_{23}c_{13}\end{array}\right)\left(\begin{array}{c}\nu _1\\ \nu _2\\ \nu _3\end{array}\right)$$
(2)
where $`c_{ij}cos\theta _{ij}`$, $`s_{ij}sin\theta _{ij}`$, $`0\theta _{ij}\frac{\pi }{2}`$ and $`\pi \delta <\pi `$. The probability that an ultrarelativistic neutrino produced as $`\nu _l`$ decays as $`\nu _l^{}`$ is:
$$P(\nu _l\nu _l^{})=|\underset{m}{}U_{lm}exp(iLM_m^2/2E)U_{l^{}m}^{}|^2=P(\overline{\nu }_l^{}\overline{\nu }_l)$$
(3)
where $`E`$ and $`L`$ are the energy and traveling distance of $`\nu _l`$, and $`M_m`$ is the mass of $`\nu _m`$. We choose $`M_1M_2M_3`$. This extension of the Standard Model introduces six parameters: $`s_{12}`$, $`s_{23}`$, $`s_{13}`$, $`\delta `$, and two mass-squared differences, e.g. $`\mathrm{\Delta }M_{21}^2M_2^2M_1^2`$ and $`\mathrm{\Delta }M_{32}^2M_3^2M_2^2`$. We vary these parameters to minimize a $`\chi ^2`$. This $`\chi ^2`$ has $`14`$ terms obtained from the solar neutrino data summarized in Table 1, the atmospheric neutrino data shown in Table 2, and the LSND data: $`P(\overline{\nu }_\mu \overline{\nu }_e)=\frac{1}{2}\mathrm{sin}^2(2\theta )=0.0031\pm 0.0013`$ for $`L`$\[km\]$`/E`$\[GeV\]$`=[P(\overline{\nu }_\mu \overline{\nu }_e)]^{1/2}/1.27\mathrm{\Delta }M^2`$\[eV<sup>2</sup>\]$`0.73`$ (here $`\mathrm{sin}^2(2\theta )`$ corresponds to “large” $`\mathrm{\Delta }M^2`$, and $`\mathrm{\Delta }M^2`$ corresponds to $`\mathrm{sin}^2(2\theta )=1`$, see discussion in ). Because one author of the LSND Collaboration is in disagreement with the conclusion, and because the result has not been confirmed by an independent experiment, we multiply the error by $`1.5`$ and take $`P(\overline{\nu }_\mu \overline{\nu }_e)=0.0031\pm 0.0020`$. We require that the astrophysical, reactor and accelerator limits be satisfied. The most stringent of these limits are listed in Table 3.
The $`\chi ^2`$ has $`8`$ degrees of freedom ($`14`$ terms minus $`6`$ parameters). Varying the parameters we obtain minimums of $`\chi ^2`$, a few of which are listed in Table 4. With $`90\%`$ confidence the neutrino mass-squared differences lie within the dots shown in Figure 1. Note that one of the mass-squared differences is determined by the solar neutrino experiments and the other one by the atmospheric neutrino observations.
If neutrinos have a hierarchy of masses (as the charged leptons, up quarks, or down quarks), then the “upper island” in Figure 1 applies, and $`M_30.07`$eV, $`M_210^5`$eV and $`M_1<M_2`$, with large uncertainties.
Note in Table 1 that the ratio of the observed-to-predicted solar neutrino flux is significantly lower for the Homestake experiment than for Sage, Gallex, Kamiokande and Super-Kamiokande which are all compatible with $`0.5`$. These latter experiments observe neutrinos within wide energy bands, while the chlorine detector in the Homestake mine observes monochromatic electron-type neutrinos from a $`{}_{}{}^{7}Be`$ line. For the Homestake experiment the spread in $`L/E`$ is due to the spread in $`L`$, which in turn is due to the excentricity of the orbit of the Earth. Therefore the interference is coherent for up to $`L\mathrm{\Delta }M^2/(2E2\pi )30`$ oscillations from the Sun to the Earth (here $`\mathrm{\Delta }M^2`$ is either $`\mathrm{\Delta }M_{21}^2`$ or $`\mathrm{\Delta }M_{32}^2`$). Due to this coherence at “small” $`\mathrm{\Delta }M^2`$ it is possible to find acceptable solutions with $`\chi ^2<13.4`$ as shown in Figure 1. For larger values of $`\mathrm{\Delta }M^2`$ coherence is lost and we find solutions with $`\chi ^2>18`$ which are unacceptable if the Homestake experimental and theoretical errors are correct.
An important test of the model would be to observe seasonal variations of the neutrino flux of the $`{}_{}{}^{7}Be`$ line. If the lower ratio measured by the Homestake experiment is real, we expect that the electron-type neutrino flux of the $`{}_{}{}^{7}Be`$ line is near a minimum of the oscillation at the average Sun-Earth distance. In other words, there are an odd number of half-wavelengths from Sun to Earth. Then we expect a modulation of the $`{}_{}{}^{7}Be`$ neutrino flux with a period of half a year, with maximums occurring at the perihelion and aphelion of the Earth orbit. We see no statistically significant Fourier component of the time dependent Homestake data from 1970.281 to 1994.388. In particular the amplitude relative to the mean of a Fourier component of period $`0.5`$ years is $`0.09\pm 0.10`$. This observation implies that there are $`8.5`$ periods of oscillation from Sun to Earth at $`90\%`$ confidence level. With a $`\chi ^2`$ with 116 degrees of freedom, including the 8 discussed earlier plus the 108 measurements by the Homestake Collaboration from 1970.281 to 1994.388 we obtain the allowed region shown in Figure 2.
The reliability of $`M_2^2M_1^2`$ depends on the correctness of the error assigned to the Homestake observed-to-predicted flux ratio. For example, if the Homestake error listed in Table 1 is doubled we obtain the solutions shown in Figure 3.
In view of the preceeding results let us assume that neutrinos indeed have mass. The question then arizes wether neutrinos are distinct from antineutrinos (Dirac neutrinos) or wether neutrinos are their own antiparticles (Majorana neutrinos). This latter possibility arizes because neutrinos have no electric charge.
Let us consider Big-Bang nucleosynthesis that determines the abundances of the light elements $`D`$, $`{}_{}{}^{3}He`$, $`{}_{}{}^{4}He`$ and $`{}_{}{}^{7}Li`$. These abundances are determined by the temperatures of freezout $`T_f1MeV`$ when the reaction rates $`T_f^5`$ become comparable to the expansion rate $`T_f^2\times (5.5+\frac{7}{4}N_\nu )^{1/2}`$. Here $`N_\nu `$ is the equivalent number of massless neutrino flavors that are ultrarelativistic at $`T_f`$ and are still in thermal equilibrium with photons and electrons at that temperature. The calculated abundances of the light elements are in agreement with observations if $`1.6N_\nu 4.0`$ at $`95\%`$ confidence level. For three generations of Majorana neutrinos, $`N_\nu =3`$. For three generations of Dirac neutrinos, $`N_\nu =6`$ while in thermal equilibrium. However, in the Standard Model only the left-handed component of neutrinos couple to $`Z`$, $`W^+`$ and $`W^{}`$. Right-handed neutrinos are not in thermal equilibrium at $`T_f`$: their temperature has lagged below the temperature of photons due to the anihilation of particle-antiparticle pairs after the decoupling of the right-handed neutrinos. Therefore for Dirac neutrinos at $`T_f`$ we have $`N_\nu 3`$. So we can not distinguish Dirac from Majorana neutrinos using available data on nucleosynthesis.
In conclusion, the minimal extension of the Standard Model with three massive Majorana or Dirac neutrinos that mix is in good agreement with all experimental constraints. However, confirmation of the model is needed, e.g. by the observation of seasonal variations of the $`{}_{}{}^{7}Be`$ spectral lines with a period of $`0.5`$ years, or spectral distortions and seasonal variations of the low energy neutrinos from the solar $`pp`$ reaction. |
warning/0002/astro-ph0002519.html | ar5iv | text | # Optical and X-ray analysis of the cluster of galaxies Abell 496. Based on observations collected at the European Southern Observatory, La Silla, Chile
## 1 Introduction
In the framework of hierarchical clustering, the Universe is believed to be made of galaxies distributed in sheets encircling voids or filaments, at the intersection of which clusters of galaxies are located. Such models can be tested through the analysis of clusters, which are likely to keep a “memory” of their formation. This is suggested for example by the alignment effects observed in some clusters, such as for example Abell 3558 (Dantas et al. 1997) or Abell 85 (Durret et al. 1998), where the cD, the brightest galaxies, the X-ray emitting gas and possibly even larger scale structures (in the case of Abell 85) all appear aligned along the same direction. Multi-wavelength studies of clusters of galaxies also allow us to draw a global and coherent portrait of these objects, which we can then use to address other questions of interest, such as the influence of mergers and environmental effects at various scales on the properties of both galaxies and X-ray gas. Large scale (i.e. cluster size) mergers are quite often observed from substructures detected in the X-ray gas; smaller scale mergers (i.e. group size) such as the infall of dwarf galaxies onto groups surrounding bright galaxies can be derived from various methods such as those of Serna & Gerbal (1996) or Gurzadyan & Mazure (1998), which require optical velocity and magnitude catalogues; the existence of subclustering also has an influence on the shape of the galaxy luminosity function, which in some cases appears to show a deficit of faint galaxies often interpreted as due to accretion of dwarf galaxies onto larger galaxies or groups (e.g. in Coma, Lobo et al. 1997, Adami et al. 2000). It therefore appears important to analyze cluster properties in detail before using them in other studies. Note in particular that the existence of substructures can lead to overestimate cluster velocity dispersions, and hence M/L ratios and the value of $`\mathrm{\Omega }_0`$ in clusters.
With the improvement of both observational means (better X-ray detectors, optical multi-object spectroscopy) and modern methods of analysis (some of which are described below), an ever increasing number of clusters showing evidence for merging and environmental effects has been found. A rather general picture has therefore emerged for clusters, with a main relaxed body on to which groups of various sizes can be falling.
Our approach these last years has been to study a small sample of nearby clusters in detail. These have the advantage of being bright, and can therefore be observed in detail within a reasonable amount of telescope time. Besides, they are free of evolution effects. We present here a detailed multi-wavelength study of Abell 496, based on optical (extensive redshift and photometric catalogues) and X-ray (ROSAT PSPC) data.
Abell 496 is a richness class 1 (Abell 1958) cD type (Struble & Rood 1987) cluster at a redshift of 0.0331. For a Hubble constant H<sub>0</sub>=50 km s<sup>-1</sup> Mpc<sup>-1</sup>, the corresponding scale is 55.0 kpc/arcmin and the distance modulus is 36.52. At optical wavelengths, an adaptive kernel map of the central region (in a 60$`\times `$60 arcmin<sup>2</sup> square) has revealed a somewhat complicated structure, with a strong concentration of galaxies in the north-south direction (Kriessler & Beers 1997). Note however that this map does not include redshift information. Thorough investigations of the X-ray properties of Abell 496 can be found in Mohr et al. (1999) and Markevitch et al. (1999); their results will be compared to ours in section 4.5.
The paper is organized as follows: we present the data in section 2; the structures along the line of sight derived from the velocity catalogue are described in section 3; the optical properties of the Abell 496 cluster itself are described in section 4; the X-ray cluster properties are described in section 5; a summary and conclusions are given in section 6.
## 2 The data
### 2.1 Optical data
Our redshift catalogue includes 466 galaxies in the direction of the cluster Abell 496, in a region covering about 160$`\times `$160 arcmin<sup>2</sup> (9.2$`\times `$9.2 Mpc for an average redshift for Abell 496 of 0.0331). It is fully described in Durret et al. (1999b).
Our photometric catalogues are described in Slezak et al. (1999). The photographic plate catalogue was obtained by scanning an SRC plate in the b<sub>J</sub> band with the MAMA measuring machine at the Observatoire de Paris; it includes 3879 galaxies located in a region of roughly $`\pm `$ 1.3 from the cluster centre. Positions are very accurate in this catalogue and were used for spectroscopy; on the other hand, b<sub>J</sub> magnitudes are not accurate, so a CCD photometric catalogue was obtained in the V and R bands in order to recalibrate photographic plate magnitudes. The R magnitudes thus estimated for the photographic plate catalogue were used to estimate the completeness of our redshift catalogue. The CCD imaging catalogue includes 239 and 610 galaxies in the V and R bands respectively, and is limited to a much smaller region of $`246`$ arcmin<sup>2</sup> in the centre of the cluster.
The cluster center will be taken to be the position of the cD galaxy, which coincides with the X-ray maximum: $`\alpha =4^h33^{mn}38.45^s,\delta =13^{}15^{}49.5^{\prime \prime }`$ (2000.0).
### 2.2 X-ray data
The ROSAT PSPC image (#800024) was taken from the archive and analyzed by Pislar (1998). The cluster was observed during 8972 s. The effective exposure time for this image, after data reduction is 5354 s. We have used the routines developed by Snowden (1995) to obtain a non cosmic background subtracted image between 0.5 and 2 keV. We have defined the image limiting radius ($`\mathrm{\hspace{0.17em}1620}h_{50}^1`$ kpc) as the radius where the surface brightness reaches the surface brightness of the cosmic background ($`\mathrm{3\hspace{0.17em}10}^4\mathrm{s}^1\mathrm{arcmin}^2`$). The different PSPC background components are detailed in Snowden et al. (1994).
The global X-ray gas temperatures derived from Einstein and EXOSAT satellite data are 3.9$`\pm `$0.2 keV (David et al. 1993) and $`4.7_{0.8}^{+1}`$ keV respectively (Edge & Stewart 1991).
## 3 Velocity distribution along the line of sight
We first discuss the overall properties derived from the velocity distribution along the line of sight.
### 3.1 Overall characteristics of the structures detected along the line of sight
A wavelet reconstruction of the velocity distribution along the line of sight is displayed in Fig. 1 (466 galaxies). We remind the reader that this type of reconstruction takes into account structures at a significance level of at least 3$`\sigma `$, and detects structures at various scales. The sample was analyzed with 256 points, and the smallest scale was excluded because it is very noisy. A more detailed description of this technique can be found in Fadda et al. (1998).
Nine “groups” or velocity substructures are found with this method, and their velocity characteristics are given in Table 1. The group number is given in col. 1, the number of galaxies in col. 2, the BWT mean velocity (Beers et al. 1990) and corresponding BWT velocity dispersion in cols. 3 and 4, and the velocity interval in col. 5. Foreground groups 1 and 2 and background groups 7 and 8 are most probably not real groups, since they are widely spread both on the sky and in velocity distribution; because of the small number of galaxies involved in the first three of these groups, we did not calculate mean velocities or velocity dispersions for these structures. For group 8 these values are only indicative, but characteristic of a low mass structure. Group 3 appears to be the cluster Abell 496 itself. Except for 7 objects, all the galaxies in group 4 appear to be located north of Abell 496. Group 5 has the same kind of shape and size as group 4 and is roughly coaligned with Abell 496 along the line of sight. Group 6 appears strongly concentrated both spatially and in velocity space, all but two galaxies being located west of Abell 496. Moreover, its velocity dispersion is also low.
These results are confirmed when we apply the same method as for the ENACS clusters (Katgert et al. 1996, Mazure et al. 1996) to detect velocity structures along the line of sight. This method consists in sorting the galaxies in order of increasing velocity, and plotting their rank as a function of velocity (hereafter the rank-velocity classification). If the distribution of galaxies in redshift space is strictly gaussian, we expect to see a regular S-shape in the sequence/gap space. When there are more than 5 galaxies between two successive gaps, we consider that the galaxies belong to a structure.
### 3.2 A finer analysis of structures 4, 5, 6 and 9
In order to understand whether groups 4, 5, 6 and 9 can be physically coherent structures, we performed a Serna & Gerbal (1996) analysis for each of these groups separately. Since this type of analysis takes into account galaxy magnitudes, we had to eliminate one galaxy in group 4 and one in group 5, for which we have no magnitude in our redshift catalogue. We also tried keeping these galaxies and assigning them an “average” magnitude R=17. The results in both cases were similar. Note that the redshift catalogue completeness is about 50% in regions 4 and 5, and 55% in region 9, within the magnitude limit R=18.8. It could therefore be argued that the Serna & Gerbal method is meaningless for these samples. However, in this type of analysis it is the brightest galaxies which mainly contribute to the dynamics of the system (since the mass to luminosity ratio is taken to be constant for all galaxies). If the samples are limited to magnitudes R$``$17.0, the redshift catalogue completeness then becomes 72%, 72% and 79% for regions 4, 5 and 9 respectively, and the Serna & Gerbal method is therefore fully applicable.
The characteristics of the substructures found with the Serna & Gerbal method are given in Table 2. Structure 4 has subgroups 4c and 4d well defined in space; they extend over 5 and 3.6 Mpc respectively and could therefore be members of two different clusters. Due to their large spatial extent, subgroups 4a and 4b respectively seem to be just forward and background galaxies, with the exception of the five galaxies of group 4a at the north west extremity (see Fig. 2).
Subgroups 5a and 5b form structures with a small velocity dispersion, but extending over about 7 Mpc, a size which appears rather large for these groups to be members of two background clusters; the extent and the velocity dispersion of 5c are even larger (Fig. 2).
While the Serna & Gerbal (1996) method finds dynamical sub-structures for the other groups along the line of sight, the same method reveals no substructures in group 6, except for two pairs of galaxies; group 6 therefore appears well defined both in velocity distribution and in space. Eighteen galaxies are included in an ellipse with a major and minor axes of about 8 and 3 Mpc, suggesting that this is a poor, diffuse and low mass cluster (Fig. 2).
Finally, three subgroups are apparently found in structure 9, 9a and 9b having very small velocity dispersions but spanning a rather large spatial region. The overall velocity field in group 9 shows an interesting pattern looking like a velocity gradient (Fig. 3). This could be a filament more or less aligned along the line of sight.
A rank-velocity classification applied to each group confirms that group 4 and possibly group 9 appear to have three substructures (two clear breaks in the curves), group 5 has three or four substructures and group 6 has no clear substructure except perhaps for the two or three first galaxies which are probably infalling objects.
To summarize, groups 4 and 9 clearly appear as filaments or at least elongated structures along the line of sight, but not really massive clusters. This analysis is confirmed by the iso-velocity contours of group 9. The continuous velocity gradient could be interpreted as the result of a merger, with the infalling groups not perfectly aligned along the line of sight. Group 6 has a low velocity dispersion and is probably a poor cluster. Group 5 exhibits two low velocity dispersion sub-structures and a moderately high velocity dispersion group (5c), but the number of galaxies in 5c is too low to provide a robust estimation and we assume that this structure is not a cluster.
We will now discuss the dynamical state of the main structure on the line of sight: the cluster Abell 496 itself (group 3).
## 4 Morphological and physical properties of Abell 496
### 4.1 Morphology of the cluster at various wavelengths
We display in Fig. 4 the superposition of the optical image of the cluster with ROSAT PSPC X-ray and radio isocontours. The X-ray contours are quite smooth, with no obvious substructures. However, there appears to be an excess of X-ray emission towards the north west, in the direction where there is also an excess of emission line galaxies (see below). A radio source is visible south east of the cluster, probably associated with a galaxy.
### 4.2 The galaxy velocity distribution in Abell 496
The cluster Abell 496 corresponds to structure 3 in Table 1; it has a BWT mean velocity of 9885 km/s and a global velocity dispersion of 715 km/s. The corresponding velocity interval is \[7813,11860 km/s\] and includes 274 galaxies. Note that the cD galaxy has a velocity of 9831 km/s, close to the mean cluster velocity, and is located very close the X-ray emission center, suggesting that the cD is at the bottom of the cluster gravitational potential well. This is an indication of a quiescent history of the cluster (see e.g. Zabludoff et al. 1993, Oegerle & Hill 1994).
The wavelet reconstruction of the velocity distribution of Abell 496 shown in Fig. 5 (274 galaxies) suggests the presence of a certain amount of substructure. The sample was analyzed with 256 points, and the two smallest scales were excluded. The corresponding velocity distribution is non-gaussian; it shows: a tiny feature at $``$8300 km/s; a main asymmetric peak in the \[8500, 10700 km/s\] range containing 232 galaxies, with a BWT mean velocity of 9769 km/s, and a BWT velocity dispersion of 518 km/s; note that this velocity structure is not quite centered on the velocity of the cD galaxy; a smaller peak at 10940 km/s with 36 galaxies in the \[10700, 11860 km/s\] range. If we only keep the largest scales, we are left with a rather symmetric velocity distribution showing an excess at high velocities. This excess corresponds to the peak at 10940 km/s, which contains a small number of galaxies (see section 4.4).
These structures are also found by applying a rank-velocity classification, which gives two breaks globally consistent with those found by analyzing the cluster velocity distribution. Such breaks probably indicate substructures with velocities coherent with the finer analysis based on the wavelet technique. However, the number of galaxies involved in these structures is small, and the velocity distribution in the main cluster therefore appears to be quite smooth, suggesting that Abell 496 is quite well relaxed.
In order to confirm the state of relaxation of Abell 496, we have applied a Serna & Gerbal (1996) analysis to the subsample of 96 galaxies located within a radius of 1800 arcsec around the cD and with magnitudes R$``$17.0; within this limited sample, the completeness of the redshift catalogue is 82% and this type of analysis is expected to give robust results. Note that galaxies in this region with measured velocities but without published magnitudes were discarded. Results are displayed in Fig. 6. At the extreme lower right of the figure, we can see the very tight pair made by the cD (#280, R=12.2) and a satellite galaxy (#292, R=15.6, in the Durret et al. 1999b catalogue): this confirms that the cD is at the bottom of the cluster potential well. We also observe a structure of 11 bright galaxies (10 galaxies with R$``$15.8, plus one with R=16.8) highly concentrated in space around the cD (mean distance to the cD: 216 arcsec, with a dispersion of 48 arcsec) but not in velocity (BWT mean velocity and velocity dispersion: 9745 and 375 km/s). This result is comparable to what is found in other clusters, where the central core is more or less well discriminated. The main body of the cluster center appears quite relaxed, with no strong subclustering within a radius of 1800 arcsec (1.65 Mpc). This picture agrees with the general shape of Abell 496 seen in X-rays (see Fig. 4).
### 4.3 Luminosity segregation in Abell 496
After violent relaxation, two-body gravitational interactions lead to a certain level of energy equipartition between galaxies of various masses, and consequently to a certain segregation in velocity dispersion with luminosity (mass). This process concerns essentially massive galaxies, and added to dynamical friction, it creates segregation with distance to the cluster center. The stage of post-violent relaxation therefore leads to a segregation in the \[L,$`\sigma _v`$\] space larger in the central regions than in the overall cluster.
In order to search for such effects in Abell 496, we have derived the velocity dispersion and the average distance of galaxies to the cluster center (defined by the position of the cD) in several magnitude bins. We restrict our sample to galaxies belonging to the cluster, i.e. in the velocity range \[7813,11860 km/s\], and within 1000 and 1800 arcsec from the cluster center, in order to have reasonably complete samples:100% and 79% complete for R$`18.5`$ respectively. The completeness is estimated by comparing the number of galaxies with measured redshifts to the number of galaxies in our photographic plate catalogue, for the same R magnitude limit. Since there are 2 galaxies with R$`<14`$ and 6 with 14$`<`$R$`<`$15, we chose to fit the data with two different “brightest” bins: one including the 8 galaxies with 12$`<`$R$`<`$15, and the other including only the 6 galaxies with 14$`<`$R$`<`$15.
The velocity dispersions estimated in several magnitude bins are different for the two samples, as shown in Fig. 7. In a 1000 arcsec radius, the velocity dispersion increases more steeply with magnitude: the corresponding slopes are $`65\pm 43`$ and $`70\pm 50`$ km s<sup>-1</sup> mag<sup>-1</sup> when the brightest bin is included or not respectively. In a 1800 arcsec radius, the velocity dispersion increases less with magnitude, the corresponding slopes being $`52\pm 30`$ and $`32\pm 41`$ km s<sup>-1</sup> mag<sup>-1</sup>.
As seen in Fig. 8, the average distance to the cluster center is somewhat smaller for the galaxies located in the brightest bin (R$``$15), then remains roughly constant with a possible decrease with increasing magnitude, specially in the broadest sample (1800 arcsec radius).
The combination of Figs. 7 and 8 seems to correspond well to a post-violent relaxation stage.
Interestingly, we can notice that both the distance to the cluster center and the overall velocity dispersion range are reduced when emission line galaxies (hereafter ELGs) are excluded. This agrees with the general scheme that ELGs are often found in the outskirts of clusters of clusters and are not as strongly tied to the cluster gravitationally (e.g. Biviano et al. 1997). We now discuss in more detail the properties of ELGs in Abell 496.
### 4.4 The emission line galaxy distribution
We now compare the distribution of emission line (ELGs) versus non-emission line (NoELGs) galaxies. There are 85 ELGs and 381 NoELGs in our velocity catalogue, among which 34 ELGs and 241 NoELGs in the velocity range of Abell 496. The global percentage of ELGs in the cluster is therefore $`12\pm 3`$%. Note that this percentage is perfectly coherent with the proportions observed by Biviano et al. (1997) in the ENACS survey.
The spatial distribution of the 211 NoELGs and 21 ELGs belonging to the main \[8500, 10700 km/s\] velocity peak is displayed in Fig. 9. The fraction of ELGs in this velocity range is 9$`\pm `$3%, consistent with the global cluster value within the error bars. NoELGs appear rather homogeneously distributed, except for a sort of linear north-south concentration towards the center. On the other hand, a large majority of the ELGs in this velocity range is located west of a north to south line crossing the center, and at least half of these ELGs even seem to be close to the west cluster edge. This agrees with the fact that ELGs tend to concentrate in the outer regions of clusters (see e.g. results based on ENACS data by Biviano et al. 1997). The presence of an excess of ELGs can at least in some cases be interpreted as due to merging events producing shocks which trigger star formation. This was shown to be the case in the zone of Abell 85 where the X-ray filament merges into the main body of the cluster (Durret et al. 1998): an excess of ELGs was observed in that region, together with a temperature increase of the X-ray gas. The ASCA X-ray gas temperature map presently available for Abell 496 (Markevitch et al. 1999, Donnelly, private communication), does not show any temperature increase for the X-ray gas in that area, so we cannot correlate the excess of ELGs towards the west cluster edge with a higher gas temperature zone. However, we can note that this excess is located roughly in the same region as the X-ray excess emission in the north west region of the cluster (Fig. 4). Such an X-ray enhancement could be due to a merging event originating from the north west, but our data cannot show this with certainty.
On the other hand, the spatial distributions of the 23 NoELGs and 13 ELGs in the \[10700, 11860 km/s\] velocity interval are comparable (Fig. 9), while the ELG fraction seems much higher: 36$`\pm `$18%. Though the small number of objects may introduce errors, there definitely seems to be an excess of ELGs with somewhat higher velocities than the bulk of the cluster; these ELGs account at least partly for the peak at 10940 km s<sup>-1</sup> in the wavelet reconstruction of the velocity distribution.
A general picture for the ELG distribution in Abell 496 is that of two samples of galaxies falling on to the main cluster, one from the back (the ELGs concentrated towards the west) and one from the front (the high velocity ELGs).
### 4.5 The galaxy luminosity function
We have seen in the previous section that Abell 496 appears to have properties common to many clusters, with a relaxed main body and ELGs probably falling on to the cluster. We therefore expect its galaxy luminosity function not to be strongly modified by environmental effects, as observed in some clusters showing more prominent substructures. We discuss below its main features.
#### 4.5.1 The bright end of the galaxy luminosity function
We have first derived the galaxy luminosity function (GLF) of Abell 496 in the R band from the redshift catalogue, within a radius of 1800 arcsec around the center, and for a limiting magnitude R=18.5. There are 196 galaxies in this sample. The completeness of the redshift catalogue within these limits is 79%, and it is 100% in that region for R$``$16.0. The obvious interest of such a GLF is that no background contribution needs to be subtracted, therefore making the results very robust. We have limited the magnitude interval to the \[13,18.5\] range, because for R$``$13 there is only one galaxy (the cD), which introduces edge effects in the wavelet reconstruction of the GLF, and for R$``$18.5 the completeness sharply decreases. This corresponds to the \[$`23.5,18.0`$\] interval in R absolute magnitude.
The GLF obtained after a wavelet reconstruction is shown in Fig. 10. The sample was analyzed with 128 points, excluding the two smallest scales. The significance level of the detected features is at least 3$`\sigma `$. A flattening is observed at R$``$16, corresponding to an absolute magnitude M$`{}_{\mathrm{R}}{}^{}20.5`$. This shape is comparable to that found in Virgo (Binggeli et al. 1988) and in Abell 963 (Driver et al. 1994), where a comparable flattening was observed at a common absolute magnitude of $`19.8`$. The GLFs of Coma and Abell 85 are more complex, with a “bump” corresponding to the brightest galaxies, followed by a “dip” at M$`{}_{\mathrm{R}}{}^{}20.5`$ (see Fig. 9 in Durret et al. 1999a). Note that the flattening of the GLF in Abell 496 occurs exactly at the same absolute magnitude as the dip in Abell 85 and Coma.
The GLF in Abell 496 suggests at least a bimodal galaxy distribution, with bright (mostly elliptical) galaxies in the bright part and dwarf galaxies in the fainter part. We therefore performed a fit of the wavelet reconstructed GLF of Abell 496 by summing two functions: a gaussian, to account for bright galaxies, and a power law (case 1) or a Schechter function (case 2) to represent faint and/or dwarf galaxies. In case 1, we fit the data as a function of apparent R magnitude; the gaussian is then found to be centered on R=15.19$`\pm `$0.01, with $`\sigma =0.8\pm 0.1`$, and the power law varies as R<sup>11.67±0.16</sup>. In case 2, we fit the data as a function of absolute R magnitude, to allow a direct comparison with other authors; the gaussian is then found to be a little broader, centered on R=15.48$`\pm `$0.04, with $`\sigma =1.0\pm 0.1`$; the Schechter function, defined as in Rauzy et al. (1988, section 3.2.3), has $`\alpha =1.19\pm 0.04`$ and M$`{}_{}{}^{}=19.43\pm 0.13`$ (in the \[$`23.5,18.0`$\] absolute magnitude range).
The GLF resulting from these various fits is shown in Fig. 10; it obviously reproduces the data very well. Except at the faint end where the sample incompleteness most probably modifies the GLF shape, both fits 1 and 2 are good but we cannot distinguish between them. In view of the obvious quality of the fit, we did not attempt to estimate error bars with Monte-Carlo simulations, as done previously e.g. for Abell 85 (see Durret et al. 1999a, Fig. 12).
The various values obtained from these fits of the GLF can be compared to those found in other clusters. The gaussian used to fit the bright part of the GLF in Abell 85 has $`\sigma =1.1`$, comparable to the value we find in case 1. The Schechter function for Abell 496 has a slope comparable to that found by Lumsden et al. (1997), but notably flatter than the values found in other surveys (e.g. Valotto et al. 1997, Rauzy et al. 1998 and references therein). Rauzy et al. (1998) argued that the flatter slope found by Lumsden was due to incompleteness at faint magnitudes; this may also be true in our case, since our sample is 100% complete only to R=16.0 (M$`{}_{\mathrm{R}}{}^{}=20.5`$), and we also find a brighter value of M than the above surveys, suggesting that we are missing faint galaxies.
We can note that the GLFs of Coma and Abell 85 were interpreted in a similar way, with the bright part mainly due to ellipticals (with a small contribution of spirals) and the faint part due to dwarfs (Durret et al. 1999a). Comparable shapes were found in several other clusters. The fact that the GLF of Abell 496 shows a flattening at the same value M$`{}_{\mathrm{R}}{}^{}=20.5`$ indicates that the galaxy population in Abell 496 is comparable to those of the above mentioned clusters.
Note that Molinari et al. (1998) have analyzed the GLF of Abell 496 from a photometric catalogue in three colors, reaching magnitudes much fainter than those of our spectroscopic catalogue. We will therefore compare our results to theirs in the next section.
#### 4.5.2 The faint end of the galaxy luminosity function
Our intent was also to derive the luminosity function from the CCD catalogue, which corresponds to a small region of $``$246 arcmin<sup>2</sup> in the cluster center. For this we first performed a wavelet reconstruction of the R magnitude distribution in the R magnitude range . Since we have no background exposure, we estimated the background contribution by connecting the counts from the Las Campanas Redshift Survey (LCRS, Lin et al. 1996) and from the ESO-Sculptor Survey (ESS, Arnouts et al. 1997), as described in our study of Abell 85 (Durret et al. 1999a, Fig. 10 and text), and we subtracted this background to the observed number of galaxies. The result is shown in Fig. 11. We have checked that the consistency of the background number counts estimated by Tyson (1988) with those of the LCRS and ESS combined as described above is good.
The difference between the observed number of galaxies and the background (Fig. 11) becomes negative for magnitudes R$`18.4`$, while the CCD catalogue is complete at least up to R=21. Therefore, this background cannot be considered as representative of the local background in our CCD field of view. Note that this was already the case for the CCD photometric data of Abell 85.
One notable feature is the dip in the galaxy magnitude distribution at R$`19.5`$ (M$`{}_{\mathrm{R}}{}^{}17`$), which is detected at a high confidence level. This dip corresponds to that observed by Molinari et al. (1998), who found a dip at R$``$19 (M$`{}_{\mathrm{R}}{}^{}17.5`$). Note that they also find a similar dip in the g band, and possibly in the i band. Molinari et al. (1998) made a second determination of the GLF by selecting cluster members in a colour-magnitude diagram. In this case, they find a small dip, or at least a flattening, for R$`18`$ (M$`{}_{\mathrm{R}}{}^{}18.5`$). This value does not agree either with the bright nor with the faint GLF that we derived. It is difficult to understand why, since their colour-magnitude relation appears quite well defined.
In order to investigate the origin of the dip seen in our data, we propose a toy model, which is not a fit but only illustrates how the dip could be accounted for. Let us first note that the contribution of the other structures detected along the line of sight is negligible. Assuming a Gaussian + a Schechter function to model the GLF (see section 4.5.1), we rescaled the number of galaxies produced by this composite function to fit the dimension of the CCD field. We then applied a magnitude cut-off to this GLF, as suggested by Adami et al. (2000), for galaxies fainter than M$`{}_{\mathrm{R}}{}^{}=19.75`$ in the inner core of the Coma cluster. This effect becomes very strong for galaxies fainter than M$`{}_{\mathrm{R}}{}^{}=17`$. The exact shape of such a cut-off is unknown, so we applied a convenient apodization function (the choice of this function influences the shape and smoothness of the dip). The background counts were modeled as the background contribution from the LCRS and ESS described above. We then summed the cluster and background contributions, and the result is shown in Fig. 12.
Such a toy model can reproduce the global GLF shape, with counts similar to the observed data and a dip comparable to the observed one. A fine-tuning of the various parameters involved could make Figs. 11 and 12 more similar, but this would push the model too far. However, we can state that a cut-off in the GLF of Abell 496 similar to that observed in Coma is a solution to account for the observed dip.
## 5 The X-ray gas
A pixel by pixel fit was performed on the X-ray image, as described by Pislar et al. (1997). The pixel size is 30 arcsec. A $`\beta `$-model and a 3D Sérsic model (Lima Neto et al. 1999) were considered for the variations of the density with radius. The global temperature estimated from these ROSAT data, using a Raymond-Smith spectrum and a Galactic absorption column density was found to be 4$`\pm `$1 keV and assumed to be constant (Pislar 1998). This is consistent with the temperatures of 3.9 and 4.7 keV previously measured with the Einstein and EXOSAT satellites respectively (David et al. 1993; Edge & Stewart 1991). The parameters corresponding to the best fits for both models are given in Table 3, and the result of the $`\beta `$-model 1 fit superimposed on the observed image is displayed in Fig. 13. We observe that in model 2 the central density is lower than in model 1 and that the $`\beta `$ and r<sub>c</sub> parameters are higher. This is because in model 2 we do not include the central region, where the cooling flow lies. The effect is the same for models 3 and 4.
Our values of $`\beta `$ and r<sub>c</sub> (in model 2) are higher than those of Markevitch et al. (1999), who found $`\beta =0.7`$ and r$`{}_{c}{}^{}=249`$ kpc. This is due to the fact that they exclude a central region smaller than ours (3 arcmin instead of 3.3 arcmin). Pislar (1998) has shown that in a cooling flow cluster the bigger the excluded central region, the higher $`\beta `$ and r<sub>c</sub>, and the lower the central density. Moreover, at the image limiting radius, the values of the dynamical and gas masses do not depend on the size of the central region excluded.
The cooling radius R<sub>cool</sub> and the mass $`\dot{M}`$ deposited in the centre were estimated as in Pislar et al. (1997) and are given in Table 3 for a temperature of 4 keV. We take $`\mathrm{2\hspace{0.17em}10}^{10}`$ yr for the cooling time.
The X-ray gas and total dynamical masses derived from the X-ray data assuming equilibrium are shown in Fig. 14 as a function of radius. Note that these curves are only valid between the cooling radius and the limiting radius of the X-ray image; within this validity range, both models are in good agreement with each other.
We find, with models 2 and 4 for a prolate geometry, at the limiting radius of the image, a gas mass of $`(6.1\pm 2.2)10^{13}M_{}`$ and a dynamical mass of $`(4.2\pm 1.1)\mathrm{\hspace{0.17em}10}^{14}M_{}`$. We overestimate the errors because we have supposed that the parameters are not correlated. The masses at $`1h_{50}^1`$ Mpc are respectively $`(3.45\pm 1.1)\mathrm{\hspace{0.17em}10}^{13}M_{}`$ and $`(2.4\pm 0.6)\mathrm{\hspace{0.17em}10}^{14}M_{}`$. The gas mass found by Mohr et al. (1999) and the dynamical mass found by Markevitch et al. (1999) are very similar if we remember that their geometry is spherical. The mass calculated by the virial theorem applied to all the galaxies in the redshift catalogue with velocities in the cluster is M$`{}_{vir}{}^{}=(7.2\pm 0.8)10^{14}`$ M. This value agrees with the dynamical mass derived above, providing the mass profiles can be extrapolated to radii larger than the X-ray image. The integrated mass of galaxies within the X-ray image limiting radius, assuming a mass to luminosity ratio of 8 M/L, is $`\mathrm{1.0\; 10}^{13}`$ M. Note that, due to incompleteness of our redshift catalogue, in particular at faint magnitudes, this is only a lower limit.
The ratio of the X-ray gas mass to the dynamical mass is shown as a function of radius in Fig. 15. The fraction of X-ray gas is about 0.15 at the limiting radius of the image and at $`1h_{50}^1`$ Mpc with a $`\beta `$-model. With a Sérsic model, the ratio is 0.12 at the limiting radius of the image. The baryon fraction that we find at $`1h_{50}^1`$ Mpc with the $`\beta `$-model is very similar to that obtained by Markevitch et al. (1999) ($`0.158\pm 0.017`$), and comparable to that found in other clusters. At the X-ray limiting radius, the stellar to X-ray gas mass ratio is 16% and the stellar to total dynamical mass ratio is 2.4%.
## 6 Discussion and conclusions
The optical analysis of the Abell 496 field has shown the existence of several structures along the line of sight. Among these, one (structure 6) is likely to be a poor, diffuse and low mass cluster, while two others (structures 4 and 9) are probably filaments more or less aligned along the line of sight, the latter presenting a smooth velocity gradient. Notice that the distances between these various structures are very large.
The cluster Abell 496 itself has quite a regular morphology. It includes 274 galaxies in the \[7813,11860 km/s\] velocity range and has a velocity dispersion of 715 km/s. Its velocity distribution implies a small amount of substructure. The analysis of the correlations between position, luminosity and velocity dispersion indicates a post-violent relaxation state. We can notice that both the distance to the cluster center and the overall velocity dispersion ranges are reduced when emission line galaxies (hereafter ELGs) are excluded. This agrees with the general scheme that ELGs are often found in the outskirts of clusters and are not as strongly tied to the cluster gravitationally (e.g. Biviano et al. 1997). There may be two samples of ELGs falling on to the main cluster, one from the back (the ELGs concentrated towards the west) and one from the front (the high velocity ELGs).
The bright luminosity function derived from our redshift catalogue shows a flattening at R$`16`$ (M$`{}_{\mathrm{R}}{}^{}20.5`$), comparable to similar shapes found in other clusters. This suggests at least a bimodal distribution, one for ellipticals and one for fainter galaxies. The fact that the flattening occurs at the same absolute magnitude as for other clusters suggests that the galaxy populations in all these clusters are comparable.
At fainter magnitudes, galaxy counts derived from CCD imaging show a dip at R$`19.5`$ (M$`{}_{\mathrm{R}}{}^{}17`$) which can be reproduced if we assume a magnitude cut-off similar to that observed in Coma (Adami et al. 2000). Notice that such a cut-off is observed in the very central regions of both clusters. Although this result is only based on imaging and remains to be confirmed spectroscopically, we may be evidencing a second example of a cut-off of the faint end of the luminosity function in a cluster.
We have modelled the X-ray gas and derived the X-ray gas mass and the dynamical mass, which we compare to the stellar mass. At the limiting radius (1.62 h$`{}_{}{}^{1}{}_{50}{}^{}`$ Mpc) of the image, we find a fraction of X-ray gas to total mass of 0.12$``$0.15 and a stellar to X-ray gas mass ratio of 0.16. We can note that Abell 496 follows exactly the two by two correlations between the X-ray luminosity (L$`{}_{\mathrm{X}}{}^{}=\mathrm{6.8\; 10}^{44}`$ erg/s, Wu et al. 1999), the X-ray temperature (T$`{}_{\mathrm{X}}{}^{}=4`$ keV) and the galaxy velocity dispersion ($`\sigma _v=715`$ km/s) described in the literature (see e.g. Wu et al. 1999 and references therein). These values are typical of a richness class 1 cluster.
Abell 496 therefore appears to be a relatively quiet and simple cluster, with no strong environmental effects, although we may see an enhancement of the X-ray emission and of the number of emission line galaxies towards the north west.
While Coma has long been the archetype of a relaxed cluster and is not believed to be relaxed any more (see Biviano 1998 and the proceedings of the Coma meeting), the results presented above suggest that Abell 496 may be such a prototype, and can be used as a “template” in the future study of more complex (i.e. substructured) clusters.
###### Acknowledgements.
The authors thank Andrea Biviano for help. C. Adami is grateful to the staff of the Dearborn Observatory for their hospitality during his postdoctoral fellowship. We acknowledge financial support from the French Programme National de Cosmologie, CNRS. |
warning/0002/hep-ph0002300.html | ar5iv | text | # References
Non-saturation of the J/$`\psi `$ suppression at large
transverse energy in the comovers approach
A. Capella and E. G. Ferreiro
Laboratoire de Physique Théorique<sup>*</sup><sup>*</sup>*Unité Mixte de Recherche - CNRS - UMR N 8627
Université de Paris-Sud, Bâtiment 210, F-91405 Orsay Cedex, France
A. B. Kaidalov
ITEP, B. Cheremushkinskaya, 25, 117259 Moscow, Russia
## Abstract
We show that, contrary to recent claims, the $`J/\psi `$ suppression resulting from its interaction with comovers does not saturate at large transverse energy $`E_T`$. On the contrary, it shows a characteristic structure - change of curvature near the knee of the $`E_T`$ distribution - which is due to the $`E_T`$ (or multiplicity) fluctuation, and agrees with recent experimental results.
LPTHE Orsay 00-21
February 2000
An interesting result of the last run (1998 data) of the NA50 collaboration at CERN on the transverse energy ($`E_T`$) dependence of $`J/\psi `$ suppression in $`PbPb`$ collisions, is the observation of a convexity at large $`E_T`$. More precisely, for $`E_T>`$ 100 GeV (which corresponds to the so-called knee of the $`E_T`$ distribution ), the slope of the ratio $`R(E_T)`$ of $`J/\psi `$ over Drell-Yan (DY) cross-sections increases with increasing $`E_T`$. In sharp contrast with this result, models of the $`J/\psi `$ suppression in non-quark-gluon plasma (QGP) scenarios \[2–9\] \- such as the one based on the interaction of the $`J/\psi `$ with comovers - exhibit a clear saturation at large $`E_T`$.
In this work, we show that the above feature of the comovers model is only true up to the knee of the $`E_T`$ distribution ($`E_T`$ 100 GeV). Beyond this value, we enter into the tail of the $`E_T`$ distribution - where the increase in $`E_T`$ is due to fluctuations. This fluctuation, which has not been taken into account in most calculations, produces a corresponding increase in the density of comovers - which, in turn, increases the $`J/\psi `$ suppression at large $`E_T`$.
In order to illustrate this phenomenon we use the model introduced in ref. . Here, as in most non-QGP models, the $`J/\psi `$ suppression is due to two mechanisms : absorption of the pre-resonant $`c\overline{c}`$ pair with nucleons (the so-called nuclear absorption) and the interaction of the $`J/\psi `$ with comovers. The corresponding $`J/\psi `$ survival probabilities are given by .
$$S^{abs}(b,s)=\frac{\left\{1\mathrm{exp}[AT_A(s)\sigma _{abs}]\right\}\left\{1\mathrm{exp}\left[BT_B(bs)\sigma _{abs}\right]\right\}}{\sigma _{abs}^2ABT_A(s)T_B(bs)},$$
(1)
$$S^{co}(b,s)=\mathrm{exp}\left\{\sigma _{co}N_y^{co}(b,s)\mathrm{ln}\left(\frac{N_y^{co}(b,s)}{N_f}\right)\right\}.$$
(2)
The survival probability $`S^{co}`$ depends on the density of comovers $`N_y^{co}(b,s)`$ in the rapidity region of the dimuon trigger $`2.9<y_{lab}<3.9`$ and $`N_f=1.15`$ fm<sup>-2</sup> is the corresponding density in $`pp`$ collisions. In order to compute $`N_y^{co}`$, various hadronic models have been used in the literature. For instance in ref. it has been assumed that the hadronic multiplicity is proportional to the number of participant nucleons (the so-called wounded nucleon model), while in ref. a formula based on the dual parton model (DPM) was used - which includes an extra term proportional to the number of binary interactions. In this paper we use the DPM formula (eq. (6) of ). In both cases, the calculations do not include the fluctuations mentioned above and, therefore, cannot be applied beyond the knee of the $`E_T`$ distribution - where the increase in $`E_T`$ (or multiplicity) is entirely due to fluctuations. In order to introduce these fluctuations, it is convenient to recall the other formulae needed to calculate the $`J/\psi `$ suppression.
At fixed impact parameter $`b`$, the $`J/\psi `$ cross-section is given by
$$\sigma _{AB}^\psi (b)=\frac{\sigma _{pp}^\psi }{\sigma _{pp}}d^2sm(b,s)S^{abs}(b,s)S^{co}(b,s),$$
(3)
where $`m(b,s)=AB\sigma _{pp}T_A(s)T_B(bs)`$. The corresponding one for DY pair production is obtained from (3) putting $`\sigma _{abs}=\sigma _{co}=0`$ (i.e. $`S^{abs}=S^{co}=1)`$ and is proportional to $`AB`$. In this way we can compute the ratio of $`J/\psi `$ over DY as a function of the impact parameter. Experimentally, however, the ratio $`R(E_T)`$, is given as a function of the transverse energy $`E_T`$ measured by a calorimeter, in the rapidity interval $`1.1<y_{lab}<2.3`$. In order to compute $`R(E_T)`$ we have to know the correlation $`P(E_T,b)`$ between $`E_T`$ and impact parameter, which is given by
$$P(E_T,b)=\frac{1}{\sqrt{2\pi qaE_T^{NF}(b)}}\mathrm{exp}[\frac{E_TE_T^{NF}(b)}{2qaE_T^{NF}(b)}]^2.$$
(4)
Here
$$E_T^{NF}(b)=qN_{cal}^{co}(b)+k[Am_A(b)]E_{in},$$
(5)
$`m_A(b)`$ is the number of participants of $`A`$ (at fixed impact parameter), $`E_{in}=158`$ GeV/c is the beam energy and $`k=1/4000`$ . In (4) and (5) $`N_{cal}^{co}(b)`$ is obtained by integrating the comover density $`N_y^{co}(b,s)`$ over $`d^2s`$, and $`dy`$ (in the rapidity range of the $`E_T`$ calorimeter). The second term in (5) was introduced in ref. in order to reproduce the correlation between $`E_T`$ and the energy $`E_{ZDC}`$ of the zero degree calorimeter. It was interpreted as due to intra-nuclear cascade - which is present here due to the location in rapidity of the NA50 calorimeter. This term is sizable for peripheral collisions, when many spectator nucleons are present, and dies away for central ones. The parameters $`q=0.56`$ and $`a=0.94`$ are obtained from a fit to the minimum bias $`E_T`$ distribution at large $`E_T`$. The parameter $`q`$ gives the relation between multiplicity of comovers (positive, negative and neutrals) and the $`E_T`$ of the NA50 calorimeter (which contains only neutrals). The product $`qa`$ controls the width of the $`E_T`$ distribution at fixed $`b0`$. The $`J/\psi `$ and DY cross-section at fixed $`E_T`$ are then given by
$$\frac{d\sigma ^{\psi (DY)}}{dE_T}=d^2b\sigma _{AB}^{\psi (DY)}P(E_T,b).$$
(6)
The quantity $`E_T^{NF}(b)`$ in eq. (5) does not contain fluctuations - hence the index $`NF`$. This is obvious from the fact that the parameter $`a`$ is not present in (5). In order to see it in a more explicit way, we plot in Fig. 1 the quantity
$$F(E_T)=E_T/E_T^{NF}(E_T),$$
(7)
where
$$E_T^{NF}(E_T)=\frac{d^2bE_T^{NF}(b)P(E_T,b)}{d^2bP(E_T,b)}.$$
(8)
We see that $`E_T^{NF}`$ coincides with $`E_T`$ only up to the knee of the $`E_T`$ distribution. Beyond it, $`E_T^{NF}`$ is smaller than the true value of $`E_T`$. This difference is precisely due to the $`E_T`$ fluctuation.
As discussed above, in order to compute the ratio $`R(E_T)`$ beyond the knee of the $`E_T`$ distribution it is necessary to introduce in $`N_y^{co}`$ the $`E_T`$ (or multiplicity) fluctuations responsible for the tail of the distribution. In order to do so, we use the experimental observation that multiplity and $`E_T`$ distributions have approximately the same shape. This indicates that the fluctuations in $`E_T`$ are mainly due to fluctuations in multiplicity - rather than in $`p_T`$. This leads to the following replacement in eq. (2):
$$N_y^{co}(b,s)N_y^{co}(b,s)F(E_T).$$
(9)
In this way the results for the ratio $`R(E_T)`$ are unchanged below the knee of the distribution (see Fig. 1). Beyond it, the $`J/\psi `$ suppression is increased as a result of the fluctuation.
We turn next to the numerical results. In ref. we used for the two parameters of the model $`\sigma _{abs}=6.7`$ mb and $`\sigma _{co}=0.6`$ mb. In this case, the computed $`J/\psi `$ suppression at $`E_T100`$ GeV is somewhat too small . Clearly, we can increase it by increasing the value of $`\sigma _{co}`$. However, we then increase the value of the suppression for peripheral collisions. This, in turn, can lead to some conflict with the SU data (see for a discussion on this point). However, recent data on the $`J/\psi `$ cross-section in $`pA`$ collisions, point to a smaller value of $`\sigma _{abs}`$ \- of 4 to 5 mb. With this value of $`\sigma _{abs}`$, we can increase $`\sigma _{co}`$ from 0.6 mb up to 1.0 mb without decreasing the $`J/\psi `$ suppression for peripheral collisions. In Fig. 2 we present the result of our calculation using $`\sigma _{abs}=4.5`$ mb and $`\sigma _{co}=1`$ mb. We see that the main features of the data are reproduced. In particular our curve shows a slight change of curvature at $`E_T`$ 100 GeV, which is entirely due to the effect of fluctuations - and is seen in the 1998 NA50 data . The physical origin of this change in the slope of $`E_T`$ is the following: when approaching the knee of the $`E_T`$ distribution from below, the number of participants approaches $`2A`$ and changes slowly. The latter is also true for the multiplicity of comovers. Beyond the knee, the multiplicity increases faster due to the fluctuations and produces a faster decrease of $`R(E_T)`$ (see Fig. 1).
We want to stress that the shape of our curve in the lower half of the $`E_T`$ region (where the ratio $`R(E_T)`$ changes rather fast with $`E_T`$) is sensitive to the relation between $`E_T`$ and impact parameter. We see from eq. (5) that this relation depends on the size of the contribution of the intra-nuclear cascade (parameter $`k`$). As mentioned above, the value $`k=1/4000`$ used here was obtained in from the best fit of the correlation between $`E_T`$ and the energy $`E_{ZDC}`$ of the zero degree calorimeter. However, since we do not have a totally reliable expression for the latter, there is an uncertainty in the value of $`k`$. In order to illustrate its effect on $`R(E_T)`$, we show in Fig. 2 (dashed line) the result with $`k=1/2000`$, i.e. doubling the (comparatively small) contribution of the intra-nuclear cascade. The effect is concentrated in the lower half of the $`E_T`$ range. This uncertainty would not be present if the $`J/\psi `$ suppression were given as a function of either $`E_T`$ or charged multiplicity at mid-rapidities.
The large $`E_T`$ structure seen by the NA50 collaboration can also be explained assuming a deconfining phase transition . At the energies of the Relativistic Heavy Ion Collider (RHIC) at Brookhaven, it will be possible to determine which of these two mechanisms is the correct one. Indeed, the transverse energy (or the corresponding energy density) where the structure has been seen by NA50, will be reached at RHIC well below the knee of the $`E_T`$ distribution, and, if our interpretation is correct, no structure will be present. It will, however, appear at higher values of $`E_T`$ \- when reaching the knee of the $`E_T`$ distribution at $`\sqrt{s}=200`$ GeV.
ACKNOWLEGMENTS
It is a pleasure to thank N. Armesto, C. Pajares, C. Gerschel, C. A. Salgado and Y. M. Shabelski for interesting discussions. We also thank B. Chaurand and M. Gonin for providing numerical tables of the NA50 data. E. G. F. thanks Ministerio de Educacion y Cultura of Spain for financial support.
FIGURE CAPTIONS
FIG. 1. The ratio $`F(E_T)`$ in eqs. (7), (8).
FIG. 2. The ratio $`R(E_T)`$ of $`J/\psi `$ over DY cross-sections, obtained with $`\sigma _{abs}=4.5`$ mb and $`\sigma _{co}=1`$ mb, compared to the NA50 data . The full curve corresponds to $`k=1/4000`$ in eq. (5). The dashed curve is obtained with $`k=1/2000`$ (see main text). The black points correspond to 1996 Pb-Pb data, the black squares correspond to 1998 Pb-Pb data, the white points to 1996 analysis with minimum bias and the white squares to 1998 analysis with minimum bias.
1 Capella
2 Capella |
warning/0002/hep-ph0002240.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Nowdays it is established beyond any doubts that the naive picture of light hadrons as made of three constituent quarks (for baryons) or $`q\overline{q}`$ pairs of constituent quarks (for mesons) is not complete. The DIS experiments revealed the rich sea structure of the nucleon, these experiments showed in particular that a considerable portion of the nucleon spin is carried by the strange component of the nucleon sea. Furthemore there are experimental facts which seems to suggest that a non-vanishing nonperturbative component of intrinsic charm is present in light hadrons .
We address the problem of intrinsic charm content of light hadrons from the point of view of the heavy quark mass expansion. The $`c\overline{c}`$ pairs in light hadrons, due to parametrically large mass of charm quarks, can appear in a light hadron as virtual state whose life time is short, of order $`1/m_c`$. The nonperturbative (with typical momenta below heavy quark mass $`m_c`$) gluon and light quark fluctuations are slowly varying from “point view” of virtual $`c\overline{c}`$ pair, hence the heavy quark mass expansion is equivalent to the semiclassical expansion. This expansion allows to rewrite operators made of heavy quarks in terms of light degrees of freedom (gluons and light quarks). For a detailed discussion of the heavy quark mass expansion see ref. .
Let us note also that in absence of a direct probe of gluons the open charm production is considered as the main source of information about nucleon’s gluon distributions. In hard leptoproduction heavy quarks are produced in the leading order via the photon-gluon fusion (PGF). The leading graph for PGF can be related directly to gluon distributions if one assumes that there is no intrinsic charm content in the nucleon (no $`c(x),\overline{c}(x)`$ and no $`\mathrm{\Delta }c(x),\mathrm{\Delta }\overline{c}(x)`$ at normalzation point $`\mu =m_c`$). However now there are many evidences that, in principle, there might be considerable intrinsic charm component in the nucleon wave function even at low normalization point. For reliable extraction of gluon distributions from open charm electroproduction experiments it is necessary to have quantitative estimates of the intrinsic charm content of the nucleon.
This paper will be organized as follows: In the first part we present the calculation of the expectation value of heavy quark currents in the background of gluon fields using a semiclassical approximation. This corresponds to an expansion in the inverse of the heavy quark mass
$$Q^{}(x)\mathrm{\Gamma }Q(x)=\underset{n}{}\frac{1}{m^n}X_\mathrm{\Gamma }^{(n)},$$
(1)
where $`\mathrm{\Gamma }`$ denotes the Lorentz structure of the current and the $`X_\mathrm{\Gamma }^{(n)}`$ are local expressions of the field strength depending on $`\mathrm{\Gamma }`$. In section 2.1 we review the large m expansion of the fermion determinant appearing in our definition of the expectation value. In section 2.2 we then outline the expansion of color singlet currents in general before we present in section 2.2.1 the expansion of the axial current used in in full detail. The connection to the expectation value of the axial vector current using the axial anomaly equation and some general restriction coming from this equation are given in section 2.2.2. As further examples we present the expansion of the scalar current in section 2.2.3, the vector current (section 2.2.4) and the tensor current (section 2.2.5), respectively. In section 2.2.6 we finally show the result of the expansion of $`Q^{}(x)_\mu \gamma _\nu Q(x)`$, appearing in the energy-momentum tensor of QCD.
In the second part we discuss the calculation of intrinsic heavy quark content of light hadrons as an application of the heavy quark mass expansion. In the case of charm content of $`\eta ^{},\eta `$ mesons and intrinsic charm contributions to the proton spin we reduce the calculations of these quantities to matrix elements which are already known either phenomenologically or were computed previously. In other cases, like intrinsic charm contribution to the nucleon tensor charge and to energy momentum tensor, the problem is reduced to matrix elements of gluon operators which can be estimated using various nonperturbative methods in QCD: lattice calculation, QCD sum rule, theory of instanton vacuum.
## 2 Heavy quark expansion of currents in the background gluon and light quark fields
The expectation value of a color-singlet quark current made of heavy quarks in the background of gluon and light quark fields can be written after integration out heavy degrees of freedom as:
$`Q^{}(x)\mathrm{\Gamma }Q(x)`$ $`=`$ $`\mathrm{det}D\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\mathrm{\Gamma }|x,`$ (2)
Here $`\mathrm{\Gamma }`$ denotes an arbitrary Lorentz-structure. Note that all calculations will be performed in the euclidean space-time, so the QCD Dirac operator reads:
$$D=\mathrm{i}\overline{)}+\mathrm{i}m$$
(3)
where the covariant derivative is defined as
$$_\mu =\left(_\mu \mathrm{i}\frac{\lambda ^a}{2}A_\mu ^a(x)\right)$$
(4)
and $`m`$ is the heavy quark mass. For the used conventions and the euclidization see the appendix. Eq.(2) can now be expanded in a power-series of the inverse heavy quark mass, $`1/m`$, under the assumption that the gradient of the background field strength is small compared to $`m`$. The expansion of determinant of the Dirac operator in eq.(2) has been calculated by a large number of authors, see e.g. . We briefly review the calculation of the determinant following since we use the result of this calculation as a check of the expansion of a scalar current of heavy quarks in Section 2.2.3.
### 2.1 Expansion of the determinant
The expansion of the determinant for heavy quarks in eq.(2) yields divergences of various types. Since most of these divergences are connected with the determinant of the free Dirac operator we normalize the determinant with that in zero external field. For the remaining infinity which can be related to the logarithmic renormalization of the coupling constant, we use the so-called $`\zeta `$-regularization. Using the identity
$$\mathrm{det}D=\mathrm{det}D^{}=\left(\mathrm{det}D^{}D\right)^{\frac{1}{2}}$$
(5)
the normalized and regularized determinant can be written as follows:
$$\left(\frac{\mathrm{det}D}{\mathrm{det}D_0}\right)_{\zeta \mathrm{reg}}=\mathrm{exp}\left[\frac{1}{2}\underset{s0}{lim}\frac{\mathrm{d}}{\mathrm{d}s}\frac{M^{2s}}{\mathrm{\Gamma }(s)}_0^{\mathrm{}}dtt^{s1}\mathrm{Tr}\left[e^{tD^{}D}e^{tD_0^{}D_0}\right]\right],$$
(6)
where $`D_0`$ denotes the Dirac operator in the absence of external gluon fields and $`M`$ is the regulator mass. The functional trace denoted by Tr in eq.(6) can be calculated with respect to any complete set of states. For further calculations it is convinient to compute functional traces in the basis of plane waves, so that
$`\mathrm{Tr}\left[e^{tD^{}D}\right]`$ $`=`$ $`\mathrm{tr}_{c,\gamma }{\displaystyle \mathrm{d}^4x\frac{\mathrm{d}^4k}{\left(2\pi \right)^4}e^{\mathrm{i}kx}\left[e^{tD^{}D}\right]e^{\mathrm{i}kx}}`$ (7)
$`=`$ $`\mathrm{tr}_{c,\gamma }{\displaystyle \mathrm{d}^4x\frac{\mathrm{d}^4k}{\left(2\pi \right)^4}\left[e^{tD^{}(_\mu _\mu +\mathrm{i}k_\mu )D(_\nu _\nu +\mathrm{i}k_\nu )}\right]1}.`$
The unity in eq.(7) points out that the operators here act on unity, so that $`_\mu 1=0`$. The further calculations are straightforward: the expression in eq.(7) can be expanded in powers of the covariant derivative, integrated with respect to $`k`$ and the Lorentz indices summed.
Since the explicit calculation is given in we present here only some useful formulae and the final result for the determinant up to order $`𝒪(1/m^2)`$. The square of the Dirac operator in eq.(7) with all differentiation operators shifted, $`_\mu _\mu +\mathrm{i}k_\mu `$, is given by
$$D^{}(_\mu _\mu +\mathrm{i}k_\mu )D(_\nu _\nu +\mathrm{i}k_\nu )=^2+\frac{\sigma }{2}F2\mathrm{i}k+k^2+m^2,$$
(8)
where we have used that
$`F_{\mu \nu }^a`$ $`=\mathrm{i}[_\mu ,_\nu ]^a=_\mu A_\nu ^a_\nu A_\mu ^a+f^{abc}A_\mu ^bA_\nu ^c`$ (9)
$``$ $`\overline{)}\overline{)}`$ $`=^2+{\displaystyle \frac{\sigma }{2}}F,`$ (10)
with the notations $`\sigma F=\sigma _{\mu \nu }\frac{\lambda ^a}{2}F_{\mu \nu }^a`$ and $`\sigma _{\mu \nu }=\frac{\mathrm{i}}{2}[\gamma _\mu ,\gamma _\nu ]`$ applied. Expanding the exponential function in eq.(7) the functional trace then reads
$`\mathrm{Tr}\left[e^{tD^{}D}\right]=\mathrm{tr}_{c,\gamma }{\displaystyle \mathrm{d}^4x\frac{\mathrm{d}^4k}{\left(2\pi \right)^4}e^{t(k^2+m^2)}\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{t^n}{n!}\left(^2+\frac{\sigma }{2}F2\mathrm{i}k\right)^n1}`$ (11)
$`={\displaystyle \frac{1}{4\pi ^2}}\mathrm{tr}_c{\displaystyle }\mathrm{d}^4xe^{tm^2}[{\displaystyle \frac{1}{t^2}}+t^0({\displaystyle \frac{1}{6}}_\alpha _\beta _\alpha _\beta {\displaystyle \frac{1}{6}}_\alpha ^2_\alpha +{\displaystyle \frac{1}{4}}F_{\alpha \beta }F_{\alpha \beta })`$
$`+t({\displaystyle \frac{1}{180}}^2^2^2{\displaystyle \frac{1}{36}}(_\alpha ^2_\alpha ^2+_\alpha ^2^2_\alpha +^2_\alpha ^2_\alpha )`$
$`+{\displaystyle \frac{1}{45}}(_\alpha _\beta _\alpha _\beta ^2+_\alpha _\beta _\alpha ^2_\beta +_\alpha _\beta ^2_\alpha _\beta `$
$`+_\alpha ^2_\beta _\alpha _\beta +^2_\alpha _\beta _\alpha _\beta +_\alpha _\beta ^2_\beta _\alpha )`$
$`{\displaystyle \frac{1}{90}}(_\alpha _\beta _\alpha _\gamma _\beta _\gamma +_\alpha _\beta _\gamma _\alpha _\beta _\gamma `$
$`+_\alpha _\beta _\gamma _\alpha _\gamma _\beta +_\alpha _\beta _\gamma _\beta _\alpha _\gamma +_\alpha _\beta _\gamma _\beta _\gamma _\alpha )`$
$`+{\displaystyle \frac{1}{6}}(F_{\alpha \beta }F_{\alpha \beta }^2+F_{\alpha \beta }^2F_{\alpha \beta }+^2F_{\alpha \beta }F_{\alpha \beta }`$
$`_\gamma F_{\alpha \beta }_\gamma F_{\alpha \beta }_\gamma F_{\alpha \beta }F_{\alpha \beta }_\gamma F_{\alpha \beta }_\gamma F_{\alpha \beta }_\gamma )`$
$`{\displaystyle \frac{\mathrm{i}}{6}}F_{\alpha \beta }F_{\gamma \beta }F_{\gamma \alpha })+𝒪(^8)].`$ (12)
Here we used that all contributions with an odd number of $`k`$’s vanish whereas all other integrals with respect to $`k`$ yield
$$\frac{\mathrm{d}^4k}{\left(2\pi \right)^4}k_{\mu _1}\mathrm{}k_{\mu _{2n}}e^{t(k^2+m^2)}=\frac{1}{4\pi ^2}\left(2t\right)^{(n+2)}e^{tm^2}\delta _{\mu _1\mathrm{}\mu _{2n}},$$
(13)
with $`\delta _{\mu _1\mathrm{}\mu _{2n}}`$ denoting all possible contractions:
$`\delta _{\mu _1\mathrm{}\mu _{2n}}=\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{^2}{\varphi _\nu \varphi _\nu }}\right]\varphi _{\mu _1}\mathrm{}\varphi _{\mu _{2n}}|_{\varphi =0}.`$ (14)
After rearranging the terms into gauge invariants and taking also the part without external gluon fields in eq.(6) into account, the determinant up to order $`1/m^2`$ can be written as follows:
$`\left({\displaystyle \frac{\mathrm{det}D}{\mathrm{det}D_0}}\right)_{\zeta \mathrm{reg}}`$ $`=`$ $`\mathrm{exp}[{\displaystyle }\mathrm{d}^4x({\displaystyle \frac{1}{48\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{M^2}{m^2}}\right)\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$ (16)
$`{\displaystyle \frac{\mathrm{i}}{720\pi ^2}}{\displaystyle \frac{1}{m^2}}\mathrm{tr}_cF_{\alpha \beta }F_{\beta \gamma }F_{\gamma \alpha }+{\displaystyle \frac{1}{1440\pi ^2}}{\displaystyle \frac{1}{m^2}}\mathrm{tr}_c[_\alpha ,F_{\alpha \beta }][_\gamma ,F_{\gamma \beta }]`$
$`{\displaystyle \frac{11}{1440\pi ^2}}{\displaystyle \frac{1}{m^2}}\mathrm{tr}_c[_\gamma ,[_\alpha ,F_{\alpha \beta }]]F_{\gamma \beta }+{\displaystyle \frac{1}{360\pi ^2}}{\displaystyle \frac{1}{m^2}}_\alpha \mathrm{tr}_c[_\alpha ,F_{\gamma \beta }]F_{\gamma \beta }`$
$`{\displaystyle \frac{1}{384\pi ^2}}{\displaystyle \frac{1}{m^2}}^2\mathrm{tr}_cF_{\gamma \beta }F_{\gamma \beta })+𝒪\left({\displaystyle \frac{1}{m^4}}\right)]`$
$`=\mathrm{exp}[{\displaystyle }\mathrm{d}^4x({\displaystyle \frac{1}{48\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{M^2}{m^2}}\right)\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$
$`{\displaystyle \frac{\mathrm{i}}{720\pi ^2}}{\displaystyle \frac{1}{m^2}}\mathrm{tr}_cF_{\alpha \beta }F_{\beta \gamma }F_{\gamma \alpha }`$
$`+{\displaystyle \frac{1}{120\pi ^2}}{\displaystyle \frac{1}{m^2}}\mathrm{tr}_c[_\alpha ,F_{\alpha \beta }][_\gamma ,F_{\gamma \beta }])+𝒪\left({\displaystyle \frac{1}{m^4}}\right)].`$
Note that in the last step partial integration has been used with all total derivatives left out. The effective action $`S_{\mathrm{eff},\mathrm{E}}=\mathrm{ln}\mathrm{det}D`$ which can be yielded from (16), rotated back to Minkowski space, corresponds exactly to the result of :
$`S_{\mathrm{eff},\mathrm{M}}`$ $`=`$ $`{\displaystyle \frac{1}{48\pi ^2}}{\displaystyle }\mathrm{d}^4x(\mathrm{ln}\left({\displaystyle \frac{M^2}{m^2}}\right)\mathrm{tr}_cF_{\alpha \beta }F^{\alpha \beta }`$ (17)
$`{\displaystyle \frac{\mathrm{i}}{15\pi ^2}}{\displaystyle \frac{1}{m^2}}\mathrm{tr}_cF_{\alpha \beta }F_{}^{\beta }{}_{\gamma }{}^{}F^{\gamma \alpha })+𝒪({\displaystyle \frac{1}{m^4}}),`$
where equation of motion terms, which vanish in pure Yang-Mills theory $`[_\alpha ,F_{\alpha \beta }]=0`$ have been neglected.
### 2.2 Expansion of heavy quark currents
In order to expand $`\mathrm{tr}_{c,\gamma }x|\frac{1}{D}\mathrm{\Gamma }|x`$ in eq.(2) in a series of the inverse heavy quark mass we can use eq.(8) to rewrite it as:
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\mathrm{\Gamma }|x`$ $`=`$ $`\mathrm{tr}_{c,\gamma }{\displaystyle \frac{\mathrm{d}^4k}{\left(2\pi \right)^4}e^{\mathrm{i}kx}\frac{1}{D^{}D}D^{}\mathrm{\Gamma }e^{\mathrm{i}kx}}`$ (18)
$`=\mathrm{tr}_{c,\gamma }{\displaystyle \frac{\mathrm{d}^4k}{\left(2\pi \right)^4}\frac{1}{k^2+m^2}\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{^2\frac{\sigma }{2}F+2\mathrm{i}k}{k^2+m^2}\right)^n\left(\mathrm{i}\overline{)}\overline{)}k\mathrm{i}m\right)\mathrm{\Gamma }1}.`$
The expansion in eq.(18) is again justified for small gradients of the gluonic fields compared to the heavy quark mass $`m`$. Depending on the Lorentz structure $`\mathrm{\Gamma }`$ some of the integrals might be divergent and need to be regularized, we choose the dimensional regularization, since the integrals in eq.(18) can then be calculated using:
$$\frac{\mathrm{d}^dk}{\left(2\pi \right)^d}\frac{k_{\mu _1}k_{\mu _2}\mathrm{}k_{\mu _{2m}}}{\left(k^2+m^2\right)^n}=\frac{1}{\left(4\pi \right)^{d/2}}\frac{\mathrm{\Gamma }\left(nmd/2\right)}{\mathrm{\Gamma }\left(n\right)2^m}\delta _{\mu _1\mathrm{}\mu _{2m}}\left(\frac{1}{m^2}\right)^{nmd/2}.$$
(19)
The number of terms contributing to a given order of $`1/m`$ is reduced by the fact that terms containing an odd number of $`\gamma `$ matrices or an odd number of $`k`$’s vanish due to the trace over Lorentz-indices and the integration with respect to $`k`$. The expansion of eq.(18) then is straightforward. The result of the expansion must be gauge invariant because we expand the gauge invariant operator. In order to obtain explicitely gauge invariant result for heavy quark mass expansion a number of helpful identities based on the Bianchi identity:
$$[_\alpha ,F_{\beta \gamma }]+[_\gamma ,F_{\alpha \beta }]+[_\beta ,F_{\gamma \alpha }]=0,$$
(20)
can be derived. They will be presented in the following sections.
We want to illustrate some technical details of expanding heavy quark currents on the example of the pseudoscalar density and the divergency of the axial-vector current, which are related to each other by the axial anomaly. Another motivation to show detailed calculation for these cases is that recently confusing results for these cases were reported in the literature . Further we present the result of the expansion of scalar, vector and tensor currents and of $`\overline{Q}_\mu \gamma _\nu Q`$ appearing in the energy-momentum tensor of QCD.
#### 2.2.1 The pseudoscalar density
For $`\mathrm{\Gamma }=\gamma _5`$ the expansion eq.(18) has the form:
$$\mathrm{tr}_{c,\gamma }x|\frac{1}{D}\gamma _5|x=\mathrm{i}m\mathrm{tr}_{c,\gamma }\frac{\mathrm{d}^4k}{\left(2\pi \right)^4}\frac{1}{k^2+m^2}\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{^2\frac{\sigma }{2}F+2\mathrm{i}k}{k^2+m^2}\right)^n\gamma _51$$
(21)
Collecting all terms which contribute up to $`𝒪(1/m^3)`$ one gets:
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\gamma _5|x`$ $`=`$ $`\mathrm{i}m\mathrm{tr}_{c,\gamma }{\displaystyle }{\displaystyle \frac{\mathrm{d}^4k}{\left(2\pi \right)^4}}[{\displaystyle \frac{1}{\left(k^2+m^2\right)^3}}{\displaystyle \frac{1}{4}}\sigma F\sigma F\gamma _5`$ (23)
$`+{\displaystyle \frac{1}{\left(k^2+m^2\right)^4}}\left({\displaystyle \frac{1}{4}}^2\sigma F\sigma F+{\displaystyle \frac{1}{4}}\sigma F^2\sigma F+{\displaystyle \frac{1}{4}}\sigma F\sigma F^2{\displaystyle \frac{1}{8}}\sigma F\sigma F\sigma F\right)\gamma _5`$
$`{\displaystyle \frac{1}{\left(k^2+m^2\right)^5}}(\sigma F\sigma Fkk+\sigma Fk\sigma Fk+\sigma Fkk\sigma F`$
$`+k\sigma Fk\sigma F+kk\sigma F\sigma F+k\sigma F\sigma Fk)\gamma _5]+𝒪({\displaystyle \frac{1}{m^5}})`$
$`={\displaystyle \frac{\mathrm{i}}{32\pi ^2m}}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_cF_{\alpha \beta }F_{\gamma \delta }{\displaystyle \frac{1}{48\pi ^2m^3}}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_cF_{\rho \alpha }F_{\rho \beta }F_{\gamma \delta }`$
$`+{\displaystyle \frac{\mathrm{i}}{192\pi ^2m^3}}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c[F_{\alpha \beta }F_{\gamma \delta }^2F_{\alpha \beta }_\rho F_{\gamma \delta }_\rho +F_{\alpha \beta }^2F_{\gamma \delta }`$
$`_\rho F_{\alpha \beta }_\rho F_{\gamma \delta }+^2F_{\alpha \beta }F_{\gamma \delta }_\rho F_{\alpha \beta }F_{\gamma \delta }_\rho ]+𝒪({\displaystyle \frac{1}{m^5}}).`$
Here we have used eq.(19) for the integration over $`k`$ and the following results for Dirac traces:
$`F_{\alpha \beta }F_{\gamma \delta }\mathrm{tr}_\gamma \left[\sigma _{\alpha \beta }\sigma _{\gamma \delta }\gamma _5\right]`$ $`=`$ $`F_{\alpha \beta }F_{\gamma \delta }\mathrm{tr}_\gamma \left[\gamma _\alpha \gamma _\beta \gamma _\gamma \gamma _\delta \gamma _5\right]`$ (24)
$`=`$ $`4\epsilon _{\alpha \beta \gamma \delta }F_{\alpha \beta }F_{\gamma \delta },`$
$`F_{\alpha \beta }F_{\gamma \delta }F_{ϵ\phi }\mathrm{tr}_\gamma \left[\sigma _{\alpha \beta }\sigma _{\gamma \delta }\sigma _{ϵ\phi }\gamma _5\right]`$ $`=`$ $`\mathrm{i}F_{\alpha \beta }F_{\gamma \delta }F_{ϵ\phi }\mathrm{tr}_\gamma \left[\gamma _\alpha \gamma _\beta \gamma _\gamma \gamma _\delta \gamma _ϵ\gamma _\phi \gamma _5\right]`$ (25)
$`=`$ $`16\mathrm{i}\epsilon _{\alpha \beta \gamma \delta }F_{\rho \alpha }F_{\rho \beta }F_{\gamma \delta }`$
Using the following identities:
$`[_\rho ,F_{\alpha \beta }][_\rho ,F_{\gamma \delta }]`$ $`=`$ $`_\rho F_{\alpha \beta }_\rho F_{\gamma \delta }+F_{\alpha \beta }_\rho F_{\gamma \delta }_\rho _\rho F_{\alpha \beta }F_{\gamma \delta }_\rho F_{\alpha \beta }^2F_{\gamma \delta },`$ (26)
$`[_\rho ,[_\rho ,F_{\alpha \beta }]]F_{\gamma \delta }`$ $`=`$ $`^2F_{\alpha \beta }F_{\gamma \delta }+F_{\alpha \beta }^2F_{\gamma \delta }2_\rho F_{\alpha \beta }_\rho F_{\gamma \delta },`$ (27)
$`F_{\alpha \beta }[_\rho ,[_\rho ,F_{\gamma \delta }]]`$ $`=`$ $`F_{\alpha \beta }^2F_{\gamma \delta }+F_{\alpha \beta }F_{\gamma \delta }^22F_{\alpha \beta }_\rho F_{\gamma \delta }_\rho `$ (28)
eq.(23) can be written as
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\gamma _5|x`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{32\pi ^2m}}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_cF_{\alpha \beta }F_{\gamma \delta }`$ (29)
$`+{\displaystyle \frac{\mathrm{i}}{192\pi ^2m^3}}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c[2[_\rho ,[_\rho ,F_{\alpha \beta }]]F_{\gamma \delta }+[_\rho ,F_{\alpha \beta }][_\rho ,F_{\gamma \delta }]`$
$`+4\mathrm{i}F_{\rho \alpha }F_{\rho \beta }F_{\gamma \delta }]+𝒪({\displaystyle \frac{1}{m^5}}).`$
From the Bianchi identity (20) we obtain the following relations<sup>1</sup><sup>1</sup>1 In the calculation of the factor of 2 was missing in the first identity, which led to the wrong result.
$`[_\rho ,[_\rho ,F_{\alpha \beta }]]=2\mathrm{i}[F_{\rho \alpha },F_{\rho \beta }]+[_\alpha ,[_\rho ,F_{\rho \beta }]][_\beta ,[_\rho ,F_{\rho \alpha }]],`$ (30)
$`\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c[_\rho ,F_{\rho \beta }][_\alpha ,F_{\gamma \delta }]=0,`$ (31)
so that
$`\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c[_\rho ,[_\rho ,F_{\alpha \beta }]]F_{\gamma \delta }`$ $`=`$ $`2\mathrm{i}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c\left[F_{\rho \alpha }F_{\rho \beta }F_{\gamma \delta }F_{\rho \beta }F_{\rho \alpha }F_{\gamma \delta }\right]`$ (32)
$`+\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c\left[[_\alpha ,[_\rho ,F_{\rho \beta }]]F_{\gamma \delta }[_\beta ,[_\rho ,F_{\rho \alpha }]]F_{\gamma \delta }\right]`$
$`=\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c\left[4\mathrm{i}F_{\rho \alpha }F_{\rho \beta }F_{\gamma \delta }+2[_\alpha ,[_\rho ,F_{\rho \beta }]]F_{\gamma \delta }\right]`$
$`=\epsilon _{\alpha \beta \gamma \delta }\left(4\mathrm{i}\mathrm{tr}_cF_{\rho \alpha }F_{\rho \beta }F_{\gamma \delta }+2_\alpha \mathrm{tr}_c[_\rho ,F_{\rho \beta }]F_{\gamma \delta }\right).`$
On the other hands it yields
$`\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c\left[[_\rho ,[_\rho ,F_{\alpha \beta }]]F_{\gamma \delta }+[_\rho ,F_{\alpha \beta }][_\rho ,F_{\gamma \delta }]\right]`$ $`=`$ $`\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_c[_\rho ,[_\rho ,F_{\alpha \beta }]F_{\gamma \delta }]`$ (33)
$`=\epsilon _{\alpha \beta \gamma \delta }_\rho \mathrm{tr}_c[_\rho ,F_{\alpha \beta }]F_{\gamma \delta }`$
$`=\epsilon _{\alpha \beta \gamma \delta }_\rho \mathrm{tr}_c\left[[_\rho ,F_{\alpha \beta }F_{\gamma \delta }]F_{\alpha \beta }[_\rho ,F_{\gamma \delta }]\right]`$
$`={\displaystyle \frac{1}{2}}\epsilon _{\alpha \beta \gamma \delta }_\rho \mathrm{tr}_c[_\rho ,F_{\alpha \beta }F_{\gamma \delta }]`$
$`={\displaystyle \frac{1}{2}}\epsilon _{\alpha \beta \gamma \delta }^2\mathrm{tr}_cF_{\alpha \beta }F_{\gamma \delta }.`$
So we finally end up with
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\gamma _5|x`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{32\pi ^2m}}\epsilon _{\alpha \beta \gamma \delta }\mathrm{tr}_cF_{\alpha \beta }F_{\gamma \delta }`$ (34)
$`+{\displaystyle \frac{\mathrm{i}}{384\pi ^2m^3}}\epsilon _{\alpha \beta \gamma \delta }^2\mathrm{tr}_cF_{\alpha \beta }F_{\gamma \delta }+{\displaystyle \frac{\mathrm{i}}{96\pi ^2m^3}}\epsilon _{\alpha \beta \gamma \delta }_\alpha \mathrm{tr}_c[_\rho ,F_{\rho \beta }]F_{\gamma \delta }+𝒪\left({\displaystyle \frac{1}{m^5}}\right)`$
$`={\displaystyle \frac{\mathrm{i}}{16\pi ^2m}}\mathrm{tr}_cF_{\alpha \beta }\stackrel{~}{F}_{\alpha \beta }`$
$`+{\displaystyle \frac{\mathrm{i}}{192\pi ^2m^3}}^2\mathrm{tr}_cF_{\alpha \beta }\stackrel{~}{F}_{\alpha \beta }+{\displaystyle \frac{\mathrm{i}}{48\pi ^2m^3}}_\alpha \mathrm{tr}_c[_\rho ,F_{\rho \beta }]\stackrel{~}{F}_{\alpha \beta }+𝒪\left({\displaystyle \frac{1}{m^4}}\right),`$
where we have introduced the common notation $`\stackrel{~}{F}_{\alpha \beta }=\frac{1}{2}\epsilon _{\alpha \beta \gamma \delta }F_{\gamma \delta }`$.
#### 2.2.2 The divergency of the axial vector current
Instead of expanding the axial vector current $`j_\mu ^5(x)=Q^{}(x)\gamma _\mu \gamma _5Q(x)`$ in the way outlined, we can use that the divergence of the axial vector current is given by
$$_\mu j_\mu ^5=2mQ^{}\gamma _5Q\frac{\mathrm{i}}{16\pi ^2}F_{\mu \nu }^a\stackrel{~}{F}_{\mu \nu }^a,$$
(35)
where the first term contains the axial current and the second is the axial anomaly term which arises due to quantum effects. The expansion of the divergence of the axial vector current in terms of the inverse of the heavy quark mass is therefore reduced to the expansion of the axial current, which we have already performed before.
Further the axial anomaly equation (35) has some general properties, which can be used to check our result for the axial current:
First the rhs of eq. (35) vanishes in the limit of infinite quark mass. This can be understood by the fact that the regulator mass cancels the physical mass in the infinite mass limit because of the different sign in the definition of the regulator. Therefore we expect the order $`𝒪(1/m)`$ term in the expectation value of the axial current multiplied by $`2m`$ exactly to cancel the anomaly term. Indeed equation (34) gives:
$`2m\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\gamma _5|x^{𝒪(1/m)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{8\pi ^2}}\mathrm{tr}_cF_{\mu \nu }\stackrel{~}{F}_{\mu \nu }.`$ (36)
Second if one thinks of the expectation value of the axial vector current as a local large $`m`$ expansion in the background of gluon fields
$$j_\mu ^5(x)=\underset{n}{}\frac{1}{m^n}X_{\mu 5}^{(n)}(x)$$
(37)
then due to equation (35) the expectation value of the axial current in the large $`m`$ expansion is
$$2m\mathrm{tr}_{c,\gamma }x|\frac{1}{D}\gamma _5|x=\underset{n}{}\frac{1}{m^n}_\mu X_{\mu 5}^{(n)}(x).$$
(38)
This means that terms appearing in the expansion of the axial current must be of the form of a total derivative. The order $`𝒪(1/m^3)`$ term in equation (34) exactly obeys this form
$`2m\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\gamma _5|x^{𝒪(1/m^3)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{96\pi ^2m^2}}_\mu R_\mu ,`$ (39)
$`R_\mu `$ $`=`$ $`_\mu \mathrm{tr}_cF_{\alpha \beta }\stackrel{~}{F}_{\alpha \beta }+4\mathrm{tr}_c[_\alpha ,F_{\alpha \nu }]\stackrel{~}{F}_{\mu \nu }.`$ (40)
The term $`f^{abc}F_{\mu \nu }^a\stackrel{~}{F}_{\nu \alpha }^bF_{\alpha \mu }^c`$ appearing falsely in the expansion of the axial current in cannot be represented as a total derivative of a local expression<sup>2</sup><sup>2</sup>2A straightforward calculation for the instanton field shows that $`\mathrm{d}^4xf^{abc}F_{\mu \nu }^a\stackrel{~}{F}_{\nu \alpha }^bF_{\alpha \mu }^c0`$. But for dimensional reasons this nonvanishing contribution can be excluded from being generated by a surface term if the instanton field is taken in the regular gauge. Therefore the integrand cannot be a total derivative. and therefore violates the general argument given above.
The expectation value for the divergence of the axial vector current in the background of gluon fields finally reads up to order $`𝒪(1/m^4)`$
$`_\mu j_\mu ^5(x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{96\pi ^2m^2}}_\mu \left(_\mu \mathrm{tr}_cF_{\alpha \beta }\stackrel{~}{F}_{\alpha \beta }+4\mathrm{tr}_c[_\alpha ,F_{\alpha \nu }]\stackrel{~}{F}_{\mu \nu }\right)+𝒪\left({\displaystyle \frac{1}{m^4}}\right).`$ (41)
#### 2.2.3 The scalar current
Following the steps outlined in the introduction to this section the expansion of a scalar current in series of the inverse heavy quark mass yields up to order $`𝒪(1/m^3)`$:
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}|x`$ $`=`$ $`\mathrm{tr}_{c,\gamma }{\displaystyle \frac{\mathrm{d}^4k}{\left(2\pi \right)^4}\frac{1}{k^2+m^2}\underset{n=0}{\overset{\mathrm{}}{}}\left(\frac{^2\frac{\sigma }{2}F+2\mathrm{i}k}{k^2+m^2}\right)^n\left(\mathrm{i}\overline{)}\overline{)}k\mathrm{i}m\right)1}`$ (42)
$`={\displaystyle \frac{\mathrm{i}}{\left(4\pi \right)^{\frac{d}{2}}}}\left({\displaystyle \frac{1}{m^2}}\right)^{1\frac{d}{2}}m\mathrm{\Gamma }\left(1{\displaystyle \frac{d}{2}}\right)d\mathrm{tr}_c1`$
$`{\displaystyle \frac{\mathrm{i}}{24\pi ^2m}}\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$
$`+{\displaystyle \frac{1}{360\pi ^2m^3}}\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \gamma }F_{\beta \gamma }{\displaystyle \frac{7\mathrm{i}}{2880\pi ^2m^3}}^2\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$
$`{\displaystyle \frac{\mathrm{i}}{720\pi ^2m^3}}\left(11\mathrm{tr}_c[_\alpha ,[_\beta ,F_{\beta \gamma }]]F_{\alpha \gamma }\mathrm{tr}_c[_\alpha ,F_{\alpha \beta }][_\gamma ,F_{\gamma \beta }]\right).`$
The infinite constant term can be cancelled by substracting the expectation value of the scalar current for vanishing gluonic background fields. Our result eq. (42) coincides with that obtained in ref. if we neglect the total derivative terms which were ignored in ref. .
Actually the result eq. (42) with the total derivative terms neglected (and hence that of ref. ) can be easily obtained from the expansion of the determinant of the Dirac operator (16), since
$`{\displaystyle \mathrm{d}^4x\mathrm{tr}_{c,\gamma }x|\frac{1}{D}|x}`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}m}}\left(\mathrm{i}\mathrm{ln}\left(\mathrm{det}D\right)\right)`$ (43)
$`=`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{d}m}}\left(\mathrm{i}\mathrm{Tr}\mathrm{ln}D\right).`$
Our expansion of the scalar current (42) is in agreement with the result for the determinat in equation (16).
#### 2.2.4 The vector current
The heavy quark expansion of the vector current up to order $`1/m^3`$ gives exactly zero
$$\mathrm{tr}_{c,\gamma }x|\frac{1}{D}\gamma _\mu |x=0+𝒪\left(\frac{1}{m^4}\right).$$
(44)
This result can be easily anticipated from the fact that the vector current is $`C`$parity odd. This implies that the first operator contributing to heavy quark mass expansion should contain at least three gluon fields, additionaly the vector current conservation requires that this operator has the following structure $`G^3`$. From counting of dimensions we see that such operator can contribute only at $`1/m^4`$ order.
#### 2.2.5 The tensor current
For the color singlet tensor current we find that the first non-vanishing order of the expansion is $`𝒪\left(1/m^3\right)`$, yielding
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\sigma _{\mu \nu }|x`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{24\pi ^2}}{\displaystyle \frac{1}{m^3}}`$ (45)
$`\times \mathrm{tr}_c\left[F_{\alpha \beta }F_{\alpha \beta }F_{\mu \nu }+F_{\alpha \nu }F_{\beta \mu }F_{\alpha \beta }F_{\alpha \mu }F_{\beta \nu }F_{\alpha \beta }\right]+𝒪\left({\displaystyle \frac{1}{m^5}}\right).`$
We note that the rhs of the above equation vanishes in the case of $`SU(2)`$ gauge group. Actually one can show that the lhs of eq. (45) is identically zero in the case of $`SU(2)`$ gauge group. Therefore the fact that rhs of eq. (45) vanishes for $`SU(2)`$ gauge group is a powerful check of our calculations.
In order to prove that lhs of eq. (45) is zero in the case of SU(2) gauge group we use the following transformation:
$$G=C\tau ^2,$$
where $`C`$ is charge conjugation matrix in Dirac spinor space and $`\tau ^2`$ is color $`SU(2)`$ matrix. Under this transformation we have:
$`G\tau ^aG^1`$ $`=`$ $`\tau ^{aT}`$
$`G\sigma _{\mu \nu }G^1`$ $`=`$ $`\sigma _{\mu \nu }^T`$
$`GDG^1`$ $`=`$ $`D^T`$
where <sup>T</sup> is the transposition operation. The lhs. of eq. (45) should be zero, since
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\sigma _{\mu \nu }|x`$ $`=`$ $`\mathrm{tr}_{c,\gamma }x|G{\displaystyle \frac{1}{D}}\sigma _{\mu \nu }G^1|x`$ (46)
$`=`$ $`\mathrm{tr}_{c,\gamma }x|\left({\displaystyle \frac{1}{D}}\right)^T\left(\sigma _{\mu \nu }^T\right)|x`$
$`=`$ $`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}\sigma _{\mu \nu }|x.`$
Nullification of the heavy quark mass expansion of the tensor current for $`SU(2)`$ gauge group implies that lhs. of eq.(45) is zero if it is computed in the field of single instanton.
#### 2.2.6 Expansion of $`\overline{Q}_\mu \gamma _\nu Q`$.
The energy-momentum tensor of QCD can be written in Minkowski-space as
$$T^{\mu \nu }=g^{\mu \nu }_{\mathrm{QCD}}F^{\mu \alpha }F_\alpha ^\nu +\frac{\mathrm{i}}{2}\overline{\psi }\stackrel{(\mu }{\stackrel{}{}}\gamma ^{\nu )}\psi ,$$
(47)
where $`(\mu \nu )`$ denotes the symmetrization of the indices. The large $`m`$ expansion of the (not symmetrized) last term in eq. (47) yields in Euclidean space
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}_\mu \gamma _\nu |x={\displaystyle \frac{2\mathrm{i}}{(4\pi )^{\frac{d}{2}}}}\left({\displaystyle \frac{1}{m^2}}\right)^{\frac{d}{2}}\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}\right)\delta _{\mu \nu }\mathrm{tr}_c1`$
$`+{\displaystyle \frac{\mathrm{i}}{(4\pi )^{\frac{d}{2}}}}\left({\displaystyle \frac{1}{m^2}}\right)^{2\frac{d}{2}}\mathrm{\Gamma }\left(2{\displaystyle \frac{d}{2}}\right)\left({\displaystyle \frac{1}{3}}\delta _{\mu \nu }\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }+{\displaystyle \frac{4}{3}}\mathrm{tr}_cF_{\alpha \nu }F_{\alpha \mu }\right)`$
$`+{\displaystyle \frac{1}{720\pi ^2}}{\displaystyle \frac{1}{m^2}}\delta _{\mu \nu }\mathrm{tr}_cF_{\alpha \beta }F_{\beta \gamma }F_{\gamma \alpha }{\displaystyle \frac{7\mathrm{i}}{5760\pi ^2}}{\displaystyle \frac{1}{m^2}}\delta _{\mu \nu }^2\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$
$`{\displaystyle \frac{\mathrm{i}}{1440\pi ^2}}{\displaystyle \frac{1}{m^2}}\delta _{\mu \nu }\left(11\mathrm{tr}_c[_\alpha ,[_\beta ,F_{\beta \gamma }]]F_{\alpha \gamma }\mathrm{tr}_c[_\alpha ,F_{\alpha \beta }][_\gamma ,F_{\gamma \beta }]\right)`$
$`+{\displaystyle \frac{1}{2880\pi ^2}}{\displaystyle \frac{1}{m^2}}(4\mathrm{tr}_cF_{\alpha \nu }F_{\beta \mu }F_{\alpha \beta }4\mathrm{tr}_cF_{\alpha \mu }F_{\beta \nu }F_{\alpha \beta }`$
$`30\mathrm{i}\mathrm{tr}_c[_\alpha ,F_{\mu \nu }][_\beta ,F_{\alpha \beta }]`$
$`+74\mathrm{i}\mathrm{tr}_c[_\alpha ,[_\beta ,F_{\beta \nu }]]F_{\alpha \mu }+14\mathrm{i}\mathrm{tr}_c[_\alpha ,[_\beta ,F_{\beta \mu }]]F_{\alpha \nu }`$
$`26\mathrm{i}\mathrm{tr}_c[_\nu ,[_\alpha ,F_{\alpha \beta }]]F_{\beta \mu }26\mathrm{i}\mathrm{tr}_c[_\mu ,[_\alpha ,F_{\alpha \beta }]]F_{\beta \nu }`$
$`+22\mathrm{i}^2\mathrm{tr}_cF_{\beta \nu }F_{\beta \mu }3\mathrm{i}_\mu _\nu \mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$
$`26\mathrm{i}\mathrm{tr}_c[_\mu ,F_{\alpha \nu }][_\beta ,F_{\alpha \beta }]26\mathrm{i}\mathrm{tr}_c[_\nu ,F_{\alpha \mu }][_\beta ,F_{\alpha \beta }]`$
$`+40\mathrm{i}\mathrm{tr}_c[_\alpha ,F_{\alpha \nu }][_\beta ,F_{\beta \mu }]+4\mathrm{i}\mathrm{tr}_c[_\nu ,F_{\alpha \beta }][_\mu ,F_{\alpha \beta }])+𝒪({\displaystyle \frac{1}{m^4}})`$ (48)
and
$$\mathrm{tr}_{c,\gamma }x|\frac{1}{D}_\mu \gamma _\nu |x^{𝒪(m^0,m^2)}=\mathrm{tr}_{c,\gamma }x|_\mu \frac{1}{D}\gamma _\nu |x^{𝒪(m^0,m^2)}.$$
(49)
The Lorentz trace of eq. (48) can be compared with the expansion of the scalar current in eq. (42) since
$$\mathrm{tr}_{c,\gamma }x|\frac{1}{D}_\nu \gamma _\nu |x=m\mathrm{tr}_{c,\gamma }x|\frac{1}{D}|x.$$
(50)
For the trace of the rhs of eq. (48) we find
$`\mathrm{tr}_{c,\gamma }x|{\displaystyle \frac{1}{D}}_\nu \gamma _\nu |x^{𝒪(m^0,m^2)}`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{24\pi ^2}}\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$ (51)
$`{\displaystyle \frac{1}{360\pi ^2m^2}}\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \gamma }F_{\beta \gamma }+{\displaystyle \frac{7\mathrm{i}}{2880\pi ^2m^2}}^2\mathrm{tr}_cF_{\alpha \beta }F_{\alpha \beta }`$
$`+{\displaystyle \frac{\mathrm{i}}{720\pi ^2m^2}}\left(11\mathrm{tr}_c[_\alpha ,[_\beta ,F_{\beta \gamma }]]F_{\alpha \gamma }\mathrm{tr}_c[_\alpha ,F_{\alpha \beta }][_\gamma ,F_{\gamma \beta }]\right),`$
which exactly coincides with the rhs of eq. (42) multiplied by $`(m)`$. Eq. (49) further agrees with the expansion of the vector current beeing zero up to $`𝒪(1/m^3)`$ since
$$\mathrm{tr}_{c,\gamma }x|_\mu \frac{1}{D}\gamma _\nu |x\mathrm{tr}_{c,\gamma }x|\frac{1}{D}_\mu \gamma _\nu |x=_\mu \mathrm{tr}_{c,\gamma }x|\frac{1}{D}\gamma _\nu |x.$$
(52)
## 3 Intrinsic heavy quarks in light hadrons
In light hadron processes heavy quarks may give contributions only through virtual effects which are supressed by the mass of the heavy quarks. Especially for the charm quark, whose mass $`m_c1.4`$ GeV is not too large, virtual processes nevertheless may give not negligible contributions. In this section we discuss the applications of heavy quark mass expansion obtained in the previous sections.
In the following sections we use the “perturbative” normalization for the gluon field strength $`G_{\mu \nu }^a=F_{\mu \nu }^a/g`$ and rotate all expressions to Minkowsky space, see the appendix.
### 3.1 Intrinsic charm in $`\eta `$ and $`\eta ^{}`$
For the decay of the $`B`$-meson into $`\eta ^{}`$ and $`K`$-mesons in a mechanism with virtual charm quarks was suggested. In this approach the Cabbibo favored process $`b\overline{c}cs`$ is followed by the conversion of the $`\overline{c}c`$ pair directly into $`\eta ^{}`$. Its contribution to the decay amplitude is therefore direct depending on the ”intrinsic charm” component of the $`\eta ^{}`$-meson which is usually characterized by the matrix element
$$0|\overline{c}\gamma _\mu \gamma _5c|\eta ^{}(q)=\mathrm{i}f_\eta ^{}^{(c)}q_\mu .$$
(53)
Using the heavy mass expansion of the divergence of the axial vector current (41), the constant $`f_\eta ^{}^{(c)}`$ can be expressed up to the order $`1/m_c^2`$ by
$$f_\eta ^{}^{(c)}=\frac{1}{12m_c^2}0|\frac{\alpha _s}{4\pi }G_{\mu \nu }^a\stackrel{~}{G}^{\mu \nu ,a}|\eta ^{}.$$
(54)
Here we neglected the term proportional to $`[_\alpha ,G_\nu ^\alpha ]`$ in (41) which vanishes in pure Yang-Mills theory. We now can estimate the value of the constant $`f_\eta ^{}^{(c)}`$:
$$f_\eta ^{}^{(c)}2\mathrm{MeV},$$
(55)
where we have used
$$0|\frac{\alpha _s}{4\pi }G_{\mu \nu }^a\stackrel{~}{G}^{\mu \nu ,a}|\eta ^{}=0.056\mathrm{GeV}^3,$$
(56)
obtained in .
In QCD the omitted term can be related to the matrix element
$$\frac{\alpha _s}{4\pi }0|g\underset{q=u,d,s}{}\overline{q}\gamma _\nu \stackrel{~}{G}_\mu ^\nu q|\eta ^{}$$
(57)
using the equation of motion. A rough order of magnitude estimate for the contribution of the omitted term (57) to $`f_\eta ^{}^{(c)}`$ using the results of indicates an deviation at the level of $`0.3`$ MeV to the value (55). A more careful analysis of the omitted matrix element (57) can be done by using the instanton methods developed in which already have been applied by to calculations of higher twist corrections to deep-inelastic scattering.
Our estimated value for $`f_\eta ^{}^{(c)}`$ is consistent with the phenomenological analysis in where the authors dervived the bound $`65\mathrm{MeV}f_\eta ^{}^{(c)}15\mathrm{MeV}`$ from the analysis of $`\gamma \eta ^{}`$ transition form factors. From the analysis of $`(\eta ,\eta ^{},\eta _c)`$-mixing in the small value $`f_\eta ^{}^{(c)}=(6.3\pm 0.6)`$ MeV was derived, taking into account off-shellness effects in the $`\overline{c}c`$ component of $`\eta ^{}`$ also, the value $`|f_\eta ^{}^{(c)}|2.4`$ MeV was found in . Further our value for $`f_\eta ^{}^{(c)}`$ is in agreement with the phenomenological bound $`|f_\eta ^{}^{(c)}|<12`$ MeV, obtained in , and corresponds to the result $`f_\eta ^{}^{(c)}2.3`$ MeV presented in . In Ref. the divergence of the axial vector current was computed using the triangle graph for the axial anomaly with massive fermions, neglecting possible $`1/m_c^2`$ contributions like
$$f^{abc}G_{\mu \nu }^a\stackrel{~}{G}_{\nu \alpha }^bG_{\alpha \mu }^c$$
(58)
from higher order diagrams. Indeed our calculation shows that such ”truly nonabelean” operators do not contribute to the order $`1/m_c^2`$ and our result (54) therefore is exactly given by the first term of the expansion in $`1/m_c^2`$ of the triangle graph .
The small value (55) for $`f_\eta ^{}^{(c)}`$ implies that the $`b\overline{c}cs`$ mechanism does not play a major role in the $`BK\eta ^{}`$ decay mode.
Bigger values of $`f_\eta ^{}^{(c)}`$ due to the operator (58) in the expansion of the axial current up to order $`1/m_c^3`$ have been given in , where $`f_\eta ^{}^{(c)}(50180)`$ MeV and in with $`f_\eta ^{}^{(c)}(12.318.4)`$ MeV. These results have been used by a number of authors for the analysis of the charm content in noncharmed hadrons (see e.g ), but since the operator (58) violates general properties of the axial anomaly and it also does not appear in explicit calculations (see section 2.2.2), results relying on should be reconsidered.
Analogously we can immediately estimate the constant $`f_\eta ^{(c)}`$ characterizing the intrinsic charm contribution to the $`\eta `$-meson. Using
$$0|\frac{\alpha _s}{4\pi }G_{\mu \nu }^a\stackrel{~}{G}^{\mu \nu ,a}|\eta =0.020\mathrm{GeV}^3,$$
(59)
obtained in we find
$$f_\eta ^{(c)}0.7\mathrm{MeV}.$$
(60)
Since in the case of the $`\eta `$ meson the contribution of the omitted term (57) can be of the same order as $`f_\eta ^{(c)}`$ itself, the estimate (60) must be considered as a poor.
### 3.2 Intrinsic charm contribution to the proton spin
Another application of our result for the heavy quark mass expansion of the divergency of axial vector current (41) has been given in . In this paper the authors have shown that the intrinsic charm contribution to the first moment of the spin structure function $`g_1(x,Q^2)`$ of the nucleon is small contrary to the result of .
In ref. it was proven that the forward matrix element of the axial current in the leading order of heavy quark mass expansion can be computed as:
$$N(p,\lambda )|\overline{c}\gamma _\mu \gamma _5c(0)|N(p,S)=\frac{\alpha _s}{48\pi m_c^2}N(p,S)|R_\mu (0)|N(p,S).$$
(61)
Here the current $`R_\mu (0)`$ is given by eq. (40). Note that the first term in $`R_\mu `$ does not contribute to the forward matrix element because of its gradient form, while the contribution of the second one is rewritten, by making use of the equation of motion, as matrix element of the operator
$`N(p,S)|\overline{c}\gamma _\mu \gamma _5c(0)|N(p,S)`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{12\pi m_c^2}}N(p,S)|g{\displaystyle \underset{\mathrm{f}=\mathrm{u},\mathrm{d},\mathrm{s}}{}}\overline{\psi }_f\gamma _\nu \stackrel{~}{G}_\mu ^\nu \psi _f|N(p,S)`$ (62)
$``$ $`{\displaystyle \frac{\alpha _s}{12\pi m_c^2}}2m_N^3S_\mu f_S^{(2)},`$
The parameter $`f_S^{(2)}`$ was determined before in calculations of the power corrections to the first moment of the singlet part of $`g_1`$. QCD-sum rule calculations gave $`f_S^{(2)}=0.09`$ , estimates using the renormalon approach led to $`f_S^{(2)}=\pm 0.02`$ and calculations in the instanton model of the QCD vacuum give a result very close to that of QCD sum rule .
Inserting these numbers we get finally for the charm axial constant of the nucleon the estimate
$$g_A^{(c)}=\frac{\alpha _s}{12\pi }f_S^{(2)}\frac{m_N^2}{m_c^2}510^4$$
(63)
with probably a 100 percent uncertainty. Note that this contribution is of non-perturbative origin (therefore we call it intrinsic), so that it is sensitive to large distances, as soon as the factorization scale is of order $`m_c`$.
### 3.3 Intrinsic charm contribution to the nucleon tensor charge
Using the results of section 2.2.5 we can estimate the intrinsic charm contribution to the tensor charge of the nucleon. The tensor charge of the nucleon is defined as:
$$N(p,S)|\overline{c}\sigma _{\mu \nu }\gamma _5c(0)|N(p,S)=2\mathrm{i}g_T^{(c)}(p_\mu S_\nu p_\nu S_\mu ).$$
(64)
Using the result of the section 2.2.5 and the identity
$$\sigma _{\mu \nu }\gamma _5=\frac{\mathrm{i}}{2}\epsilon _{\mu \nu \alpha \beta }\sigma ^{\alpha \beta },$$
we obtain for the charm contribution to the nucleon tensor charge the following result:
$`g_T^{(c)}={\displaystyle \frac{1}{96\pi ^2}}{\displaystyle \frac{1}{m_c^3m_N^2}}`$
$`\times \epsilon ^{\lambda \rho \mu \nu }p_\lambda S_\rho N(p,S)|\mathrm{tr}_cg_s^3\left[G_{\alpha \beta }G^{\alpha \beta }G_{\mu \nu }+2G_{\alpha \nu }G_{\beta \mu }G^{\alpha \beta }\right]|N(p,S)+𝒪\left({\displaystyle \frac{1}{m_c^5}}\right).`$ (65)
The matrix element in the rhs of above equation can be rouhgly estimated in the instanton vacuum using method of . As discussed in section 2.2.5 the gluonic operator in rhs of eq. (65) is identically zero on one instanton. Therefore the first nontrivial contribution one can get from instanton–anti-instanton pair. If we compare the expression for the charm contribution to the nucleon charge with that for axial charge we see that the charm contribution to the tensor charge is suppressed by additional power of $`m_N/m_c`$ and one power of instanton packing fraction ($`\pi ^2\overline{\rho }^4/\overline{R}^4`$), however the tensor charge is enchanced by one power of $`\alpha _s(m_c)`$ <sup>3</sup><sup>3</sup>3 Let us note that the expansion parameter in the heavy quark mass expansion is $`g_sG/m_c`$, because the non-perturbative gluon field strength $`G1/g_s`$ ($`cf.`$ instanton field). Therefore $`g_s(\mu )`$ accompanied by gluon field strength is not counted as suppression.. This allows us to make a rough estimate for the charm contribution of the tensor charge:
$`g_T^{(c)}{\displaystyle \frac{m_N}{\alpha _s(m_c)m_c}}{\displaystyle \frac{(N_c2)\pi ^2\overline{\rho }^4}{N_c\overline{R}^4}}g_A^{(c)}10^4.`$ (66)
Factors of $`N_c`$ are written in a way to reproduce the large $`N_c`$ behaviour of the matrix element and to account for the fact that the operator in rhs of eq. (65) is identically zero at $`N_c=2`$.
### 3.4 Intrinsic charm contribution to the nucleon momentum
The charm contribution to the nucleon momentum can be defined as:
$`M_2^{(c)}(\mu ^2)={\displaystyle _0^1}dx_Bx_B[c(x_B)+\overline{c}(x_B)]={\displaystyle \frac{\mathrm{i}}{2(Pn)^2}}N(P)|\overline{c}n/(n)c(0)|N(P),`$ (67)
where $`c(x_B)`$ is the charm parton distribution normalized at the scale $`\mu `$, which is assumed to be $`\mu m_c`$. The light cone vector $`n`$ is arbitrary non-collinear to nucleon momentum $`P`$.
Now we can use the result of eq. (48) in order to estimate the charm contribution to the nucleon momentum carried by intrinsic charm quarks.
$`M_2^{(c)}(\mu )={\displaystyle \frac{\mathrm{i}}{2(Pn)^2}}[{\displaystyle \frac{\mathrm{i}\alpha _s(\mu )}{4\pi }}{\displaystyle \frac{1}{\left(2\frac{d}{2}\right)}}{\displaystyle \frac{4}{3}}N(P)|n^\mu n^\nu \mathrm{tr}_cG_{}^{\alpha }{}_{\nu }{}^{}G_{\alpha \mu }|N(P)`$
$`+{\displaystyle \frac{1}{120\pi ^2}}{\displaystyle \frac{g_s^3(\mu )}{m_c^2}}N(P)|n^\mu n^\nu \mathrm{tr}_cG_{\alpha \nu }G_{\beta \mu }G^{\alpha \beta }|N(P)]+𝒪\left({\displaystyle \frac{1}{m_c^4}}\right).`$ (68)
In derivation of this expression we neglected terms which are proportional to $`[^\alpha ,G_{\alpha \beta }]`$ which are suppressed by one power of $`g_s(\mu )^2`$. The first term in eq. (68) is divergent<sup>4</sup><sup>4</sup>4We show only the most singular term and actually is related to the mixing of quark and gluon operators. We can rewrite eq. (68) as follows:
$`M_2^{(c)}(\mu )={\displaystyle \frac{4}{3}}{\displaystyle \frac{\alpha _s(\mu )}{4\pi }}{\displaystyle \frac{1}{\left(2\frac{d}{2}\right)}}M_2^{(G)}(\mu )`$ (69)
$`+`$ $`{\displaystyle \frac{\mathrm{i}}{2(Pn)^2}}{\displaystyle \frac{1}{120\pi ^2}}{\displaystyle \frac{g_s^3(\mu )}{m_c^2}}N(P)|n^\mu n^\nu \mathrm{tr}_cG_{\alpha \nu }G_{\beta \mu }G^{\alpha \beta }|N(P)+𝒪\left({\displaystyle \frac{1}{m_c^4}}\right).`$
Here the first term, which is proportional to the momentum fraction carried by gluons $`M_2^{(G)}(\mu )`$ accounts for extrinsic charm. Note that the coefficient in front of this term is exactly the leading anomalous dimension $`\gamma _{qG}=4/3`$ which accounts for mixing quark and gluon twist-2 operators under QCD evolution. The intrinsic charm contribution is given by the second term, so that we have estimates:
$`M_2^{(c),\mathrm{intrinsic}}(\mu )={\displaystyle \frac{\mathrm{i}}{2(Pn)^2}}{\displaystyle \frac{1}{120\pi ^2}}{\displaystyle \frac{g_s^3(\mu )}{m_c^2}}N(P)|n^\mu n^\nu \mathrm{tr}_cG_{\alpha \nu }G_{\beta \mu }G^{\alpha \beta }|N(P)+𝒪\left({\displaystyle \frac{1}{m_c^4}}\right).`$ (70)
We see that the momentum fraction carried by intrinsic charm in the nucleon is related to the value of nucleon matrix element:
$$N(P)|n^\mu n^\nu \mathrm{i}g_s(\mu )^3\mathrm{tr}_cG_{\alpha \nu }G_{\beta \mu }G^{\alpha \beta }|N(P).$$
One can easily see that this matrix element in the theory of instanton vacuum is zero in one-instanton approximation, the same as matrix element
$$N(P)|n^\mu n^\nu g_s(\mu )^2\mathrm{tr}_cG_{}^{\alpha }{}_{\nu }{}^{}G_{\alpha \mu }|N(P).$$
Keeping in mind that for instanton field $`G1/g_s`$ we can write:
$$\frac{N(P)|n^\mu n^\nu \mathrm{i}g_s(\mu )^3\mathrm{tr}_cG_{\alpha \nu }G_{\beta \mu }G^{\alpha \beta }|N(P)}{N(P)|n^\mu n^\nu g_s(\mu )^2\mathrm{tr}_cG_\nu ^\alpha G_{\alpha \mu }|N(P)}=\mathrm{\Lambda }^2,$$
(71)
where $`\mathrm{\Lambda }`$ is parameter of the dimension of mass whose value can be obtained using various nonperturbative methods in QCD: lattice calculation, QCD sum rule, theory of instanton vacuum. Generically we expect that this mass parameter is of order of typical strong interaction scale $`\mathrm{\Lambda }1`$ GeV. Now we can rewrite eq. (70) in terms of this parameter and momentum fraction carried by gluons in the nucleon at scale $`\mu m_c`$ as:
$`M_2^{(c),\mathrm{intrinsic}}(\mu )={\displaystyle \frac{\alpha _s(\mu )}{30\pi }}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_c^2}}M_2^G(\mu )+𝒪\left({\displaystyle \frac{1}{m_c^4}}\right).`$ (72)
If we assume that $`\mathrm{\Lambda }^2=`$few GeV<sup>2</sup> than we get the estimate for the charm contribution to the nucleon momentum:
$$M_2^{(c),\mathrm{intrinsic}}(\mu )=\mathrm{few}\times 10^3.$$
(73)
We see that the heavy quark mass expansion of local currents allows us to reduce the problem of estimate of intrinsic charm content of the nucleon to the calculation of the ratio (71). The latter ratio can be computed using various methods of nonperturbative QCD, probably the most promising would be a calculation of this ratio in lattice QCD.
Recent analysis of refs gives for $`M_2^{(c),\mathrm{intrinsic}}`$ values at the level of fraction of percent what is in agrrement with our estimate (73).
Let us note that, since we performed the heavy quark mass expansion of heavy quark part of energy momentum tensor not neglecting total derivatives, one can compute also its non-forward nucleon matrix element. From the non-forward matrix element of energy momentum one can obtain the total angular momentum carried by intrinsic heavy quarks in the nucleon using Ji’s sum rules . The corresponding estimates we shall report elsewhere.
## 4 Conclusions
In this paper we have computed the heavy quark mass expansion of various local heavy quark currents. The details of the technique are illustrated on the example of heavy quark mass expansion of the pseudoscalar density $`\overline{Q}\gamma _5Q`$. This operator plays an important role in problems related to intrinsic charm contribution to the proton spin and to intrinsic charm content of $`\eta ,\eta ^{}`$ mesons.
We corrected the mistakes in refs. for heavy quark mass expansion of the operator $`\overline{Q}\gamma _5Q`$ . In these papers large intrinsic charm contribution to the proton spin and to intrinsic charm content of $`\eta ,\eta ^{}`$ mesons was obtained due to contribution of the operator $`f^{abc}G_{\mu \nu }^a\stackrel{~}{G}_{\nu \alpha }^bG_{\alpha \mu }^c`$ which appeared in heavy quark mass expansion of the operator $`\overline{Q}\gamma _5Q`$ presented in refs. . We showed that coefficient in front of this operator is identically zero (the result which actually follows from general properties of the axial anomaly ), so that the physical effects based on presence of the above operator discussed in refs. are absent.
For the first time we presented the full results<sup>5</sup><sup>5</sup>5Not neglecting total derivatives and terms proportional to $`[^\mu ,G_{\mu \nu }]`$ for heavy quark mass expansion of the operators $`\overline{Q}Q`$ (to the order $`1/m^3`$), $`\overline{Q}\gamma _5Q`$ (to the order $`1/m^3`$), $`^\mu \overline{Q}\gamma _\mu \gamma _5Q`$ (to the order $`1/m^2`$), $`\overline{Q}\gamma _\mu Q`$ (to the order $`1/m^3`$), $`\overline{Q}\sigma _{\mu \nu }Q`$ (to the order $`1/m^3`$), and $`\overline{Q}\gamma _\mu _\nu Q`$ (to the order $`1/m^2`$).
The results obtained for heavy quark mass expansion allowed us to estimate the intrinsic charm content of $`\eta ^{},\eta `$ mesons as well the charm contribution to the proton spin, nucleon tensor charge and to the fraction of nucleon momentum carried by intrinsic charm. In the case of charm content of $`\eta ^{},\eta `$ mesons and intrinsic charm contributions to the proton spin we reduce the calculations of these quantities to matrix elements which are already known either phenomenologically or were computed previously. In other cases, like intrinsic charm contribution to the nucleon tensor charge and to energy momentum tensor, the problem is reduced to matrix elements of gluon operators which can be estimated using various nonperturbative methods in QCD: lattice calculation, QCD sum rule, theory of instanton vacuum.
We made here rough order of magnitude estimate of matrix elements of gluon operators appearing in heavy quark mass expansion of tensor current and of energy momentum tensor using instanton model of QCD vacuum. More quantitative estimates will be given elsewhere. The predictions for intrinsic charm contribution to various observables are summarized in Table 1.
Acknowledgments:
We gratefully acknowledge useful discussions with F. Araki, P.V. Pobylitsa, A. Schäfer, E. Shuryak and O.V. Teryaev. The work has been supported in parts by DFG and BMFB.
## 5 Appendix
For the euclidization we use the follwing conventions:
$$\begin{array}{cccccccccc}\hfill \mathrm{i}x_\mathrm{M}^0& =& x_{4,\mathrm{E}},\hfill & \hfill x_\mathrm{M}^k& =& x_{k,\mathrm{E}}\hfill & \mathrm{d}^4x_\mathrm{M}& \hfill =& \mathrm{i}\mathrm{d}^4x_\mathrm{E},& \\ \hfill _\mathrm{M}^0& =& \mathrm{i}_{4,\mathrm{E}},\hfill & \hfill _\mathrm{M}^k& =& _{k,\mathrm{E}},\hfill & & & & \\ \hfill A_\mathrm{M}^0& =& \mathrm{i}A_{4,\mathrm{E}},\hfill & \hfill A_\mathrm{M}^k& =& A_{k,\mathrm{E}}.\hfill & & & & \end{array}$$
(74)
The covariant derivative therefore reads in Minkowski and in Euclidean space-time:
$`_\mathrm{M}^\mu `$ $`=`$ $`\left(^\mu \mathrm{i}A^\mu (x)\right)_\mathrm{M},`$ (75)
$`_{\mu ,\mathrm{E}}`$ $`=`$ $`\left(_\mu \mathrm{i}A_\mu (x)\right)_\mathrm{E}.`$ (76)
The field strength, defined as
$$F_{\mu \nu }^a=_\mu A_\nu _\nu A_\mu +f^{abc}A_\mu ^bA_\nu ^c=\mathrm{i}[_\mu ,_\nu ]$$
(77)
transforms as
$$F_{ij,\mathrm{M}}=F_{ij,\mathrm{E}},F_{0j,\mathrm{M}}=\mathrm{i}F_{4j,\mathrm{E}}.$$
(78)
For the Dirac matrices we choose the conventions:
$$\gamma _\mathrm{M}^0=\gamma _{4,\mathrm{E}},\gamma _\mathrm{M}^k=\mathrm{i}\gamma _{k,\mathrm{E}},$$
(79)
and $`\gamma _5`$ is defined within this paper as:
$$\gamma _{5,\mathrm{M}}=\gamma _\mathrm{M}^5=\mathrm{i}\left(\gamma ^0\gamma ^1\gamma ^2\gamma ^3\right)_\mathrm{M}=\left(\gamma _1\gamma _2\gamma _3\gamma _4\right)_\mathrm{E}=\gamma _{5,\mathrm{E}}.$$
(80)
With
$$\epsilon _\mathrm{M}^{0123}=\epsilon _{0123,\mathrm{M}}=+1=\epsilon _{1234,\mathrm{E}},$$
(81)
it yields
$$\mathrm{tr}_\gamma \left[\gamma _5\gamma _\alpha \gamma _\beta \gamma _\gamma \gamma _\delta \right]_\mathrm{E}=4\epsilon _{\alpha \beta \gamma \delta ,\mathrm{E}}.$$
(82)
The fermionic fields transform as
$$\psi _\mathrm{M}=\psi _\mathrm{E},\overline{\psi }_\mathrm{M}=\mathrm{i}\psi _\mathrm{E}^{},$$
(83)
so the Dirac operator in Euclidean space-time reads:
$$D=\mathrm{i}\overline{)}+\mathrm{i}m.$$
(84)
In section 2.1 we have used the following transformation properties for the effective action and the appearing operators:
$`S_{\mathrm{eff},\mathrm{M}}`$ $`=`$ $`\mathrm{i}S_{\mathrm{eff},\mathrm{E}},`$ (85)
$`\left(F_{\alpha \beta }F^{\alpha \beta }\right)_\mathrm{M}`$ $`=`$ $`\left(F_{\alpha \beta }F_{\alpha \beta }\right)_\mathrm{E},\left(F_{\alpha \beta }F_{}^{\beta }{}_{\gamma }{}^{}F^{\gamma \alpha }\right)_\mathrm{M}=\left(F_{\alpha \beta }F_{\beta \gamma }F_{\gamma \alpha }\right)_\mathrm{E}.`$ (86) |
warning/0002/cond-mat0002112.html | ar5iv | text | # A reconstruction from small-angle neutron scattering measurements of the real space magnetic field distribution in the mixed state of Sr2RuO4.
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## Abstract
We have measured the diffracted neutron scattering intensities from the square magnetic flux lattice in the perovskite superconductor Sr<sub>2</sub>RuO<sub>4</sub>, which is thought to exhibit p-wave pairing with a two-component order parameter. The relative intensities of different flux lattice Bragg reflections over a wide range of field and temperature have been shown to be inconsistent with a single component Ginzburg-Landau theory but qualitatively agree with a two component p-wave Ginzburg-Landau theory.
preprint: Preprint
\]
The discovery of superconductivity at temperatures near 1K in strontium ruthenate has excited great interest because it is a superconducting layered perovskite which does not contain copper. However it shows great differences from the High-T<sub>c</sub> cuprates: it is a stoichiometric undoped compound with a long mean free path, in which the electrons form a Fermi liquid with a well-established quasi-two-dimensional Fermi surface . Furthermore, it was suggested that the strongly interacting electrons pair in a triplet p-wave state (rather than the singlet, mainly d-wave state which is believed to occur in hole-doped cuprates). Clear evidence of non s-wave pairing in this compound has been provided by the observation that nonmagnetic impurities strongly suppress T<sub>c</sub>, which extrapolates to $`1.5`$ K in the clean limit. Strong support for triplet (p-wave) pairing is given by the results of Ishida et al. who have measured the Knight shift with a field parallel to the RuO<sub>2</sub> planes ; the spin susceptibility measured by the Knight shift is not suppressed below T<sub>c</sub>,unlike a singlet superconductor. Also, $`\mu `$SR measurements in the Meissner state in zero field have revealed spontaneous fields, which can be generated by domain boundaries, surfaces and impurities in a superconductor which breaks time-reversal symmetry . Such states can arise most naturally with p-wave pairing, but also are possible with d-wave singlet pairing.
Agterberg argued that if the pairing was time-reversal symmetry breaking p-wave, then in tetragonal symmetry the d-vector has the symmetry $`\widehat{𝐳}\mathrm{exp}(\pm \mathrm{i}\phi )`$ ($`\phi `$ is the azimuthal angle about the tetragonal c axis), and a two-component Ginzburg Landau (TCGL) theory would be expected to describe the superconductor. In zero field, this gives two degenerate states which are related by time reversal; with a field applied in the c-direction perpendicular to the planes, one is dominant, but the other is also present . Under these conditions, a small amount of anisotropy in the Fermi surface would lead to a square flux lattice instead of a triangular one, with the orientation of the square flux line lattice (FLL) relative to the crystal axes determined by the orientation of the fourfold anisotropy of the paired electrons. The FLL structure has been observed in this material and is observed to be square over a wide range of field and temperature. The nearest-neighbour directions in the square FLL are at 45 to the Ru-O-Ru directions in the crystal lattice .
These results are consistent with the pairing wavefunction described above. However, a square FLL is also seen in borocarbide superconductors, which are definitely non-p-wave . Also, one can measure spontaneous fields in a superconductor by $`\mu `$SR due to other causes or from other states than that proposed, and application of a strong field in the basal plane to observe the Knight shift might alter the pairing state. Hence, it is important to obtain further evidence as to what kind of superconductivity occurs in strontium ruthenate. Here we present a detailed study of the scattered neutron intensities from the FLL. We show that they are not consistent with a single component Ginzburg-Landau model. Also we demonstrate how the local $`B(𝐫)`$ may be reconstructed from our data and show that the FLL structure is quite different from the Abrikosov one.
We shall present measurements of intensities of higher-order Bragg reflections from the FLL so we consider how they are related to the FLL structure. The formula for the integrated intensity $`I_{hk}`$ of a $`(h,k)`$ diffracted peak of wavevector $`𝐪_{hk}`$ gives:
$$I_{hk}\frac{F_{hk}^2}{q_{hk}},$$
(1)
where $`F_{hk}`$ is a spatial Fourier component of the local field $`B(𝐫)`$ in the mixed state:
$$B(𝐫)=\underset{h,k}{}F_{hk}\mathrm{exp}(\mathrm{i}𝐪_{hk}𝐫).$$
(2)
In the Abrikosov solution of the Ginzburg-Landau (GL) equations (as applied to a square lattice) , the $`F_{hk}`$ are given by:
$$F_{hk}(1)^{(h^2+k^2+hk)}\mathrm{exp}\left(\frac{\pi }{2}(h^2+k^2)\right);$$
(3)
this rapidly falls off with $`q`$ (see Table I).
The Abrikosov solution is only valid near $`B_{c2}`$. In high-$`\kappa `$ superconductors, with the field not close to $`B_{c2}`$, the London expression is appropriate instead. This gives $`F_{hk}1/(1+q_{hk}^2\lambda ^2)`$. Note that unlike the Abrikosov solution, all the $`F_{hk}`$ are positive. Table I shows that the Fourier components fall off much less rapidly with $`q`$. However, strontium ruthenate has a value of the Ginzburg Landau parameter $`\kappa =\lambda /\xi 2.0`$ for the field along the c axis, which means that the London approach is not realistic except at very low inductions. Therefore, to see what conventional GL theory predicts for this material at lower fields, one must use the Brandt numerical solution of the GL equations . Typical results are given in Table I.
Next, we consider the Agterberg TCGL solution , which is equivalent to the Abrikosov one, except that there are two complex order parameters instead of only one. In the mixed state with $`B`$ parallel to c both components are automatically present because of mixed gradient terms in the free energy functional . Typical values from this theory for $`F_{hk}`$, relative to $`F_{10}`$ and the resulting SANS intensities are given in Table I. It may be seen that the two-component theory gives intensities that fall off much less rapidly with $`q`$ than those given by the one-component Abrikosov solution.
Under the conditions of our experiments, where the field is not close to $`B_{c2}`$, it may be argued that the Abrikosov approximation used by Agterberg is not appropriate. However, recently Heeb and Agterberg have solved numerically the GL equations at all fields for the TCGL case . We also give in Table I a list of Fourier components from these calculations, using values of parameters that appear to describe our results quite well.
The corresponding vortex structures in real space are shown in Figures 1-4. Note that there is a minimum field point in the two-component theory (for the conditions of our experiment) which lies between the positions of the flux line cores, not in the centre of the square. We give results of this theory for two values of the parameter $`\nu `$ ($`1<\nu <1`$) which describes the degree of fourfold anisotropy of superconducting electrons ($`\nu =0`$ corresponds to a cylindrical Fermi surface). We note that the results do not change greatly with $`\nu `$. Hence the qualitative difference between Figures 1 and 3 is due to the difference between TCGL and GL theories rather than effects of fourfold crystal anisotropy. It may be that $`\nu `$ is quite small since $`|\nu |>0.0114`$ is sufficient to stabilise a square FLL and align it to the crystal lattice with an orientation determined by the sign of $`\nu `$ .
We now turn to measurements of the FLL structure. Single crystal Sr<sub>2</sub>RuO<sub>4</sub> was prepared by the floating zone technique with excess RuO<sub>2</sub> as a flux . Six plates of total mass 556 mg were cleaved from the as-grown crystal and annealed for 72 hours in air at 1420C to remove defects and increase T<sub>c</sub>, which was 1.39K with a width (10-90%) of $``$50 mK. With the field applied parallel to the c-axis at 100 mK, the value of $`B_{c2}`$ was 58mT. For the small angle neutron scattering (SANS) measurements, the samples were mounted with conducting silver paint as an aligned mosaic with their c-axes perpendicular to a copper plate, which was mounted on the mixing chamber of a dilution refrigerator. This was placed between the poles of an electromagnet, which had holes parallel to the field for transmission of neutrons. The magnetic field was parallel to the c-axes of the crystals within 2, and the FLL was observed using long-wavelength neutrons on instrument D22 at the Institut Laue Langevin. Typical wavelengths employed were 14.6 Å, with a wavelength spread (FWHM) of 12%; the neutron beam was incident nearly parallel to the applied field, and the transmitted neutrons were registered at a $`128\times 128`$ pixel multidetector (pixel size $`7.5\times 7.5`$ mm<sup>2</sup>) placed 17.71 m beyond the sample. Typical results are shown in Fig 5. In addition to the strongest $`\{10\}`$ reflections, the $`\{11\}`$ reflections are strong, and higher orders are present. The intensity of the strongest diffraction spot is $`<`$10<sup>-3</sup> of the incident beam intensity, so these higher order reflections are not due to multiple scattering. Their intensities are recorded in Table I: it will be noted that they are much larger than those given by the Abrikosov structure.
To reconstruct the $`B(𝐫)`$ of the FLL corresponding to these results, we require the sign of $`F_{hk}`$ relative to $`F_{10}`$ (the FLL is centrosymmetric, so all the $`F_{hk}`$ are real). The most important component after $`F_{10}`$ is $`F_{11}`$. If it has the same sign as $`F_{10}`$, then the $`\{11\}`$ components add in phase at the flux line cores to give a field peak that is sharper than the field minimum. Measurements of the field distribution in strontium ruthenate by $`\mu `$SR show that this is the case. This sign for $`F_{11}`$ is not surprising, since all models in Table I give it as positive. For the small contributions of $`F_{20}`$ and $`F_{21}`$, we may assume the same signs as given by the Agterberg and Abrikosov solutions: taking the London sign makes a large difference to $`B(𝐫)`$, and also can be ruled out by $`\mu `$SR results . The reconstruction of $`B(𝐫)`$ is shown in Fig. 4. Note that it is completely different from the Abrikosov or Brandt solutions to the GL equations, and in good qualitative agreement with the TCGL predictions.
The results we have given so far correspond to low temperature and a particular magnetic field value. In Table II we present the values of the form factors $`F_{10}`$ and $`F_{11}`$ for a range of fields at 100mK. Also, in Fig 3 we plot versus temperature the ratio of the Fourier components for the strongest two reflections $`F_{11}/F_{10}`$ at 10, 20 and 30mT. Remarkably, this ratio varies little with field and temperature and does not tend to the Abrikosov value as T$``$T<sub>c</sub>. Non-local effects and deviations from GL theory in ultra pure superconductors should both die away at high temperatures. Therefore, these effects are not expected to be the cause of the flux line shapes we report, although they may affect the details of $`B`$(r) at low temperatures.
In conclusion, the strength of the higher order reflections from the FLL in strontium ruthenate and their temperature dependence certainly show that a standard one component Ginzburg Landau model is insufficient to explain the observed diffraction pattern. However, our results are in good qualitative but not perfect agreement with a two component Ginzburg Landau theory. Unconventional flux line shapes in this material are strong evidence for unconventional superconductivity in Sr<sub>2</sub>RuO<sub>4</sub>.
We thank J-L Ragazzoni of the ILL for setting up the dilution refrigerator, E.H. Brandt for a copy of his code to solve the GL equations and G.M. Luke for communicating his results prior to publication. This work was supported by the U.K. E.P.S.R.C., CREST of Japan Science and Technology Corporation, and the neutron scattering was carried out at the Institut Laue-Langevin, Grenoble.
Email address: P.G.Kealey@bham.ac.uk |
warning/0002/astro-ph0002008.html | ar5iv | text | # ISO-SWS spectroscopy of NGC 1068 1footnote 11footnote 1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.
## 1 Introduction
NGC 1068 is one of the nearest and probably the most intensely studied Seyfert 2 galaxy. Observations in all wavelength bands from radio to hard X-rays have formed a uniquely detailed picture of this object. NGC 1068 has played a key role in the development of unified scenarios for Seyfert 1 and Seyfert 2 galaxies (Antonucci & Miller (1985)), in the study of molecular gas in the nuclear region of Seyferts (e.g. Myers & Scoville (1987); Tacconi et al. (1994)), and in elucidating the importance of star formation activity coexistent with the AGN, both on larger (e.g. Telesco & Decher (1988)) and smaller (Macchetto et al. (1994); Thatte et al. (1997)) scales. NGC 1068 hosts a prominent Narrow Line Region (NLR) that is approximately cospatial with a linear radio source with two lobes (Wilson & Ulvestad (1983)). The narrow emission line region has been extensively characterized from subarcsecond clouds probed by HST (Evans et al. (1991); Macchetto et al. (1994)), the $`5`$ arcseconds of the NLR hosting most of the line flux (e.g. Walker (1968); Shields & Oke (1975); Cecil et al. (1990)), and the ionization cone and extended emission line region (Pogge (1988); Unger et al. (1992)) extending to radii of at least 30″ (1″ = 72 pc at the distance of 14.4 Mpc, Tully (1988)). The velocity field is complex, with an ensemble of rapidly moving clouds dominating the inner arcseconds and a more quiescent rotation pattern prevailing at larger radii (e.g. Walker (1968); Alloin et al. (1983); Meaburn & Pedlar (1986); Cecil et al. (1990)). While most of the excitation of the narrow line region and the extended emission line region is likely through photoionization by the central AGN (Marconi et al. (1996)), high resolution observations suggest kinematic disturbance and possibly shock excitation of regions close to the radio outflow (e.g. Axon et al. (1998)).
With ESA’s Infrared Space Observatory ISO, sensitive mid-infrared spectroscopy of AGNs became possible, with detections of a broad range of low- and high-excitation fine structure lines, recombination lines, and pure rotational lines from such sources (Moorwood et al. (1996); Sturm et al. (1999); Alexander et al. (1999))). Model predictions of the mid-IR spectra of AGN had been obtained prior to ISO (e.g. Spinoglio & Malkan (1992)), but observations were restricted by limited sensitivity and focussed primarily on the continuum emission and broad features rather than emission lines.
In this paper, we present the ISO-SWS spectra of NGC 1068 and draw conclusions on the structure of the narrow line region that can be obtained mainly from the comparison of optical and reddening-insensitive infrared lines, and discuss the nature of the mid-infrared molecular hydrogen emission. The ISO-SWS data of NGC 1068 are analysed further in several companion papers. Alexander et al. (2000) use photoionization modelling based on the ISO fine structure line set and other NLR lines to model the shape of the AGN’s spectral energy distribution. Lutz et al. (2000) analyze limits on emission from the obscured broad line region. Finally, Sturm et al. (2000) discuss continuum energy distribution and features of NGC 1068 in conjunction with ISO-SWS spectra of other galaxies.
Our paper is organised as follows. In §2 we discuss the ISO-SWS observations and data reduction. §3 presents results and implications of the density of the narrow line region. §4 uses infrared line profiles in comparison to optical ones to constrain the structure of the narrow line region. We discuss the mid-infrared molecular hydrogen emission in §5 and summarize in §6.
## 2 Observations and data reduction
We have used the Short Wavelength Spectrometer SWS (de Graauw et al. (1996)) on board the Infrared Space Observatory ISO (Kessler et al. (1996)) to observe the nuclear region of NGC 1068. In Table 1 we present a log of our observations. We carried out observations in the SWS01 mode which provides a full 2.4 – 45$`\mu `$m scan at slightly reduced spectral resolving power, as well as observations in the SWS02 and SWS06 modes targeted at full resolution observations of individual lines or short ranges. Because of the large width of the emission lines in the NGC 1068 narrow line region, we have mostly relied on the SWS06 mode which can be set up to provide wider continuum baselines than standard SWS02 line scans. We supplement observations from our ISO guaranteed and open time with serendipitous information on some fine structure lines obtained in another ISO project (PI G. Stacey), the main results of which are to be presented elsewhere.
Table 1 includes also the position angle of the long axis of the SWS apertures for the various observations. Our pointing was always centered on the nucleus of NGC 1068, but the SWS apertures range from 14″$`\times `$20″ at short wavelengths to 20″$`\times `$33″ at the longest wavelengths (de Graauw et al. (1996)). At the long wavelengths, we partly include the $``$15″ radius ring of star forming regions encircling the nucleus of NGC 1068. The apertures were always oriented in approximately north-south (or south-north) direction, with position angles between -11° and -23°.
We have analyzed the data using the SWS Interactive Analysis (IA) system (Lahuis et al. (1998); Wieprecht et al. (1998)) and calibration files of July 1998. A preliminary account of part of the observations is given by Lutz et al. (1997). Since then, calibration files have been updated for wavelength calibration and in particular with respect to the SWS relative spectral response function, leading to more reliable intercalibration between the ‘AOT bands’ forming a full SWS spectrum. Our data reduction started in the standard way and continued with steps of (interactive) dark current subtraction and matching up- and downscans. We eliminated data from those detectors of band 3 that were most noisy during a particular revolution, and from interactively identified regions with single detector signal jumps in bands 1 and 2, and simultaneous 12-detector signal jumps in band 3. After relative spectral response correction and flux calibration, we ‘flatfielded’ the 12 detectors of a band to a consistent level, corrected for the ISO velocity, and extracted the AAR data product. Redundant scans of the same line were shifted to a consistent level. Single-valued spectra were produced by kappa-sigma clipping the AAR dot cloud and rebinning it with a resolution of typically 3000 which does not lead to significant smearing for NGC 1068 linewidths. For those ranges affected by fringes, the single-valued spectra were defringed using the iterative sine fitting option of the aarfringe module within the SWS Interactive Analysis.
The large number of observations required special treatment of redundant data. In addition, observations from revolution 285 where apparently affected by a slight ($`23\mathrm{}`$?) pointing problem, which occured occasionally in the earlier phase of the ISO mission. Since SWS beam profiles in some AOT bands are peaked and slightly offset with respect to the nominal pointing (A. Salama, 1999, priv. comm.), modest pointing offsets can cause noticeable flux losses in some AOT bands and resulting band mismatches. Such mismatches were evident in revolution 285 band 3 data. Since most of the NGC 1068 mid-infrared flux comes from a small region (Cameron et al. (1993); Braatz et al. (1993); Bock et al. (1998)), we corrected for this problem and the small scatter between other observations by the following scaling procedure: For the SWS01 full spectrum obtained in revolution 285, the individual AOT bands were scaled to obtain both good match at band limits, and good agreement with the overall flux level as estimated from our other SWS data and ground-based photometry (Lebofsky et al. (1978); Rieke & Low (1975)). At wavelengths below 10$`\mu `$m photometry from different epochs should be used with great caution because of the known variability (Glass (1997)), our fluxes are however in good agreement with the photometry of Glass for the ISO epoch. All other data were then scaled to this SWS01 spectrum by the ratio inferred from the continuum flux densities. We believe the final flux scale to be accurate within the 20-30% typical for SWS data (Schaeidt et al. (1996)).
Accurate wavelength calibration of the SWS grating spectrometer is central for part of our line profile analysis, since shifts between lines in NGC 1068 tend to be of the order 300 km/s or less (Marconi et al. (1996)). Valentijn et al. (1996) deduce an accuracy of $``$30km/s from extensive calibrations during the SWS performance verification phase. Since then, a slow secular drift in SWS wavelength calibration has been calibrated to similar accuracy. We have tested the wavelength calibration of the NGC 1068 data, using identical calibration files to analyze spectra of the planetary nebulae NGC 7027 and NGC 6543 taken close to revolution 633 where some of the most important NGC 1068 lineprofiles were taken. We confirm the excellent accuracy from these observations, the largest error not exceeding the value given by Valentijn et al. (1996). This test and the good internal consistency of velocities measured in NGC 1068 for different lines from the same species spanning most of the SWS wavelength range (e.g. H<sub>2</sub>, see Table 2) leads us to adopt an upper limit of 50 km/s for any systematic errors in our wavelength scale, taking into account a margin for mispointing.
Two emission features, which were tentatively detected in the preliminary analysis of Lutz et al. (1997) could not be confirmed with the larger observational database and the improved calibration. A broad emission feature near 19$`\mu `$m, which might be interpreted as silicate emission, is not confirmed with the new spectral response calibration. An emission line at 28$`\mu `$m was identified as an unusually strong H<sub>2</sub> S(0) line. This identification was later found suspect because of the line’s larger width compared to the other H<sub>2</sub> lines observed in NGC 1068. The line was not confirmed in deeper follow up observations. Detailed inspection of the original data indeed shows that it is an artifact of a highly unlikely coincidence of detector ‘glitches’ at the expected wavelength of the S(0) transition.
## 3 Results
The 2.4-45$`\mu `$m full spectrum of NGC 1068 is displayed in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.. Some solid state features are superposed on the strong AGN-heated mid-infrared continuum. These include 3.4$`\mu `$m C-H absorption, 9.6$`\mu `$m silicate absorption, and 7.7, 8.6, 11.3$`\mu `$m ‘PAH’ emission (see Sturm et al (2000) for a discussion in conjunction with other SWS spectra of galaxies). Bright fine structure emission lines, mainly originating in the narrow line region, are already visible in the full spectrum. Figures ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. and ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. show individual emission lines. The display range is chosen to be $`\pm `$2500 km/s around systemic velocity for recombination and fine structure lines, and $`\pm `$1000 km/s for the much narrower molecular lines. Throughout this paper, we adopt a systemic velocity of 1148km/s (Brinks et al. (1997)). Good rest wavelengths are available for the observed transitions from the literature and from recent ISO determinations (Feuchtgruber et al. (1997)). Many lines were observed repeatedly, Figures ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. and ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. show only the best quality data. Table 2 lists the measured line fluxes and limits, presenting averages of independent measurements with higher weight given to better data. Table 2 includes also upper limits for some transitions (not shown in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.) that were observed with good enough signal-to-noise ratio. We list such limits for transitions from elements like Na or Ar where other ionization stages are detected.
Fluxes of relatively narrow lines, for example from H<sub>2</sub> and \[Si II\], were measured by direct integration of the continuum subtracted line profiles. The large linewidth makes this procedure error prone for faint lines from the NLR, where continuum definition is the main source of measurement error. The relative constancy of NLR line profiles over a wide range of lower ionization potentials (see below) lead us to adopt a different procedure to measure fluxes for the NLR lines: We derived a simple two-gaussian template from the brightest NLR lines (\[O IV\] 25.89$`\mu `$m (note nearby \[Fe II\]), \[Ne V\] 24.32$`\mu `$m, \[Ne VI\] 7.652$`\mu `$m) and used fits of this template (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.) plus a linear continuum to measure the fluxes of fainter NLR lines. The two gaussian components have FWHM 333 and 1246 km/s, peak ratio narrow/wide 1.34, and the wider component is blueshifted by 100km/s. In fitting, we varied only continuum flux and slope, total line flux, and total velocity. The fluxes measured this way agreed very well ($``$10%) with those determined by direct integration not only for the lines used to derive the template, but also for other bright NLR lines like \[Mg VIII\] 3.028$`\mu `$m. The fit thus preserves the fluxes for bright lines and is preferable for faint lines where continuum subtraction is the dominant source of error.
For the blended lines of \[Mg VII\] and H<sub>2</sub> (0-0) S(7) near 5.5$`\mu `$m, we list the fluxes resulting from a tentative gaussian fit using two components for \[Mg VII\] and one component for H<sub>2</sub>. The uncertainty of the S(7) flux is particularly large, up to a factor 2. The H<sub>2</sub> excitation diagram (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.) in fact suggests that it may be overestimated. Similar caution has to be applied to the tentative flux listed for the \[Fe II\] 25.99$`\mu `$m line. A slight shoulder appears in the long wavelength wing of the \[O IV\] 25.89$`\mu `$m transition, but its flux is very uncertain and may at best be good enough to serve for consistency checks with other \[Fe II\] lines. Feuchtgruber et al. (1997) discuss evidence that the lines of \[Ar III\] and \[Mg VII\] at 9.0$`\mu `$m are blended at the resolution of SWS. We list only a total flux in Table 2.
We detect a weak unidentified feature at rest wavelength about 7.555$`\mu `$m. If real, its width would suggest a NLR origin. An instrumental origin due to an imperfection of the relative spectral response function (RSRF) cannot be excluded but is unlikely since the RSRF shows very little structure at that wavelength. Also, instrumental ‘ghosts’ to strong SWS lines are not known at such a level. A possible interpretation of this feature is that it is a blueshifted ($``$3800 km/s) component of the nearby strong \[Ne VI\] line containing $``$1.5% of the total line flux. Residual instrumental fringing prevents us from looking for analogous components near other strong lines such as \[Ne V\] or \[O IV\]. At this point, the feature must be considered as possibly real but without an obvious identification by a line that is potentially strong in AGN spectra. Also, no line is seen at this wavelength in archival ISO spectra of the high excitation planetary nebula NGC 6302.
We postpone a detailed discussion of the line profiles to section 4. Figures ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. and ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. already suggest, however, that we are dealing with three distinct components: Fine structure transitions from species with lower ionization potential $``$ 40 eV (i.e. from \[Ne III\] upwards) have similar wide profiles and apparently originate in the NLR. Fine structure lines from lower ionization stages, in particular \[Ne II\] and \[S III\], are narrower and are most likely contaminated by star formation within their beams. This limits their use in modelling of the AGN-excited NLR spectrum (Alexander et al. (2000)) to upper limits rather than measurements. Inspection of Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. suggests that the starburst contribution strongly dominates the low excitation lines like \[S III\] 33.48$`\mu `$m and \[Si II\] 34.81$`\mu `$m which are measured with the largest aperture. Lines measured with intermediate apertures like \[Ne II\] 12.81$`\mu `$m and in particular \[S III\] 18.71$`\mu `$m still show strong wings and will have a considerable NLR contribution. The smallest line widths are measured for transitions of molecular hydrogen.
### 3.1 Density of the narrow line region
The fine structure lines detected by SWS can be used to determine the density and, in conjunction with optical forbidden lines, the electron temperature of the line emitting gas in the narrow line and coronal line regions. Differences in infrared and optical lineprofiles (§4) discourage a determination of electron temperatures from the integrated fluxes, which would not account for significant variations in extinction across the NLR. A reliable average density can be determined, however, from the mid-infrared lines alone which are insensitive to electron temperature and extinction variations. The contribution of starburst excitation to the density-sensitive forbidden lines can be estimated using the large line width variation between the NLR and the circumnuclear ring of star formation regions.
The most suitable NLR density diagnostic is provided by the ratio of the \[Ne V\] transitions at 14.32 and 24.32$`\mu `$m. These lines cannot be diluted significantly by circumnuclear star formation since they are undetected in starburst galaxies (Genzel et al. (1998)). They were observed with the same SWS aperture size and with good signal-to-noise. Adopting the same atomic data as Alexander et al. (1999, see also their Figure 3 for diagrams of several density sensitive ratios) and an electron temperature of 10000 K, the observed \[Ne V\] ratio of 1.39 corresponds to an electron density $`\mathrm{n}_\mathrm{e}2000`$ cm<sup>-3</sup> in the region of the NLR where species with lower ionization potential near 100eV prevail.
A seemingly discrepant result is obtained from the \[S III\] transitions at 18.71 and 33.48$`\mu `$m – the observed ratio of 0.73 is consistent with the low density limit and corresponds to $`\mathrm{n}_\mathrm{e}500`$ cm<sup>-3</sup>. But, Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. shows the line profiles of the two transitions to be quite different: \[S III\] 18.71$`\mu `$m shows strong wide wings and is apparently NLR-dominated with small starburst contamination. In contrast, the larger aperture of \[S III\] 33.48$`\mu `$m collects more emission from the starburst ring showing up as a strong narrow component of the profile. If all 18.71$`\mu `$m emission were from a NLR at 2000 cm<sup>-3</sup>, the NLR contribution to the 33.48$`\mu `$m flux would be about 1/3, consistent with the weaker wings of this line.
A similar problem may affect, to a lesser degree, the density sensitive ratio of the \[Ne III\] transitions at 15.55 and 36.01$`\mu `$m. The observed ratio of 8.9 is lower than but probably still consistent with the low density limit ($``$12) which applies up to the $`\mathrm{n}_\mathrm{e}2000`$ cm<sup>-3</sup> derived from \[Ne V\]. The modest signal-to-noise ratio of the 36.01$`\mu `$m line makes it impossible to use the line profile to assess starburst contamination in the large aperture. A crude estimate can be obtained assuming that the ratio of starburst \[Ne III\] 36.01$`\mu `$m and \[S III\] 33.48$`\mu `$m seen additionally in the large aperture is 0.03–0.04 as in the prototypical starburst M 82 (Förster-Schreiber (1998)). Then, $``$10% of the 36.01$`\mu `$m line would be starburst contamination, bringing the ratio closer to its low density limit value. The density in NLR regions dominated by lower excitation species like \[S III\] and \[Ne III\] hence appears consistent with that for the higher excitation region containing \[Ne V\].
With respect to the coronal line region, the observed ratio 0.55 of the \[Si IX\] lines at 2.584 and 3.936$`\mu `$m is close to its low density limit which implies $`\mathrm{n}_\mathrm{e}10^6`$ cm<sup>-3</sup>. The same limit is found for the Circinus galaxy (Moorwood et al. (1996)) and NGC 4151 (Sturm et al. (1999)). Such a density limit is consistent with all popular scenarios for coronal line formation in AGN except for origin in a very dense transition region between NLR and BLR.
## 4 Line profiles and the structure of the narrow line region
Integrated emission lines profiles are an indirect tool to constrain the dynamical structure and extinction properties of the narrow line region. Different lines probe different parts of the NLR and the velocity field is generally far from uniform. With the advent of linear optical detectors, considerable effort was devoted to studies of both forbidden and permitted optical line profiles in Seyfert galaxies. Although there is still no full consensus among different studies of the NLR forbidden lines, the emerging picture is as follows.
(1) The forbidden lines in most cases show blue asymmetries in the sense of a sharper falloff to the red than to the blue. Line centroids are blueshifted with respect to the systemic velocity, whereas line peaks in high resolution spectra are close to systemic velocity. This has been most thoroughly studied in moderate excitation species including \[O III\] 5007Å (e.g. Heckman et al. (1981); Vrtilek & Carleton (1985); Whittle 1985a ; Dahari & De Robertis (1988)) but holds also for the higher excitation coronal lines (Penston et al. (1984)).
(2) Line widths and blueshifts often vary between different species observed in the same source. Line profiles appear to be correlated with the ionization potential and/or the critical density. There are indications, but no complete consensus, that the correlation with the critical density may be the fundamental one (e.g. Pelat et al. (1981); Penston et al. (1984); Whittle 1985b ; De Robertis & Osterbrock (1986); Appenzeller & Östreicher (1988)).
Various scenarios have been put forward to explain these trends. Most of them invoke extinction to explain blue asymmetries, and the most popular ones assume outflow in a dusty NLR, with higher excitation species probably originating closer to the central source in regions of higher velocity and obscuration. It is obvious that observations of infrared NLR emission are a powerful independent method to test such scenarios: near- and mid- infrared lines suffer more than an order of magnitude less extinction than in the optical. The combination of optical and infrared data should hence elucidate the role of dust obscuration. Comparison of recombination line profiles in the optical with infrared ones would be advantageous because of the relative insensitivity of recombination line emissivities to local gas conditions. The line-to-continuum ratio of recombination lines in the infrared is low, however, and better profiles are obtained for coronal and fine structure lines which additionally cover a wide range of excitations. Sturm et al. (1999) have presented a first such analysis using ISO-SWS observations of NGC 4151. On the basis of the similarity of optical and infrared profiles, they ruled out the most simple scenario of an outflowing NLR with pervasive dust, and suggested either a geometrically thin but optically highly thick obscuring disk, or an intrinsic asymmetry of the NLR.
Because of the large flux and width of its ‘narrow’ lines, NGC 1068 is best suited for a line profile analysis at the modest resolving power ($``$2000) of ISO-SWS. Optical line profiles have been observed at very high resolving power by various groups (e.g. Pelat & Alloin (1980); Alloin et al. (1983); Meaburn & Pedlar (1986); Veilleux (1991); Dietrich & Wagner (1998)) and show complex, multi-peaked structure related to individual cloud complexes within the narrow line region of NGC 1068. Marconi et al. (1996) have extended this work into the near infrared. At lower resolving power, they do not discriminate the fine details of the best optical profiles but derive the line centroids by Gaussian fits for a large set of near-infrared and optical lines. They find that all optical and near-infrared high excitation lines are significantly blueshifted with respect to systemic velocity ($`>`$200km/s for lower ionization potential $``$ 20eV). They interpret this significant trend as a consequence of non-isotropic flows or ionization patterns rather than selective extinction effects.
We extracted line profiles for five high signal-to-noise SWS lines by subtracting a linear continuum fitted outside 2500km/s from the line center and normalizing to the peak of the line. These five lines originate in species spanning a wide range of excitation potentials ranging from 55 to 303eV and are shown in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.. The remaining uncertainty of these profiles is dominated by noise for the high excitation lines of \[Mg VIII\] and \[Si IX\] and by continuum uncertainties for the other lines. These could be both due to weak underlying real continuum features (e.g. PAH near \[Ne VI\]) and due to residual fringing (\[Ne V\] and \[O IV\]). We do not show low excitation lines with significant starburst contribution, and the \[Ne V\] 14.32$`\mu `$m and \[Ne III\] 15.55$`\mu `$m lines which are consistent with those shown in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. but more uncertain due to fringing. Brackett $`\alpha `$ has a much lower line to continuum ratio but still good signal-to-noise ratio, and a line profile similar to the fine structure lines, as discussed by Lutz et al. (2000) in the context of putting limits on a broad line region contribution.
For lines too faint to derive a good line profile, we fitted a single gaussian plus linear continuum to derive at least a centroid velocity (Table 2). While such a gaussian is not a good approximation to the intrinsic NLR profile, we adopted it for simplicity and for consistency with the optical/near-infrared data of Marconi et al. (1996). We also fitted the two-component NLR profile of Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. but do not list the derived velocities since they agree with the simple gaussian fit except for an offset that is constant within the uncertainties. From repeated observations for some of these lines, we estimate an error of $``$50km/s. For lines with no velocity listed in Table 2, we estimate that the uncertainty of deriving the centroid of a broad noisy line is too large to include it into an analysis of NGC 1068. None of them, however, is discrepant by more than $``$300km/s which would suggest misidentification.
In the following subsections, we will derive a large aperture optical NLR line profile for comparison with the ISO data, compare optical and infrared profiles and centroid velocities, and interpret the differences found.
### 4.1 A large aperture optical line profile
Mismatch between the typically small optical apertures and the large mid-IR ones is important when attempting to compare optical and mid-IR line profiles. Datacubes from imaging spectroscopy would be ideal to extract optical line profiles matching the ISO apertures. At this point, however, published imaging spectroscopy of NGC 1068 is either limited in field size (Pécontal et al. (1997)) or lacks wavelength coverage: The datacube of Cecil et al. (1990) has been obtained with 2600km/s total coverage in the \[N II\] lines that are additionally heavily blended with H$`\alpha `$, making it difficult to determine the extent to which broad components are missing in their total line profile (their Fig. 7).
We make use of two auxiliary large aperture optical spectra to address the problem of aperture mismatch: A 4000 to 7800Å spectrum from Wise Observatory (WO, S. Kaspi 1999, priv. communication), providing good fluxes of the brightest lines in a $`10\mathrm{}\times 15\mathrm{}`$ aperture (position angle 0°), and a high spectral resolution Coudé Echelle spectrum from Karl Schwarzschild Observatory Tautenburg (KSO, E. Guenther 1999, priv. communication), providing a good \[O III\] line profile (though not good fluxes) in a $`6.8\mathrm{}\times 15\mathrm{}`$ aperture (mean position angle -28°, varying during integration). In addition, we estimated relative emission line fluxes in our apertures by integrating over the corresponding regions of a narrow band \[O III\] map (R. Pogge, M.M. deRobertis, 1999, unpublished data). Both the Wise spectrum and the \[O III\] map confirm that ISO line fluxes of the NGC 1068 NLR can be sensibly compared to smaller aperture optical data, since those already sample most of the flux in the narrow line region. For example, a 4″ diameter aperture will already get $``$70% of the flux in the ISO aperture. The fluxes measured in the large Wise aperture for the brightest optical lines agree within $``$30% with published smaller aperture ones (e.g. Shields & Oke (1975); Koski (1978); Marconi et al. (1996)).
Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. displays our KSO $`6.8\mathrm{}\times 15\mathrm{}`$ aperture \[O III\] 5007Å line profile in comparison to its $`2.5\mathrm{}\times 2.5\mathrm{}`$ equivalent (Veilleux (1991)). The line profile changes induced by this more than tenfold increase in aperture area are relatively modest and fit expectations from high resolution longslit spectroscopy. While the major components of Veilleux’ spectrum are well reproduced in the KSO data, their ratios differ somewhat leading to an overall slightly wider profile. This is fully consistent with observations of relatively broad components over larger regions not sampled by Veilleux’ aperture (Pelat & Alloin (1980); Alloin et al. (1983); Meaburn & Pedlar (1986)). The only feature in the KSO profile not present in the Veilleux profile is an additional narrow feature at or slightly redshifted from systemic velocity. This feature almost certainly corresponds to the ‘velocity spike’ in the NE region of the NLR detected by many authors but seen perhaps most clearly in the data of Meaburn & Pedlar (1986). This feature is missed by Veilleux’ aperture but partly covered by the KSO data.
The KSO aperture is still about three times smaller in area than the SWS apertures through which the best fine structure line profiles have been taken. The drop in \[O III\] surface brightness with radius is so rapid (e.g. Fig. 1 of Meaburn & Pedlar 1986) that only modest differences in the total line profile are expected. An exception to this is the NE region of the NLR about 6″ from the nucleus which was incompletely covered. The KSO slit orientation cannot be chosen freely and was approximately aligned with the ISO apertures but not with the NLR (PA -28° instead of PA $``$30°), missing part of the NE end of the NLR. Long slit spectroscopy (e.g. Meaburn & Pedlar (1986)) shows this NE region to be dominated by the narrow ‘velocity spike’ near systemic velocity which is already seen in the comparison of KSO and Veilleux (1991) profiles. We hence expect this spike to be more prominent in an optical line profile fully equivalent to the ISO aperture. Integrating the \[O III\] map over the ISO and KSO apertures we estimate a need to add $``$11% to the KSO flux to account for the NE region and other low surface brightness emission near systemic velocity. We have taken this into account by adding such a narrow (FWHM 150km/s) component to the KSO profile, and will use this in the following as basis of our optical line profile comparison (see also Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.). Use of such a modified profile is supported by the spectrum of Pelat & Alloin (1980) which was obtained with a rotating longslit sweeping across the NE region of the NLR, and showing a similar narrow spike (their component 5).
### 4.2 Line profile variations
The optical and infrared line profiles in NGC 1068 differ strongly. This is most evident in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. which compares the profiles of \[O IV\] 25.89$`\mu `$m and \[O III\] 5007Å. We chose \[O IV\] as the representative infrared line since it was observed with very good S/N and is close to \[O III\] in lower ionization potential of the emitting species (55 vs. 35eV), ensuring origin in a similar region of the NLR. Critical densities ($`1.0\times 10^4`$ cm<sup>-3</sup> vs. $`7.0\times 10^5`$ cm<sup>-3</sup>) match less well than \[Ne V\] or \[Ne VI\] would, but this is less relevant given the low NLR density we have inferred (see §3.1).
The optical \[O III\] line profile is both blueshifted and broader than the infrared \[O IV\] line profile. There are however also significant similarities. Shoulders near -900 km/s and $``$350 km/s are present in both profiles. Adopting different relative scalings (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.), the impression arises that the two profiles in fact agree fairly well over parts of their extent if the scaling is set properly. The main difference lies in different relative strengths of blue wing, center, and red wing, the infrared profile having a stronger red wing and center.
The four highest quality infrared profiles are compared in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. using a spread velocity scale. The most obvious and significant variation is in \[Mg VIII\] which is both broader and more blueshifted than the lower excitation lines. The same is observed for the more noisy \[Si IX\] line. As already noted, the \[O IV\],\[Ne V\], and \[Ne VI\] profiles with excitation energies ranging from 55 to 126eV are very similar but the overplot shows some variation in detail. There is a minor shift in the narrow core of the \[Ne VI\] line which is, however, not significant compared to the quoted systematic uncertainty. Comparing \[Ne V\] to \[O IV\] taken from the same observation there is even less shift. Concerning the broader wings, there are significant differences in addition to the possible presence of \[Fe II\] at $``$1100km/s in the \[O IV\] profile. There is a trend from \[O IV\] to \[Ne VI\] in the blue wing becoming stronger and the red wing fainter.
Such profile variations determine the line centroids derived from gaussian fits (Table 2) to which we add a centroid of 1015 km/s derived in the same way for our extrapolated large aperture optical \[O III\] profile. Anticipating that the shifts may reflect several partially degenerate influences, we show in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. centroid velocities as a function of lower ionization potential, critical density, and extinction for the particular line. The extinction values are relative and based on a preliminary ISO-based extinction curve for the center of our Galaxy (Lutz et al. 1997a ).
The trend observed in the ISO data for the velocity centroids as a function of ionization potential (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.) is markedly different from the equivalent dataset obtained in the optical and near-infrared (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA., data taken from Marconi et al. 1996). In both data sets the velocity centroids of the lowest excitation non-NLR lines (H<sub>2</sub>, \[Fe II\]) are close to systemic. However, none of the ISO lines reach the large blueshifts observed consistently over a wide excitation range in the optical and near-infrared. The data sets are least discrepant at the high excitation end. Here, \[Si IX\] 3.936$`\mu `$m is common to both sets and agrees within the errors, though being somewhat less blueshifted in the ISO data. The discrepancy is largest at intermediate excitation (20-200eV) where the optical/NIR lines are all strongly blueshifted while the ISO lines only slowly deviate from systemic velocity as excitation increases.
### 4.3 Interpretation of the profile differences
In addition to the cloud distribution and kinematics, line profiles reflect the emissivities of fine structure or forbidden lines in the narrow line region clouds, which are affected by many parameters: ionization equilibrium, density, electron temperature, and extinction. If any of these parameters vary among kinematically distinct structural components of the NLR, line profile variations will result that in turn can help elucidate the NLR structure. An important aspect is that some of these parameters are partially degenerate in infrared datasets: Shorter wavelength (2-5$`\mu `$m) lines which suffer higher extinction are also typically high excitation coronal lines with high critical densities, whereas the longer wavelengths are dominated by ions of lower excitation with transitions with lower critical densities.
Here we have assumed that the observed lines are emitted locally, with no contribution of scattered light. This is not strictly correct for part of the NGC 1068 narrow line region, in particular the NE region (Capetti et al. (1995); Inglis et al. (1995)). The impact of scattering on our analysis depends on the properties of the scatterer. If scattering is wavelength-independent, our profile comparison is unaffected since scattering would effectively only redistribute the line emission spatially within our large apertures. If (dust) scattering decreases with wavelength, shorter wavelength line profiles would be modified more strongly. We estimate the effect on the integrated line profiles is not large, since the \[O III\] polarisations measured by Inglis et al. (1995) reach at most a few percent in regions that are in addition minor contributors to the total flux, and since the considerable spatial variation in line profiles is not suggestive of scattered radiation originating in a central source.
It is unlikely that density variations play a direct role in the profile variations. Our density estimate of $`\mathrm{n}_\mathrm{e}2000\mathrm{c}\mathrm{m}^3`$ is too far below the relevant critical densities, making much higher densities over a significant part of the NLR unlikely, which would be required to create the variations. Strictly speaking however, this estimate applies only to the moderately excited (100eV) gas, and our upper limit for the coronal region density is still consistent with densities higher than the critical densities of some lower excitation line. Strong collisional suppression of part of some low critical density fine structure line profiles is also unlikely from the similarity of their line centroids to that of the recombination line Brackett $`\alpha `$. We hence believe that the clear correlation of centroid velocities with fine structure line critical density (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.b) is mostly a secondary consequence of the correlation with ionization potential.
Variations in electron temperature can also affect the line profiles. Emissivities of optical forbidden lines like \[O III\] 5007Å strongly vary with electron temperature while the infrared fine structure lines originate close to the ground state and are much less sensitive to temperature. Hence, at least those profile variations seen among the infrared lines (Figures ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA. and ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.) must be unrelated to temperature fluctuations. For the optical/IR profile variations there is a basic ambiguity of an optical component being faint due to low electron temperature or due to high extinction. Optical electron temperature determinations e.g. from \[O III\] 4363Å/5007Å are not available for the various kinematic components of NGC 1068. The line profiles of the temperature insensitive recombination lines in the optical and IR do not resemble each other, but rather follow the shapes of the forbidden optical and IR lines, respectively (Veilleux (1991), Fig. ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.). This is inconsistent with temperature variations being the main origin of profile variations.
The exclusion of other factors suggests that ionization structure and extinction are the main source of the observed variation of optical/IR line profiles. Previous studies of the spatial and kinematical structure of the NGC 1068 NLR (e.g. Cecil et al. (1990), Marconi et al. (1996)) point to the existence of two spatial and kinematical components. The first is a strong ionization cone with associated blueshifted outflow, with the highest excitation species likely concentrated towards the central and fastest part of this cone. The second is an extended system of photoionized clouds closer to systemic velocity. The correlation between the ISO centroid velocities and the ionization potential (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.a) fits this picture well. If the relative importance of the fast outflow gradually increases towards high excitation lines, the gradual centroid shift is easily explained. However, the marked differences between optical \[O III\] and infrared \[O IV\] profile (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.) and the different optical and infrared centroid velocities at similar ionization potential (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.a, including data from Marconi et al. 1996) show that this picture must be incomplete. We suggest that these remaining differences are due to extinction variations across the NLR, with a general trend of higher exctinction in the redshifted (SW) than in the blueshifted (NE) part. The optical profile can be explained by an intrinsic profile similar to that of the IR lines whose line center and red wing are reddened by a few magnitudes, leading to the obscuration of about half of the total line flux. A similar extinction pattern is suggested by the increase of polarization from the blue to the red wing of the narrow lines in the central arcseconds (Antonucci & Miller (1985); Bailey et al. (1988); Inglis et al. (1995)), attributed to absorption by aligned grains.
Considering that the optical profiles are modified by extinction, it is important to recognize that the less obscured near- and mid-infrared lines remain blueshifted with respect to systemic velocity, high excitation ones more strongly than lower excitation ones. Also, coronal line emission observed in the little obscured near-infrared is still much stronger in the northeast cone than in the southwest (Thompson & Corbin (1999)). This might be due to heavy extinction of an intrinsically symmetric NLR, or due to a real asymmetry. No distinction can be made from the present data. A trend with IR extinction is not clear in Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.c and would be difficult to separate from the trends with ionization potential and critical density because of the mentioned degeneracy. An intrinsically asymmetric NLR with outflow preferentially towards the observer is fully consistent with the observations of NGC 1068. The difference in strength of the two radio lobes (Wilson & Ulvestad (1983)) may be in support of such an asymmetric scenario, although the relation between radio lobe flux and NLR emission is certainly not simple. Consistency with an asymmetric NLR was also noted for the ISO spectra of NGC 4151 (Sturm et al. (1999)). However, large randomly oriented AGN samples should not show the preferential blueshift which is noted at least in optical samples. An alternative scenario of a geometrically small but optically highly thick screen (disk or torus) obscuring part of the receding NLR region was proposed by Sturm et al. (1999) for NGC 4151. Such a scenario can also fit the NGC 1068 data provided the screen obscures a relatively larger fraction of the coronal line region than of the larger NLR which is dominated by medium excitation species. This is not implausible given the similar (arcsecond) spatial scale of the central concentration of high column density molecular gas (e.g. Tacconi et al. (1994)) and the coronal line region as mapped in \[Si VI\] (Thompson & Corbin (1999); Thatte et al. (2000)). While studies of few individual sources will remain ambiguous, a search for preferential shifts using high resolution near- or mid-infrared spectroscopy should address this issue provided the sample of Seyferts is large enough and not biased by orientation of a putative asymmetric outflow, as may be the case when identifying Seyferts in the optical.
Overall, the structure of the NGC 1068 narrow line region seems to determine the optical/infrared line profiles via differences in weight of outflowing cone and extended components, and additional extinction variations that are significant for the optical wavelength range. The relative weight of the cone is larger for higher excitation species. A tantalizing ambiguity remains that cannot be resolved from a single source study: Is the NLR intrinsically one-sided or is there a very high obscuration screen blocking also part of the IR emission from our view?
## 5 Molecular hydrogen emission
NGC 1068 was the first galaxy detected in the rovibrational transitions of molecular hydrogen (Thompson et al. (1978)) and has been studied since in considerable detail both in these lines tracing fairly excited molecular material (e.g. Oliva & Moorwood (1990); Blietz et al. (1994); Davies et al. (1998)), and by millimeter wave interferometry tracing colder components (e.g. Tacconi et al. (1994); Sternberg et al. (1994); Helfer & Blitz (1995); Schinnerer (1999)). The system of dense, warm cloud cores in the central few arcseconds inferred from the millimeter studies calls for observations in the pure rotational transitions of molecular hydrogen, which trace gas of typically a few 100K, intermediate between the near-infrared and millimeter wave tracers.
Our SWS observations of these lines (Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA., Table 2) are summarized in the H<sub>2</sub> excitation diagram of Figure ISO-SWS spectroscopy of NGC 1068 <sup>1</sup><sup>1</sup>1Based on observations with ISO, an ESA project with instruments funded by ESA Member States (especially the PI countries: France, Germany, the Netherlands and the United Kingdom) with the participation of ISAS and NASA.. The diagram has a curved shape suggestive of a mixture of temperatures, as expected from other galaxies observed with ISO in the rotational lines of molecular hydrogen (Rigopoulou et al. (1996); Sturm et al. (1996); Valentijn et al. 1996a ; Kunze et al. (1996); Spoon et al. (2000)). The location of the S(3) point slightly below the general trend suggests a moderate extinction towards this line whose wavelength is near the center of the silicate absorption feature ($`\mathrm{A}_{9.6\mu \mathrm{m}}1`$). As noted earlier, the flux for the heavily blended S(7) line is very uncertain, so that the shorter wavelength rovibrational lines of similar excitation will probably be a more trustworthy representation of the excitation diagram at these upper level energies.
While the higher rotational lines like S(5) and S(7) probe the same excited but low mass component as the near-infrared rovibrational lines, the bulk of the warm gas observed with ISO will reside in the component traced by the S(1) line. In estimating its mass, we will assume a hydrogen ortho/para (O/P) ratio in equilibrium at the local temperature (Sternberg & Neufeld (1999)). Combining the S(1) flux with the S(0) limit on one hand and the S(3) detection on the other, the temperature of the S(1) emitting gas is found to lie in the range 140K$``$T$``$375K, assuming the extinction correction for S(3) is modest and aperture effects are minor. Because of the very steep temperature sensitivity of the rotational line emissivities, the corresponding mass varies drastically, between about $`4\times 10^6\mathrm{M}_{\mathrm{}}`$ for the 375 K case and $`1.5\times 10^8\mathrm{M}_{\mathrm{}}`$ for the 140 K case. For comparison with other mass estimates, we will adopt $``$200 K and $`2.5\times 10^7\mathrm{M}_{\mathrm{}}`$ as a possible approximation to the curvature of the excitation diagram. This mass would be of the order 5% of the total gas mass estimated from the CO interferometric map (Helfer & Blitz (1995) and priv. comm.) within the 17$`\mu `$m ISO beam. Compared to a similar estimate for the starburst galaxy NGC 3256 (Rigopoulou et al. (1996), 3% for 150K warm gas temperature), the fraction of warm gas and/or its temperature must be higher, but not exceeding on average that inferred by Kunze et al. (1996) for the highly active star forming region in the ‘overlap region’ of the antenna galaxies (8% for 200K warm gas temperature).
The NGC 1068 molecular hydrogen emission sampled by SWS may represent a mixture of relatively cool gas from the 15″ radius molecular ring partly covered by the longer wavelength apertures, and a warmer component from the unusually dense and warm central few arcseconds (Tacconi et al. (1994); Blietz et al. (1994)). Knowing from the S(3) measurement that the extinction to the H<sub>2</sub> emitting region cannot be very large, we can compare our 1-0 Q(3) flux to 1-0 S(1) fluxes measured in smaller apertures (1.4 to 2.2$`\times 10^{20}`$W cm<sup>-2</sup>; Blietz et al. (1994); Thompson et al. (1978)). Our Q(3) line flux falls in the middle of that range. Since the intrinsic Q(3)/S(1) flux ratio is 0.7, the implication is that the central few arcseconds dominate at least for the more highly excited hydrogen lines, though there may be some extended contribution. Some of the molecular hydrogen emission in NGC 1068 may originate in X-ray irradiated gas. We include in Table 2 upper limits for two transitions of H$`{}_{}{}^{+}{}_{3}{}^{}`$ which have been proposed as a signature of X-ray heated molecular gas (Draine & Woods (1990)). We note, however, that these limits are not at all stringent and are fully consistent with even the ‘high’ end of H$`{}_{}{}^{+}{}_{3}{}^{}`$ flux expectations for X-ray illumination. Estimates are uncertain, see e.g. the much lower predictions of Maloney et al. (1996).
The question can be raised whether the NGC 1068 H<sub>2</sub> spectra in fact contain a direct signature of a parsec-scale molecular torus. Krolik and Lepp (1989) have modelled molecular line emission of such an X-ray illuminated torus, predicting emission in some molecular hydrogen lines like (0-0) S(5) that may under favorable conditions be detectable at ISO-SWS sensitivities. For NGC 1068 it is evident that larger scale emission may swamp any possible torus emission, as already cautioned by Krolik and Lepp (1989). The ISO data smoothly complement the near-infrared emission originating in larger scale ($``$100pc) clouds, following an excitation diagram plausibly ascribed to the same clouds. If any torus emission were present at lower level, it may be difficult to discriminate from the larger scale emission since, depending on black hole mass and spatial scale, its velocity width could be very similar to the larger scale emission. The widths of our H<sub>2</sub> rotational lines are consistent with those for CO observed on 100pc and larger scales, with no evidence for other kinematic components. However, even if the rotational emission does not trace a compact torus, it may still be excited to a significant fraction by UV radiation, X-rays or shocks related to the AGN.
## 6 Summary
ISO-SWS spectroscopy provides the first detailed census of the mid-infrared spectrum of the prototypical Seyfert 2 galaxy NGC 1068. We have detected 36 emission lines on top of the strong AGN-heated continuum. Most lines originate in the NLR characterized by a density of $``$2000 cm<sup>-3</sup>.
We have compared the mid-infrared ISO line profiles with optical emission line profiles produced in the NLR. The line profiles are consistent with a model where the NLR is a combination of a highly ionized outflow and lower excitation extended emission, with extinction significantly affecting the optical line profiles. Remaining blueshift and asymmetry of the least obscured lines may reflect either intrinsic asymmetry of the NLR or an additional very high column density obscuring component.
We detect strong emission from warm molecular hydrogen, which most likely originates on the 100pc to kpc scale, and which is also probed by emission in near-infrared and millimeter wave tracers of molecular material. This emission masks any possible emission from a putative parsec-scale molecular torus.
Companion papers use the SWS data to model the spectral energy distribution of the active nucleus, to put limits on emission from the obscured broad line region, and discuss the continuum and its features.
We are grateful to Eike Guenther and Shai Kaspi for obtaining optical spectra that were invaluable in the interpretation of the ISO spectra, and to Richard Pogge and M.M. de Robertis for providing us with an unpublished \[O III\] image of NGC 1068. SWS and the ISO Spectrometer Data Center at MPE are supported by DLR (DARA) under grants 50 QI 8610 8 and 50 QI 9402 3. We acknowledge support by the German-Israeli Foundation (grant I-0551-186.07/97). |
warning/0002/hep-ph0002195.html | ar5iv | text | # I Introduction
## I Introduction
In recent years the study of the deeply virtual Compton scattering (DVCS) became one of the most popular topics in QCD due to the fact that it is determined by skewed parton distributions \[1-3\] which generalize usual parton densities introduced by Feynman. These new probes of the nucleon structure are accessible in exclusive processes such as DVCS and potentially they can give us more information than the traditional parton densities. In this paper we consider the small-x DVCS where the energy of the incoming virtual photon $`E`$ is very large in comparison to its virtuality $`Q^2`$. (The first study of the small-x DVCS was undertaken in Ref. ). To be specific, we calculate the DVCS amplitude in the region
$$sQ^2tm^2$$
(1)
where $`s=2mE`$, $`m`$ is the nucleon mass, and $`t`$ is the momentum transfer. The DVCS in this region is a semihard processes which can be described by the BFKL (Balitsky-Fadin-Kuraev-Lipatov) pomeron . It turns out that at large momentum transfer the coupling of the BFKL pomeron to the nucleon is essentially equal to the Dirac form factor of the nucleon $`F_1(t)`$, so the DVCS amplitude in the region (1) can be calculated without any model assumptions. The results obtained in this region can be used for the estimates of the amplitude at experimentally accessible energies where one or more conditions in Eq. (1) are relaxed. To be specific, we have in mind the HERA kinematics where $`x10^2÷10^4`$, $`Q^26`$ GeV<sup>2</sup>, and $`t1÷5`$ GeV<sup>2</sup> . Since there are only model predictions for the small-x DVCS in current literature , even the approximate calculations of the cross section in QCD are very timely.
## II Small-x DVCS in the lowest order in perturbation theory
Similarly to the case of deep inelastic scattering (DIS), the amplitude of DVCS is determined by the matrix element
$$H^{AB}=ie_\nu ^Ae_\mu ^B𝑑ze^{iq^{}z}p^{}|T\{j^\mu (z)j^\nu (0)\}|p$$
(2)
where $`q,p`$ and $`q^{},p^{}`$ are the initial and the final momenta of the photon and the nucleon, respectively. The momentum transfer is defined as $`r=p^{}p`$. Since $`Q^2=q^2`$ is large we can use perturbation theory for the hard part of the DVCS process . The typical diagram for the DVCS amplitude in the lowest order in perturbation theory is shown in Fig.1 (recall that the diagrams with gluon exchanges dominate at high energies).
It is convenient to calculate at first the imaginary part of the amplitude $`H^{AB}`$
$$V^{AB}=\frac{1}{\pi }\mathrm{Im}T^{AB}.$$
(3)
In the leading order in perturbation theory the amplitude at high energy is purely imaginary up to the $`\frac{Q^2}{s}`$ corrections (see e.g. the review ). At high orders in perturbation theory the amplitude will be purely imaginary in the leading logarithmic approximation (LLA) and we will restore the real part using the dispersion relations.
At high energies it is convenient to use the Sudakov variables. Let us define the light-like vectors $`p_1=q^{}`$, $`p_2=p^{}\frac{m^2}{s}p_1`$, then
$`q`$ $`=`$ $`p_1(1{\displaystyle \frac{r_{}^2}{s}})xp_2r_{}q^{}=p_1`$ (4)
$`p`$ $`=`$ $`p_2(1+x)+{\displaystyle \frac{m^2+r_{}^2}{s}}p_1+r_{}p^{}=p_2+{\displaystyle \frac{m^2}{s}}p_1`$ (5)
where $`x\frac{Q^2+t}{s}\frac{Q^2}{s}=x_{Bj}`$ and $`tr_{}^2`$ at large energies. Consider the integral over gluon momentum $`k=\alpha _kp_1+\beta _kp_2+k_{}`$
$$V^{AB}=\frac{2}{\pi }g^4\frac{d^4k}{16\pi ^4}\frac{1}{k^2}\frac{1}{(r+k)^2}\mathrm{Im}\mathrm{\Phi }_{\xi \eta }^{ab}(k+r,k)\mathrm{Im}\mathrm{\Phi }_N^{\xi \eta ab}(kr,k)$$
(6)
where $`\mathrm{\Phi }_{\xi \eta }^{ab}(k,r+k)`$ and $`(\mathrm{\Phi }_N)_{\xi \eta }^{ab}(k,r+k)`$ are the upper and the lower blocks of the diagram in Fig. 2 (stripped of the strong coupling constant $`g`$). Here $`a,b`$ and $`\xi ,\eta `$ are the color and Lorentz indices, respectively.
It is well known that in the Regge kinematics ($`s`$ everything else) $`\alpha _k\frac{m^2}{s}`$, and $`\beta _kx`$ so $`k^2k_{}^2`$. Moreover, $`\alpha `$’s in the upper block are $`1`$ so one can drop $`\alpha _k`$ in the upper block. Similarly, $`\beta `$’s in the lower block are $`1`$ and one can neglect $`\beta _k`$ in the lower block. We get ($`\mathrm{\Phi }^{ab}=\frac{\delta _{ab}}{8}\mathrm{\Phi }^{cc}`$):
$$V^{AB}=\frac{g^4}{4\pi }\frac{d^4k}{16\pi ^4}\frac{1}{k_{}^2}\frac{1}{(r+k)_{}^2}\mathrm{Im}\mathrm{\Phi }_{\xi \eta }^{aa}(k+r,k)|_{\alpha _k=0}\mathrm{Im}\mathrm{\Phi }_N^{\xi \eta bb}(kr,k)|_{\beta _k=0}.$$
(7)
At high energies, the metric tensor in the numerator of the Feynman-gauge gluon propagator reduces to $`g^{\mu \nu }\frac{2}{s}p_2^\mu p_1^\nu `$ so the integral (7) for the imaginary part factorizes into a product of two “impact factors” integrated with two-dimensional propagators
$$V^{AB}=\frac{2s}{\pi }g^4\left(e_q^2\right)\frac{d^2k_{}}{4\pi ^2}\frac{1}{k_{}^2}\frac{1}{(r+k)_{}^2}I(k_{},r_{})I_N(k_{},r_{})$$
(8)
where
$`I(k_{},r_{})`$ $`=`$ $`{\displaystyle \frac{1}{2s}}p_2^\xi p_2^\eta \left({\displaystyle e_q^2}\right){\displaystyle \frac{d\beta _k}{2\pi }\mathrm{Im}\mathrm{\Phi }_{\xi \eta }^{aa}(k+r,k)}|_{\alpha _k=0}`$ (9)
$`I_N(k_{},r_{})`$ $`=`$ $`{\displaystyle \frac{1}{2s}}p_1^\xi p_1^\eta {\displaystyle \frac{d\alpha _k}{2\pi }\mathrm{Im}\mathrm{\Phi }_{N\xi \eta }^{aa}(kr,k)}|_{\beta _k=0}`$ (10)
and $`\left(e_q^2\right)`$ is the sum of squared charges of active flavors ($`u,d,s`$, and possibly $`c`$). The photon impact factor is given by the two one-loop diagrams shown in Fig. 3.
The standard calculation of these diagrams yields
$$I^{AB}(k_{},r_{})=\overline{I}^{AB}(k_{},r_{})\overline{I}^{AB}(0,r_{})$$
(11)
where
$`\overline{I}^{AB}(k_{},r_{})={\displaystyle \frac{1}{2}}{\displaystyle _0^1}{\displaystyle \frac{d\alpha }{2\pi }}{\displaystyle _0^1}{\displaystyle \frac{d\alpha ^{}}{2\pi }}\left\{P_{}^2\alpha ^{}\overline{\alpha ^{}}+Q^2\alpha ^{}\alpha \overline{\alpha }\right\}^1`$ (12)
$`\left\{(12\alpha \overline{\alpha })P_{}^2(e^A,e^B)_{}+4\alpha \overline{\alpha }\overline{\alpha ^{}}[P_{}^2(e^A,e^B)2(e^A,P)_{}(e^B,P)_{}]4\alpha \overline{\alpha }(12\alpha )(r,e^A)_{}(P,e^B)_{}\right\}`$ (13)
for the transverse polarizations $`A,B=1,2`$ (cf. ) and
$`\overline{I}^{3B}(k_{},r_{})={\displaystyle \frac{1}{2Q}}{\displaystyle _0^1}{\displaystyle \frac{d\alpha }{2\pi }}{\displaystyle _0^1}{\displaystyle \frac{d\alpha ^{}}{2\pi }}\left\{P_{}^2\alpha ^{}\overline{\alpha ^{}}+Q^2\alpha ^{}\alpha \overline{\alpha }\right\}^1`$ (14)
$`\left\{(12\alpha \overline{\alpha })P_{}^2(r,e^B)_{}+4\alpha \overline{\alpha }\overline{\alpha ^{}}[P_{}^2(r,e^B)_{}2(r,P)_{}(e^B,P)_{}]4\alpha \overline{\alpha }(12\alpha )Q^2(P,e^B)_{}\right\}`$ (15)
for the longitudinal polarization
$$e^3(q)=\frac{1}{Q}(p_1+xp_2)$$
(16)
. Here $`P_{}k_{}+r_{}\alpha `$ and $`(a,b)_{}`$ denotes the (positive) scalar product of transverse components of vectors $`a`$ and $`b`$. At large transverse momenta $`k_{}^2r_{}^2`$ the impact factor (11) reduces to
$$I^{AB}(k_{},r_{})\frac{(e^A,e^B)_{}}{4\pi ^2}\frac{k_{}^2}{Q^2}\mathrm{ln}\frac{Q^2}{r_{}^2}.$$
(17)
The impact factor for the proton which decribes the pomeron-nucleon coupling cannot be calculated in the perturbation theory. However, in the next section we demonstrate that at high momenta $`k_{}^2m^2`$ this impact factor reduces to
$$I_N(k_{},r_{})\stackrel{k_{}^2m^2}{=}F_1^{p+n}(t)$$
(18)
where $`F_1^{p+n}(t)`$ is the sum of the proton and neutron Dirac form factors. As we shall see below, the characteristic transverse momenta in our gluon loop are large so the estimate (18) is sufficient for our purposes. Substituting the nucleon impact factor (18) into Eq. (8) we obtain
$$V^{AB}=\frac{2s}{\pi }g^4(e_q^2)F_1^{p+n}(t)\frac{d^2k_{}}{4\pi ^2}\frac{I^{AB}(k_{},r_{})}{k_{}^2(r+k)_{}^2}.$$
(19)
Performing the final integration over $`k_{}`$, one gets
$`V^{AB}=`$ (22)
$`{\displaystyle \frac{2}{x}}\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)F_1^{p+n}(t)`$
$`\left((e^A,e^B)_{}\left({\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{|t|}}+2\right)(e^A,e^B)_{}+{\displaystyle \frac{2}{r_{}^2}}(e^A,r)_{}(e^B,r)_{}+O(t/Q^2)\right)`$
for the transverse polarizations and
$`V^{3B}=`$ (24)
$`{\displaystyle \frac{2}{x}}\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)F_1^{p+n}(t){\displaystyle \frac{(r,e^B)_{}}{Q}}\left({\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{|t|}}5\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}+{\displaystyle \frac{15}{2}}{\displaystyle \frac{\pi ^2}{3}}+O(t/Q^2)\right)`$
for the longitudinal one. The longitudinal amplitude (24) is twist-suppressed as $`\frac{\sqrt{|t|}}{Q}`$ in comparison to the transverse amplitude (22) (as it should, due to the fact that $`t0`$ corresponds to real incoming photon).
Since the integral over $`k_{}`$ (19) converges at $`k_{}Q`$ the region $`k_{}m`$, where we do not know the nucleon impact factor, contributes to the terms $`O(t/Q^2)`$ which we neglect.
## III Nucleon impact factor
In the lowest order in perturbation theory there is no difference between the diagrams for the nucleon impact factor shown in Fig. 4
and similar diagrams with two gluons replaced by two photons (up to the trivial numerical factor $`c_F=\frac{4}{3}`$ and replacement of $`eg`$). In this case the lower part of the diagram can be formally written as follows:
$$\mathrm{\Phi }_N(kr,k)\stackrel{\mathrm{def}}{}\frac{1}{2}\frac{p_1^\xi p_1^\eta }{s}(\mathrm{\Phi }_N)_{\xi \eta }^{bb}(kr,k)=\frac{2}{3}ip_1^\mu p_1^\nu 𝑑ze^{ikz}p^{}|T^{}\{J_\mu (z)J_\nu (0)\}|p$$
(25)
where $`J_\mu =\overline{u}\gamma _\mu u+\overline{d}\gamma _\mu d`$. The $`T^{}`$ means the T-product where the diagrams with pure gluon exchanges in t-channel are excluded; by definition, such diagrams contribute to subsequent ranks of BFKL ladder rather than to impact factor. (This is the reason why we have not included in $`J`$ the contribution of strange quarks). Since $`k^2`$ in our case is large and negative (-$`k^2=k_{}^2m^2`$) we can expand the T-product of two currents near the light cone (see e.g. )
$$\mathrm{\Phi }_N(k,r+k)=\frac{2}{3s}𝑑ze^{ikz}\frac{zp_1}{\pi ^2z^4}p^{}|\overline{\psi }(z)[z,0]\overline{)}p_1\psi (0)+\overline{\psi }(0)[0,z]\overline{)}p_1\psi (z)|p_{z^2=0}^{}$$
(26)
where again $`\mathrm{}^{}`$ stands for the matrix element with pure gluon exchanges excluded. Here $`[x,y]`$ denotes the gauge link connecting the points $`x`$ and $`y`$ ($`[x,y]Pexp(ig_0^1du(xy)^\mu A_\mu (ux+(1u)y)`$). The matrix element (19) can be parametrized in terms of skewed parton distributions as follows
$`p^{},\lambda ^{}|\overline{q}(z)[z,0]\overline{)}p_1q(0)|p,\lambda _{z^2=0}^{}=`$ (27)
$`\overline{u}(p^{},\lambda ^{})\overline{)}p_1u(p,\lambda ){\displaystyle _0^1}𝑑Xe^{i(Xx)pz}𝒱_x^q(X,t)+{\displaystyle \frac{1}{2m}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda ){\displaystyle _0^1}𝑑Xe^{i(Xx)pz}𝒲_x^q(X,t)`$ (28)
$`p^{},\lambda ^{}|\overline{q}(0)[0,z]\overline{)}p_1q(z)|p,\lambda _{z^2=0}^{}=`$ (29)
$`\overline{u}(p^{},\lambda ^{})\overline{)}p_1u(p,\lambda ){\displaystyle _0^1}𝑑Xe^{iXpz}𝒱_x^q(X,t)+{\displaystyle \frac{1}{2m}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda ){\displaystyle _0^1}𝑑Xe^{iXpz}𝒲_x^q(X,t),`$ (30)
where $`𝒱_x^u(X,t)`$ and $`𝒲_x^u(X,t)`$ are the nonflip and spin-flip skewed parton distributions for the valence $`u`$ quark (recall that we must take into account only valence quarks since we forbid diagrams with pure gluon exchanges). Similarly, $`𝒱_x^d(X,t)`$ and $`𝒲_x^d(X,t)`$ refer to the valence $`d`$-quark distributions. At large energies $`\overline{u}(p^{},\lambda ^{})\overline{)}p_1u(p,\lambda )=s\delta _{\lambda \lambda ^{}}`$, so
$`p^{},\lambda ^{}|\overline{q}(0)[0,z]\overline{)}p_1q(z)\overline{q}(z)[z,0]\overline{)}p_1q(0)|p,\lambda _{z^2=0}^{}=`$ (32)
$`{\displaystyle _0^1}𝑑X\left(e^{iXpz}e^{i(Xx)pz}\right)\left[s\delta _{\lambda \lambda ^{}}𝒱_x^q(X,t)+{\displaystyle \frac{1}{2m}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda )𝒲_x^q(X,t)\right].`$
After integration over $`z`$ the lower block (25) reduces to
$`\mathrm{\Phi }_N(kr,k)=`$ (35)
$`{\displaystyle \frac{2}{3s}}{\displaystyle _0^1}𝑑X\left[{\displaystyle \frac{(Xx)s+2p_1k}{k^22pk(Xx)iϵ}}{\displaystyle \frac{Xs+2p_1k}{k^2+2pkXiϵ}}\right]`$
$`\left(\delta _{\lambda \lambda ^{}}(𝒱_x^u(X,t)+𝒱_x^d(X,t))+{\displaystyle \frac{1}{2ms}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda )(𝒲_x^u(X,t)+𝒲_x^d(X,t))\right).`$
The nucleon impact factor (10) is the integral of the imaginary part of r.h.s. of eq. (35) over energy
$`I_N(k_{},r_{})={\displaystyle _0^1}{\displaystyle \frac{d\alpha _k}{2\pi }}\mathrm{Im}\mathrm{\Phi }_N((\alpha _k{\displaystyle \frac{r_{}^2}{s}})p_1k_{}r_{},\alpha _kp_1+k_{})=`$ (39)
$`{\displaystyle \frac{1}{3}}{\displaystyle _0^1}𝑑\alpha _k{\displaystyle _x^1}𝑑X\left[s(Xx)\delta (k_{}^2\alpha _ks(Xx))sX\delta (k_{}^2+\alpha _ksX)\right]`$
$`\left(\delta _{\lambda \lambda ^{}}(𝒱_x^u(X,t)+𝒱_x^d(X,t))+{\displaystyle \frac{1}{2ms}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda )(𝒲_x^u(X,t)+𝒲_x^d(X,t))\right)`$
$`={\displaystyle \frac{1}{3}}{\displaystyle _x^1}𝑑X\left(\delta _{\lambda \lambda ^{}}(𝒱_x^u(X,t)+𝒱_x^d(X,t))+{\displaystyle \frac{1}{2ms}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda )(𝒲_x^u(X,t)+𝒲_x^d(X,t))\right).`$
Since valence quark distributions decrease at $`x0`$ we can extend the lower limit of integration in r.h.s. of eq. (39) to 0 and obtain
$`I_N(k_{},r_{})\stackrel{k_{}^2m^2}{=}`$ (41)
$`{\displaystyle \frac{1}{3}}{\displaystyle _0^1}𝑑X\left(\delta _{\lambda \lambda ^{}}(𝒱_x^u(X,t)+𝒱_x^d(X,t))+{\displaystyle \frac{1}{2ms}}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda )(𝒲_x^u(X,t)+𝒲_x^d(X,t))\right).`$
Let us recall the sum rules ,
$`{\displaystyle _0^1}𝑑X\left(_x^q(X,t)_x^{\overline{q}}(X,t)\right)`$ $`=`$ $`F_1^q(t)`$ (42)
$`{\displaystyle _0^1}𝑑X\left(𝒦_x^q(X,t)𝒦_x^{\overline{q}}(X,t)\right)`$ $`=`$ $`F_2^q(t)`$ (43)
where $`_x^q(X,t)`$ and $`𝒦_x^q(X,t)`$ are the total (valence $`+`$ sea) nonflip and spin-flip skewed quark distributions while $`_x^{\overline{q}}(X,t)`$ and $`𝒦_x^{\overline{q}}(X,t)`$ are the antiquark ones. Here $`F_1^q(t)`$ and $`F_2^q(t)`$ stand for the $`q`$-quark components of the Dirac and Pauli form factors of the proton). Since the contribution of sea quarks drops from the difference $`^q^{\overline{q}}`$ we can rewrite eqs. (43) as the sum rules for valence quark distributions
$$_0^1𝑑X𝒱_x^q(X,t)=F_1^q(t),_0^1𝑑X𝒲_x^q(X,t)=F_2^q(t).$$
(44)
Substituting this estimate to eq. (41) and using the isospin invariance, we get the final result for the nucleon impact factor at large transverse momenta
$$I_N(k_{},r_{})\stackrel{k_{}^2m^2}{=}\delta _{\lambda \lambda ^{}}F_1^{p+n}(t)+\frac{1}{2ms}\overline{u}(p^{},\lambda ^{})\overline{)}p_1\overline{)}r_{}u(p,\lambda )F_2^{p+n}(t),$$
(45)
where $`F_1^{p+n}(t)F_1^p(t)+F_1^n(t)`$ and $`F_2^{p+n}(t)F_2^p(t)+F_2^n(t)`$. As usual, $`F_1^{p(n)}`$ and $`F_2^{p(n)}`$ are the Dirac and Pauli form factors of the proton (neutron), respectively. With our accuracy they can be approximated by the dipole formulas
$$\begin{array}{ccccccc}F_1^p+\frac{t}{4m^2}F_2^p\hfill & =\hfill & G_E^p=\hfill & \frac{1}{\left(1+\frac{|t|}{0.7\mathrm{GeV}^2}\right)^2}\hfill & F_1^p+F_2^p=G_M^p\hfill & =\hfill & \frac{2.79}{\left(1+\frac{|t|}{0.71\mathrm{GeV}^2}\right)^2}\hfill \\ F_1^n+\frac{t}{4m^2}F_2^n\hfill & =\hfill & G_E^n=\hfill & 0\hfill & F_1^n+F_2^n=G_M^n\hfill & =\hfill & \frac{1.91}{\left(1+\frac{|t|}{0.71\mathrm{GeV}^2}\right)^2}\hfill \end{array},$$
(46)
which leads to <sup>*</sup><sup>*</sup>*Literally, one obtains
$$F_1^{p+n}(t)=\frac{1}{1+\left(\frac{|t|}{0.71GeV^2}\right)^2}\frac{1+0.88\frac{|t|}{4m^2}}{1+\frac{|t|}{4m^2}},F_2^{p+n}=\frac{0.12}{1+\left(\frac{|t|}{0.71GeV^2}\right)^2}$$
(47) but with our accuracy we can use the estimate (48).
$$F_1^{p+n}(t)=\frac{1}{1+\left(\frac{|t|}{0.7GeV^2}\right)^2},F_2^{p+n}=0$$
(48)
. Note that the spin-flip term turned out to be negligible for our values of $`t`$. Moreover, it vanishes at $`t=0`$ which suggests that it is numerically small at all $`t`$.
Our final estimate of the nucleon impact factor is
$$I_N(k_{},r_{})\stackrel{k_{}^2m^2}{=}\delta _{\lambda \lambda ^{}}F_1^{p+n}(t)$$
(49)
where $`F_1^{p+n}`$ is given by the dipole formula (48) The dipole formula for the neutron form factor does not seem to work as well as the dipole formula for the proton form factor. As a measure of the uncertainty we can compare the results obtained from eq. (48) to those obtained using the model from Ref. (which was fit only to the proton form factor) $`F_1^{p+n}(t)`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle _0^1}𝑑X\left(𝒱_x^u(X,t)+𝒱_x^d(X,t)\right)`$ (50) $`𝒱_x^u(X,t)`$ $`=`$ $`1.89X^{0.4}\overline{X}^{3.5}(1+6X)\mathrm{exp}\left({\displaystyle \frac{\overline{X}}{X}}{\displaystyle \frac{|t|}{2.8\mathrm{GeV}^2}}\right)`$ (51) $`𝒱_x^d(X,t)`$ $`=`$ $`0.54X^{0.6}\overline{X}^{4.2}(1+8X)\mathrm{exp}\left({\displaystyle \frac{\overline{X}}{X}}{\displaystyle \frac{|t|}{2.8\mathrm{GeV}^2}}\right)`$ (52) The results for the DVCS cross section in this model are about $`1.5`$ times bigger than the results obtained from the dipole formula (48). . In what follows we shall omit the factor $`\delta _{\lambda \lambda ^{}}`$ (as it was done in eq. (18)) since all our amplitudes will always be diagonal in the proton’s spin.
## IV The BFKL ladder
In the next order in perturbation theory the most important diagrams are those of the type shown in Fig. (5) Actually, this diagram gives the total contribution in LLA if one replaces the three-gluon vertex in Fig. (5) by the effective Lipatov’s vertex .
Calculation of this diagrams in the leading log approximation yields
$`V^{AB}=`$ (54)
$`{\displaystyle \frac{2sg^4}{\pi }}({\displaystyle e_q^2})\left(6\alpha _s\mathrm{ln}{\displaystyle \frac{1}{x}}\right){\displaystyle \frac{d^2k_{}}{4\pi ^2}\frac{d^2k_{}^{}}{4\pi ^2}\frac{I^{AB}(k_{},r_{})}{k_{}^2(r+k)_{}^2}K(k_{},k_{}^{},r_{})\frac{I_N(k_{}^{},r_{}^{})}{(k_{}^{})^2(r+k^{})_{}^2}}`$
where $`K(k_{},k_{}^{},r_{})`$ is the BFKL kernel
$`K(k_{},k_{}^{},r_{})=`$ (55)
$`r_{}^2+{\displaystyle \frac{k_{}^2(rk^{})_p^2}{(kk^{})_{}^2}}+{\displaystyle \frac{k_{}^2(rk^{})_p^2}{(kk^{})_{}^2}}+k_{}^2(kp)_{}^2{\displaystyle \frac{1}{2}}\delta (k_{}k_{}^{}){\displaystyle 𝑑p_{}\left(\frac{k_{}^2}{p_{}^2(kp)_{}^2}+\frac{(kr)_{}^2}{(pr)_{}^2(kp)_{}^2}\right)}`$ (56)
As we shall see below, the integral over $`k_{}^{}`$ converges at $`|k_{}^{}|m`$ so we can again use the approximation (18) for the nucleon impact factor. One obtains
$$d^2k_{}^{}K(k_{},k_{}^{},r_{})\frac{1}{(k_{}^{})^2(r+k^{})_{}^2}I_N(k_{}^{},r_{}^{})=\pi F_1^{p+n}(t)\left(\mathrm{ln}\frac{k_{}^2}{r_{}^2}+\mathrm{ln}\frac{(kr)_{}^2}{r_{}^2}\right)$$
(57)
and therefore the amplitude (54) takes the form
$$V^{AB}=\frac{g^4s}{\pi }F_1^{p+n}(t)\left(\frac{3\alpha _s}{\pi }\mathrm{ln}\frac{1}{x}\right)\frac{d^2k_{}}{4\pi ^2}\frac{I(k_{},r_{})}{k_{}^2(r+k)_{}^2}\left(\mathrm{ln}\frac{k_{}^2}{r_{}^2}+\mathrm{ln}\frac{(kr)_{}^2}{r_{}^2}\right).$$
(58)
Finally, the integration over $`k`$ yields
$`V^{AB}=`$ (61)
$`{\displaystyle \frac{2}{x}}\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)F_1^{p+n}(t)\left({\displaystyle \frac{3\alpha _s}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\right)`$
$`\left((e^A,e^B)_{}\left({\displaystyle \frac{1}{6}}\mathrm{ln}^3{\displaystyle \frac{Q^2}{|t|}}+2\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}2+\zeta (3)\right)+\left({\displaystyle \frac{2}{r_{}^2}}(e^A,r)_{}(e^B,r)_{}(e^A,e^B)_{}\right)\right)`$
where the accuracy is $`O(\frac{1}{\mathrm{ln}x})`$.
In the next order in BFKL approximation (see Fig. 6) it is still possible to obtain the DVCS amplitude (3) in the explicit form (we have not obtained the explicit expressions for higher-order terms in the BFKL expansion (61) <sup>§</sup><sup>§</sup>§It is possible to write down the result of the summation of the BFKL ladder in the form of Mellin integral over complex momenta using the Lipatov’s conformal eigenfunctions of the BFKL equation in the coordinate space. Unfortunately, we were not able to perform explicitly the integration of the Lipatov’s eigenfunctions with impact factors and without it the Mellin representation of the DVCS amplitude is useless for practical applications.).
The amplitude in this order is
$`V^{AB}`$ $`=`$ $`{\displaystyle \frac{g^4s}{\pi }}({\displaystyle e_q^2})\left(6\alpha _s\mathrm{ln}{\displaystyle \frac{1}{x}}\right)^2{\displaystyle \frac{d^2k_{}}{4\pi ^2}\frac{d^2k_{}^{}}{4\pi ^2}\frac{d^2k\mathrm{"}_{}}{4\pi ^2}I(k_{},r_{})}`$ (63)
$`{\displaystyle \frac{1}{k_{}^2(r+k)_{}^2}}K(k_{},k\mathrm{"}_{},r_{}){\displaystyle \frac{1}{(k\mathrm{"}_{})^2(r+k\mathrm{"})_{}^2}}K(k\mathrm{"}_{},k_{}^{},r_{}){\displaystyle \frac{1}{(k_{}^{})^2(r+k^{})_{}^2}}I_N(k_{}^{},r_{}^{}).`$
Once again, if we use the fact that the integral over $`k_{}^{}`$ converges at $`|k_{}^{}|m`$ we can approximate the nucleon impact factor by eq. (49), and obtain
$`{\displaystyle \frac{d^2k_{}^{}}{4\pi ^2}\frac{d^2k\mathrm{"}_{}}{4\pi ^2}K(k_{},k\mathrm{"}_{},r_{})\frac{1}{(k\mathrm{"})_{}^2(r+k\mathrm{"})_{}^2}K(k\mathrm{"}_{},k_{}^{},r_{})\frac{1}{(k^{})_{}^2(r+k^{})_{}^2}I_N(k_{}^{},r_{}^{})}=`$ (64)
$`{\displaystyle \frac{1}{4\pi }}F_1^{p+n}(t){\displaystyle \frac{d^2k\mathrm{"}_{}}{4\pi ^2}\frac{K(k_{},k\mathrm{"}_{},r_{})}{(k\mathrm{"})_{}^2(r+k\mathrm{"})_{}^2}\left(\mathrm{ln}\frac{(k\mathrm{"}_{})^2}{r_{}^2}+\mathrm{ln}\frac{(k\mathrm{"}r)_{}^2}{r_{}^2}\right)}=`$ (65)
$`{\displaystyle \frac{1}{16\pi ^2}}F_1^{p+n}(t)\left(\mathrm{ln}^2{\displaystyle \frac{k_{}^2}{r_{}^2}}+\mathrm{ln}^2{\displaystyle \frac{(kr)_{}^2}{r_{}^2}}\right).`$ (66)
The resulting integration over $`k_{}`$ yie lds
$`V^{AB}=`$ (69)
$`{\displaystyle \frac{9}{x}}\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^4({\displaystyle }e_q^2)F_1^{p+n}(t)\mathrm{ln}^2x[(e^A,e^B)_{}({\displaystyle \frac{1}{24}}\mathrm{ln}^4{\displaystyle \frac{Q^2}{|t|}}+\mathrm{ln}^2{\displaystyle \frac{Q^2}{|t|}}2\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}+`$
$`2(\zeta (3)1)+1.46)+({\displaystyle \frac{2}{r_{}^2}}(e^A,r)_{}(e^B,r)_{}(e^A,e^B)_{})].`$
As we mentioned, we were not able to obtain the explicit expressions for the amplitude in higher orders in perturbation theory. It turns out, however, that for HERA energies the achieved accuracy is reasonably good; the estimation of the next term gives $``$ 30% of the answer at not too low $`x`$ (see the discussion in next section). Our final result for the DVCS amplitude with transversely polarized photons is In the leading logarithmic approximation it is not possible to distinguish between $`\alpha _s(Q)`$ and $`\alpha _s(\sqrt{|t|})`$ – to this end one needs to use the NLO BFKL approximation (see also ) which is beyond the scope of this paper.
$`V^{AB}=`$ (71)
$`{\displaystyle \frac{2}{x}}\left({\displaystyle \frac{\alpha _s(Q)}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)F_1^{p+n}(t)\left[(e^A,e^B)_{}v+\left({\displaystyle \frac{2}{r_{}^2}}(e^A,r)_{}(e^B,r)_{}(e^A,e^B)_{}\right)v^{}\right],`$
where
$`v(x,Q^2/t)`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{|t|}}+2\right)+{\displaystyle \frac{3\alpha _s(Q)}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\left({\displaystyle \frac{1}{6}}\mathrm{ln}^3{\displaystyle \frac{Q^2}{|t|}}+2\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}2+\zeta (3)\right)`$ (72)
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{3\alpha _s(Q)}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\right)^2\left({\displaystyle \frac{1}{24}}\mathrm{ln}^4{\displaystyle \frac{Q^2}{|t|}}+\mathrm{ln}^2{\displaystyle \frac{Q^2}{|t|}}+2(\zeta (3)1)\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}+1.46\right)`$ (73)
$`v^{}(x,Q^2/t)`$ $`=`$ $`1+{\displaystyle \frac{3\alpha _s(Q)}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{3\alpha _s(Q)}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\right)^2.`$ (74)
Note that the spin-dependent part $`v^{}`$ does not contain any $`\mathrm{ln}\frac{Q^2}{|t|}`$ and is hence much smaller than the spin-independent part $`v`$. For the longitudinal polarization (16) the amplitude is twist-suppressed as $`\sqrt{\frac{|t|}{Q^2}}`$ so we have not calculated any terms beyond eq. (24). In the numerical analysis carried out in next sections we keep only the spin-independent part of the amplitude
$`V_{}{\displaystyle \frac{1}{4}}{\displaystyle e_{}^Ae_{}^BV^{AB}}={\displaystyle \frac{2}{x}}\left({\displaystyle \frac{\alpha _s(Q)}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)F_1^{p+n}(t)v(x,Q^2,t).`$ (75)
The above expressions give us the imaginary part of the DVCS amplitude. For the calculation of the DVCS cross section we need to know also the real part of this amplitude which can be estimated via the dispersion relation. For the positive-signature amplitude $`H_{}`$ ($`\frac{1}{4}e_{}^Ae_{}^BH^{AB}`$) we get (see also )
$$\mathrm{Re}H_{}(s)=\frac{\pi }{2}\mathrm{tan}\left(s\frac{d}{ds}\right)\mathrm{Im}H_{}(s),$$
(76)
which amounts to the substitution
$$\mathrm{ln}s\frac{1}{2}\left(\mathrm{ln}(siϵ)+\mathrm{ln}s\right)$$
(77)
in our amplitude (75). Thus, the real part is
$`R{\displaystyle \frac{1}{\pi }}\mathrm{Re}H_{}={\displaystyle \frac{2}{x}}\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)(F_1^p(t)+F_1^n(t))r(x,Q^2,t)`$ (80)
$`r(x,Q^2,t)={\displaystyle \frac{\pi }{2}}[{\displaystyle \frac{3\alpha _s}{\pi }}({\displaystyle \frac{1}{6}}\mathrm{ln}^3{\displaystyle \frac{Q^2}{|t|}}+2\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}2+\zeta (3))+`$
$`\left({\displaystyle \frac{3\alpha _s}{\pi }}\right)^2\mathrm{ln}{\displaystyle \frac{1}{x}}({\displaystyle \frac{1}{24}}\mathrm{ln}^4{\displaystyle \frac{Q^2}{|t|}}+\mathrm{ln}^2{\displaystyle \frac{Q^2}{|t|}}+2(\zeta (3)1)\mathrm{ln}{\displaystyle \frac{Q^2}{|t|}}+1.46)].`$
## V Comparison with the deep inelastic scattering
It is instructive to compare the DVCS amplitude $`V^{AB}`$ given by eq. (3) with the corresponding amplitude for the forward $`\gamma ^{}`$ scattering
$$T^{AB}=ie_\nu ^Ae_\mu ^B𝑑ze^{iqz}p|T\{j^\mu (z)j^\nu (0)\}|p.$$
(81)
The imaginary part of this amplitude is the total cross section for deep inelastic scattering (DIS)
$`{\displaystyle \frac{1}{\pi }}\mathrm{Im}T^{AB}=W^{AB}=`$ (82)
$`e_\nu ^Ae_\mu ^B\left[\left({\displaystyle \frac{q_\mu q_\nu }{q^2}}g_{\mu \nu }\right)F_1(x,Q^2)+{\displaystyle \frac{1}{pq}}\left(p_\mu q_\mu {\displaystyle \frac{pq}{q^2}}\right)\left(p_\nu q_\nu {\displaystyle \frac{pq}{q^2}}\right)F_2(x,Q^2)\right]`$ (83)
For example $`W^{AB}`$ averaged over the transverse polarizations of the photons is
$$W_{}\stackrel{\mathrm{def}}{}\frac{1}{4}e_{}^Ae_{}^BW^{AB}=F_1(x,Q^2)=\frac{1}{2x}F_2(x,Q^2)$$
(84)
(at the leading twist level we have the Callan-Gross relation $`F_2=2xF_1`$). We will compare the imaginary part of the DVCS amplitude $`V_{}`$ given by eq. (75) to the result for $`W_{}`$ calculated with the same accuracy. (We use the notation $`W_{}(x)`$ rather than $`F_1(x)`$ in order to avoid confusion with $`F_1(t)`$).
Similarly to the DVCS case, the DIS amplitude has the form (cf. eqs.(19,)(54), and (63)):
$`W_{}`$ $`=`$ $`{\displaystyle \frac{2g^2s}{\pi }}\left({\displaystyle e_q^2}\right){\displaystyle \frac{d^2k_{}}{4\pi ^2}\frac{1}{k_{}^4}I_{}(k_{},0)}`$ (87)
$`[1+{\displaystyle \frac{3g^2}{8\pi ^3}}\mathrm{ln}{\displaystyle \frac{1}{x}}{\displaystyle }d^2k_{}^{}K(k_{},k_{}^{},0)+`$
$`{\displaystyle \frac{9g^4}{128\pi ^6}}\mathrm{ln}^2{\displaystyle \frac{1}{x}}{\displaystyle }d^2k_{}^{}{\displaystyle }d^2k\mathrm{"}_{}K(k_{},k\mathrm{"}_{},0){\displaystyle \frac{1}{(k_{}^{\prime \prime })^2}}K(k\mathrm{"}_{},k_{}^{},0)]{\displaystyle \frac{1}{(k_{}^{})^2}}I_N(k_{}^{},0)`$
where $`I_{}(k_{},0)`$ is the virtual photon impact factor averaged over the transverse polarizations
$$I_{}(k_{},0)=\frac{1}{2}_0^1\frac{d\alpha }{2\pi }_0^1\frac{d\alpha ^{}}{2\pi }\frac{k_{}^2(12\alpha \overline{\alpha })(12\alpha ^{}\overline{\alpha ^{}})}{k_{}^2\alpha ^{}\overline{\alpha ^{}}+Q^2\alpha ^{}\alpha \overline{\alpha }}$$
(88)
The nucleon impact factor $`I_N(k_{}^{},0)`$ cannot be calculated in perturbation theory since it is determined by the large-scale nucleon dynamics. However, we know the asymptotics at large $`k_{}m`$
$$I_N(k_{},0)\stackrel{k_{}^2m^2}{=}F_1^{p+n}(0)=1$$
(89)
Also, at $`I_N(k_{},0)0`$ at $`k0`$ due to the gauge invariance. It seems reasonable to model this impact factor by the simple formula
$$I_N(k_{},0)=\frac{k_{}^2}{k_{}^2+m^2}$$
(90)
which has the correct behavior both at large and small $`k_{}`$. With this model, the DIS amplitude (87) takes the form
$`W_{}={\displaystyle \frac{F_2}{2x}}=`$ (94)
$`{\displaystyle \frac{4}{3x}}\left({\displaystyle \frac{\alpha _s(Q)}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)[({\displaystyle \frac{1}{2}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{7}{6}}\mathrm{ln}{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{77}{18}})+`$
$`{\displaystyle \frac{3\alpha _s}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\left({\displaystyle \frac{1}{6}}\mathrm{ln}^3{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{7}{12}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{77}{18}}\mathrm{ln}{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{131}{27}}+2\zeta (3)\right)+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{3\alpha _s}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\right)^2`$
$`({\displaystyle \frac{1}{24}}\mathrm{ln}^4{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{7}{36}}\mathrm{ln}^3{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{77}{36}}\mathrm{ln}^2{\displaystyle \frac{Q^2}{m^2}}+({\displaystyle \frac{131}{27}}+4\zeta (3))\mathrm{ln}{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{1396}{81}}{\displaystyle \frac{\pi ^4}{15}}+{\displaystyle \frac{14}{3}}\zeta (3))].`$
Note that the coefficients in front of leading logs of $`Q^2`$, determined by the anomalous dimensions of twist-2 operators, coincide up to the factor $`2/3`$. The graph of the model (94) versus the experimental data is presented in Fig. 7 for $`Q^2=10`$GeV<sup>2</sup> and $`Q^2=35`$GeV<sup>2</sup> (we take $`e_q^2=\frac{10}{9}`$).
In the case of DIS it is possible to calculate explicitly the next term in BFKL series (94) For DIS it is possible to write down the total BFKL sum as a Mellin integral and unlike DVCS the integrals of impact factors with the BFKL eigenfunctions $`(k_{}^2)^{\frac{1}{2}+i\nu }`$ can be calculated explicitly. Eqs. (94) and (96) correspond to the expansion of this explicit expression in powers of $`\alpha _s\mathrm{ln}x`$. . It has the form
$`{\displaystyle \frac{4}{3x}}\left({\displaystyle \frac{\alpha _s(Q)}{\pi }}\right)^2({\displaystyle \underset{\mathrm{flavors}}{}}e_q^2)[{\displaystyle \frac{1}{6}}\left({\displaystyle \frac{3\alpha _s}{\pi }}\mathrm{ln}{\displaystyle \frac{1}{x}}\right)^3({\displaystyle \frac{1}{120}}\mathrm{ln}^5{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{7}{144}}\mathrm{ln}^4{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{77}{108}}\mathrm{ln}^3{\displaystyle \frac{Q^2}{m^2}}+({\displaystyle \frac{131}{54}}+`$ (95)
$`+3\zeta (3))\mathrm{ln}^2{\displaystyle \frac{Q^2}{m^2}}+({\displaystyle \frac{1396}{81}}{\displaystyle \frac{\pi ^4}{15}}+7\zeta (3))\mathrm{ln}{\displaystyle \frac{Q^2}{m^2}}+{\displaystyle \frac{4736}{243}}{\displaystyle \frac{7\pi ^4}{90}}+{\displaystyle \frac{77}{3}}\zeta (3)+6\zeta (5))]`$ (96)
The ratio of this $`\left(\alpha _s\mathrm{ln}x\right)^3`$ term to the sum of the first three ones (94) is presented in Fig. 8 for $`Q^2=10`$GeV<sup>2</sup> and $`Q^2=35`$GeV<sup>2</sup>.
From these graphs we see that the sum of the first tree terms gives the reliable estimate of the DIS amplitude at not too low $`x`$ and it is expected that the same will also be true for DVCS amplitude <sup>\**</sup><sup>\**</sup>\** At very small $`x10^3÷10^5`$ the full BFKL result for $`F_2`$ in our model is growing more rapidly than Fig. 7. On the other hand if one takes into account the NLO BFKL corrections the result for $`F_2`$ at very small x goes well under the experimental points. This indicates that at such $`x`$ we need to unitarize the BFKL pomeron, which is currently an unsolved problem. (The best hope is to find the effective action for the BFKL pomeron (see e.g. ,)). On the contrary, at “intermediate” $`x0.1÷0.001`$, we see from Fig. 7 that, since the corrections almost cancel each other, it makes sense to take into account only a few first terms in BFKL series. .
It is instructive to compare the t-dependence of our DVCS amplitude (73) with the model used in the paper
$`V_1(x,t,Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{R}}F_1(x,Q^2)e^{bt/2}`$ (97)
$`V_2(x,t,Q^2)`$ $`=`$ $`{\displaystyle \frac{1}{R}}F_1(x,Q^2){\displaystyle \frac{1}{\left(1+\frac{|t|}{0.71}\right)^2}}`$ (98)
where $`R0.5`$ for our energies. (Literally, the model used in ref. corresponds to $`V_1`$ but it is more natural to approximate the $`t`$ \- dependence by the dipole formula). The comparison is shown in Fig. 9 for $`Q^2=10`$GeV<sup>2</sup>, $`Q^2=35`$GeV<sup>2</sup> and $`x`$=0.01, $`x`$=0.001.
## VI DVCS cross section
In order to estimate the cross section for DVCS at HERA kinematics ($`Q^2>6`$GeV<sup>2</sup> and $`x<10^2`$) we will use formulas from Ref. (see also Ref. ) with the trivial substitution $`\frac{1}{2x}F_2(x)R^1e^{bt/2}V_{}(x,Q^2,t)`$. The expressions for the DVCS cross section, the QED Compton (Bethe-Heitler) cross section, and the interference term have the form ($`\overline{y}1y`$) <sup>††</sup><sup>††</sup>††The expression for the interference term from ref. is corrected by factor 2 , :
$`{\displaystyle \frac{d\sigma ^{\mathrm{DVCS}}}{dxdydtd\varphi _r}}`$ $`=`$ $`\pi \alpha ^3x{\displaystyle \frac{1+\overline{y}^2}{Q^4y}}(V_{}^2(x,Q^2,t)+R_{}^2(x,Q^2,t))`$ (99)
$`{\displaystyle \frac{d\sigma ^{\mathrm{QEDC}}}{dxdydtd\varphi _r}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^3}{\pi x}}{\displaystyle \frac{y(1+\overline{y}^2)}{|t|Q^2\overline{y}}}\left((F_1^p(t))^2+{\displaystyle \frac{|t|}{4m^2}}(F_p^2(t))^2\right)`$ (100)
$`{\displaystyle \frac{d\sigma ^{\mathrm{INT}}}{dxdydtd\varphi _r}}`$ $`=`$ $`2\alpha ^3{\displaystyle \frac{(1+\overline{y}^2)}{Q^3\sqrt{\overline{y}|t|}}}R_{}(x,Q^2,t)F_1^p(t)\mathrm{cos}\varphi _r.`$ (101)
Here $`y=1\frac{E^{}}{E}`$ ($`E`$ and $`E^{}`$ are the incident and scattered electron energies, respectively, as defined in the proton rest frame) and $`\varphi _r=\varphi _e+\varphi _N`$ where $`\varphi _N`$ is the azimuthal angle between the plane defined by $`\gamma ^{}`$ and the final state proton and the $`xz`$ plane and $`\varphi _e`$ is the azimuthal angle between the plane defined by the initial and final state electron and $`xz`$ plane (see Ref. ). As mentioned above, we approximate the Dirac and Pauli form factors of the proton by the dipole formulas (46).
At first let us discuss the relative weight of the above cross sections. We start with the asymmetry defined in ref.
$$A=\frac{_{\pi /2}^{\pi /2}𝑑\varphi _r𝑑\sigma ^{\mathrm{DQI}}_{\pi /2}^{3\pi /2}𝑑\varphi _r𝑑\sigma ^{\mathrm{DQI}}}{_0^{2\pi }𝑑\varphi _r𝑑\sigma ^{\mathrm{DQI}}}$$
(102)
where
$$d\sigma ^{\mathrm{DQI}}d\sigma ^{\mathrm{DVCS}}+d\sigma ^{\mathrm{QEDC}}+d\sigma ^{\mathrm{INT}}.$$
(103)
The asymmetry shows the relative importance of the interference term, which is proportional to the real part of the DVCS amplitude. In our approximation the asymmetry is
$$A(y,t)=\frac{4y\sqrt{\frac{Q^2}{|t|\overline{y}}}(e_q^2)\left(\frac{\alpha _s}{\pi }\right)^2\left(1+2.8\frac{|t|}{4m^2}\right)r}{4\pi ^2(e_q^2)^2(v^2+r^2)\left(\frac{\alpha _s}{\pi }\right)^4\left(1+\frac{|t|}{4m^2}\right)+\frac{y^2Q^2}{\overline{y}|t|}\left(1+7.84\frac{|t|}{4m^2}\right)}$$
(104)
The plots of asymmetry versus $`y`$ and $`|t|`$ are given by Fig. 10.
Second, we define the ratio of the DVCS and Bethe-Heitler cross sections
$$D(y,t)\frac{d\sigma _{DVCS}}{d\sigma _{QEDC}}=\frac{4\pi ^2(e_q^2)^2(v^2+r^2)\left(\frac{\alpha _s}{\pi }\right)^4\left(1+\frac{|t|}{4m^2}\right)\overline{y}\frac{|t|}{Q^2}}{y^2\left(1+7.84\frac{|t|}{4m^2}\right)}$$
(105)
This ratio is presented on Fig. 11.
We see that there is a sharp dependence on $`y`$; at $`y>0.2`$ the DVCS part is negligible in comparison to Bethe-Heitler background whereas at $`y<0.05`$ the QEDC background is small in comparison to DVCS.
Finally let us estimate the relative weight of the DVCS signal (starting from $`|t|=1`$ GeV<sup>2</sup>) as compared to the DIS background. We define (cf. ref. )
$`R_\gamma ={\displaystyle \frac{\sigma (\gamma ^{}+p\gamma +p)}{\sigma (\gamma ^{}+p\gamma ^{}+p)}}`$ (106)
$`{\displaystyle \frac{4\pi \alpha }{Q^2F_2(x,Q^2)}}\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^4\left({\displaystyle e_q^2}\right)^2{\displaystyle _1^{Q^2}}𝑑t\left(F_1^{p+n}(t)\right)^2(v^2(x,Q^2/t)+r^2(x,Q^2/t))`$ (107)
At $`Q^2=10`$GeV<sup>2</sup> we find $`R_\gamma =1.56\times 10^5`$ for $`x=0.01`$ and $`R_\gamma =2.36\times 10^5`$ for $`x=0.001`$, while for $`Q^2=35`$GeV<sup>2</sup> we find $`R_\gamma =0.62\times 10^5`$ for $`x=0.01`$ and $`R_\gamma =0.71\times 10^5`$ for $`x=0.001`$.
The expressions (60)-(62) are correct if $`Q^2|t|`$ up to $`O(\frac{|t|}{Q^2})`$ accuracy with the notable exception of the correction $`O(\frac{\sqrt{|t|}}{Q})`$ coming from the expansion of electron propagator in the u-channel of the Bethe-Heitler amplitude. As suggested in ref. , at intermediate $`t`$ one can keep the propagator in unexpanded form (and expand the rest of the amplitude, as we have done above). This amounts to the replacement
$$\overline{y}\overline{y}\left[(1+\frac{|t|}{Q^2\overline{y}})(1+\frac{|t|\overline{y}}{Q^2})2\frac{(2y)}{\sqrt{\overline{y}}}\sqrt{\frac{|t|}{Q^2}}\mathrm{cos}\varphi _r+4\frac{|t|}{Q^2}\mathrm{cos}^2\varphi _r\right]$$
(108)
in the numerator in eqs. (61) and (62) (see ref. ).
The resulting asymmetry (63) is presented in Fig. 12. We see that the correction factor (68) crucially changes the behavior of the asymmetry due to the fact that it restores the azimuthal dependence of the QEDC amplitude which was not taken into account in eqs. (60-62). In order to find asymmerty at these $`Q^2`$ and $`t`$ with greater accuracy one should take into account other twist-4 contributions as well. On the contrary, the ratio $`D(x,Q^2/t)`$ does not change much (see Fig. 13) so we hope that our leading-twist results for the ratio presented in Fig. 11 are reliable.
## VII Conclusion
The DVCS in the kinematical region (1) is probably the best place to test the momentum transfer dependence of the BFKL pomeron. Without this dependence, the model (98) would be exact, hence the upper curves in Fig. 9 indicate how important is the $`t`$-dynamics of the pomeron. We see that the $`t`$-dependence of the BFKL pomeron changes the cross section at $`t>2`$GeV<sup>2</sup> by orders of magnitude and therefore it should be be possible to detect it.
The pQCD calculation of the DVCS amplitude in the region (1) is in a sense even more reliable than the calculation of usual DIS amplitudes since it does not rely on the models for nucleon parton distributions . Indeed, all the non-perturbative nucleon input is contained in the Dirac form factor of the nucleon <sup>‡‡</sup><sup>‡‡</sup>‡‡ There are, of course, the non-perturbative corrections to the BFKL pomeron itself. At present, it is not clear how to take them into account., which is known to a pretty good accuracy. (Of course any reasonable models of nucleon parton distributions such as (39) should reproduce the form factor after integration over $`X`$).
Finally, let us discuss uncertainties in our approximation and possible ways to improve it. One obvious improvement would be to calculate (at least numerically) the next $`\left(\alpha _s\mathrm{ln}x\right)^3`$ term in the BFKL series for the DVCS amplitude. Hopefully, it will be as small as the corresponding calculation of the DIS amplitude suggests. Second, there are non-perturbative corrections to the BFKL pomeron which we mention above. These non-perturbative corrections correspond to the situation like the “aligned jet model” when one of the two gluons in Fig. (1) is soft and all the momentum transfers through the other gluon. It is not clear how to take these corrections into account, but one should expect them to be smaller than the corresponding corrections to $`F_2(x)`$ which come from two non-perturbative gluons in Fig. 1 (in other words, from the “soft pomeron” contribution to $`F_2(x)`$).
The biggest uncertainty in our calculation is the argument of coupling constant $`\alpha _s`$ which we take to be $`Q^2`$. As we mention above, it is not possible to fix the argument of $`\alpha _s`$ in the LLA, so we could have used $`\alpha _s(|t|)`$ instead. We hope to overcome this difficulty by using the results of NLO BFKL in our future work.
## Acknowledgments
The authors are grateful to A.V. Belitsky, C.E. Hyde-Wright, I.V. Musatov, A.V. Radyushkin, and M.I. Strikman for valuable discussions. This work was supported by DOE contract DE-AC05-84ER40150 under which the Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility.
## References |
warning/0002/hep-ph0002213.html | ar5iv | text | # 1 Introduction
## 1 Introduction
One of the most striking phenomenological implications of Supersymmetry (SUSY) is the prediction of a relatively light Higgs boson, which is common to all supersymmetric models whose couplings remain in the perturbative regime up to a very high energy scale . The search for the lightest Higgs boson thus allows a crucial test of SUSY and especially the Minimal Supersymmetric Standard Model (MSSM) . Therefore it is one of the main goals at the present and the next generation of colliders. A precise knowledge of the dependence of the masses and mixing angles of the Higgs sector of the MSSM on the relevant SUSY parameters is necessary for a detailed analysis of SUSY phenomenology at LEP2, the upgraded Tevatron, and also for the LHC and a future linear $`e^+e^{}`$ collider, where high-precision measurements in the Higgs sector might become possible.
The mass of the lightest $`𝒞𝒫`$-even MSSM Higgs boson, $`m_h`$, is bounded from above by about $`m_h\stackrel{<}{}\mathrm{\hspace{0.33em}135}\mathrm{GeV}`$, including radiative corrections up to the two-loop level . The mixing angle $`\alpha _{\mathrm{eff}}`$ that diagonalizes the neutral $`𝒞𝒫`$-even Higgs sector receives the same kind of corrections. By incorporating $`\alpha _{\mathrm{eff}}`$ into the Higgs decay widths, the leading electroweak corrections to the decay of the neutral $`𝒞𝒫`$-even Higgs bosons are taken into account.
In this note we present the Fortran code FeynHiggsFast. Using low energy MSSM parameters as input, it evaluates the masses of the neutral $`𝒞𝒫`$-even Higgs bosons, $`m_h`$ and $`m_H`$, as well as the corresponding mixing angle, $`\alpha _{\mathrm{eff}}`$, at the two-loop level. In addition the mass of the charged Higgs boson, $`m_{H^\pm }`$, is evaluated at the one-loop level. The $`\rho `$-parameter, leading to constraints in the scalar fermion sector of the MSSM, is evaluated up to $`𝒪(\alpha \alpha _s)`$, taking into account the gluon exchange contribution at the two-loop level . FeynHiggsFast is based on a compact analytical approximation formula, containing at the two-loop level the dominant corrections in $`𝒪(\alpha \alpha _s)`$ obtained in the Feynman-diagrammatic (FD) approach and subdominant corrections of $`𝒪(G_F^2m_t^6)`$ obtained with renormalization group (RG) methods . In comparison with the FD result, consisting of the complete one-loop and the two-loop contributions as given in Ref. (incorporated into the Fortran code FeynHiggs ), the approximation formula is much shorter. Thus FeynHiggsFast is about $`3\times 10^4`$ times faster than FeynHiggs. Agreement between these two codes of better than $`2\mathrm{GeV}`$ is found for most parts of the MSSM parameter space.
The following sections are organized as follows: In Sect. 2 we shortly summarize the analytical approximation formula and perform a comparison with the full FD result. A description of how to use FeynHiggsFast is given in Sect. 3. The conclusions can be found in Sect. 4.
## 2 The analytical approximation
In order to fix our notations, we first list the conventions for the MSSM scalar top sector: the mass matrix in the basis of the current eigenstates $`\stackrel{~}{t}_L`$ and $`\stackrel{~}{t}_R`$ is given by (neglecting contributions $`M_Z^2`$):
$$_{\stackrel{~}{t}}^2=\left(\begin{array}{cc}M_{\stackrel{~}{t}_L}^2+m_t^2& m_tX_t\\ m_tX_t& M_{\stackrel{~}{t}_R}^2+m_t^2\end{array}\right),$$
(1)
where
$$m_tX_t=m_t(A_t\mu \mathrm{cot}\beta ).$$
(2)
Furtheron the following simplification is used:
$$M_{\stackrel{~}{t}_L}=M_{\stackrel{~}{t}_R}:=m_{\stackrel{~}{q}}.$$
(3)
In this simplified case we define<sup>1</sup><sup>1</sup>1 The more general case with $`M_{\stackrel{~}{t}_L}M_{\stackrel{~}{t}_R}`$ can be found in eq. (25).
$$M_S^2:=m_{\stackrel{~}{q}}^2+m_t^2.$$
(4)
For further details see Ref. .
At the tree level the mass matrix of the neutral $`𝒞𝒫`$-even Higgs bosons in the $`\varphi _1,\varphi _2`$ basis can be expressed as follows:
$`M_{\mathrm{Higgs}}^{2,\mathrm{tree}}`$ $`=`$ $`\left(\begin{array}{cc}m_{\varphi _1}^2& m_{\varphi _1\varphi _2}^2\\ m_{\varphi _1\varphi _2}^2& m_{\varphi _2}^2\end{array}\right)`$ (7)
$`=`$ $`\left(\begin{array}{cc}M_A^2\mathrm{sin}^2\beta +M_Z^2\mathrm{cos}^2\beta & (M_A^2+M_Z^2)\mathrm{sin}\beta \mathrm{cos}\beta \\ (M_A^2+M_Z^2)\mathrm{sin}\beta \mathrm{cos}\beta & M_A^2\mathrm{cos}^2\beta +M_Z^2\mathrm{sin}^2\beta \end{array}\right).`$ (10)
$`M_A`$ denotes the mass of the $`𝒞𝒫`$-odd Higgs boson, $`M_Z`$ is the mass of the $`Z`$ boson, and $`\mathrm{tan}\beta =v_2/v_1`$ is the ratio of the two vacuum expectation values of the two Higgs doublets in the MSSM, see Ref. .
At higher orders, the Higgs mass matrix (10) is supplemented by the renormalized self-energies $`\widehat{\mathrm{\Sigma }}_s(q^2),s=\varphi _1,\varphi _2,\varphi _1\varphi _2`$. Here we use the approximation of neglected external momentum: $`\widehat{\mathrm{\Sigma }}_s\widehat{\mathrm{\Sigma }}_s(0)`$. The masses $`m_h`$ and $`m_H`$ are then obtained by diagonalizing the higher order corrected mass matrix with the mixing angle $`\alpha _{\mathrm{eff}}`$:
$`\left(\begin{array}{cc}m_{\varphi _1}^2\widehat{\mathrm{\Sigma }}_{\varphi _1}& m_{\varphi _1\varphi _2}^2\widehat{\mathrm{\Sigma }}_{\varphi _1\varphi _2}\\ m_{\varphi _1\varphi _2}^2\widehat{\mathrm{\Sigma }}_{\varphi _1\varphi _2}& m_{\varphi _2}^2\widehat{\mathrm{\Sigma }}_{\varphi _2}\end{array}\right)\stackrel{\alpha _{\mathrm{eff}}}{}\left(\begin{array}{cc}m_H^2& 0\\ 0& m_h^2\end{array}\right),`$ (15)
$`\alpha _{\mathrm{eff}}=\mathrm{arctan}\left[{\displaystyle \frac{(M_A^2+M_Z^2)\mathrm{sin}\beta \mathrm{cos}\beta \widehat{\mathrm{\Sigma }}_{\varphi _1\varphi _2}}{M_Z^2\mathrm{cos}^2\beta +M_A^2\mathrm{sin}^2\beta \widehat{\mathrm{\Sigma }}_{\varphi _1}m_h^2}}\right],{\displaystyle \frac{\pi }{2}}<\alpha _{\mathrm{eff}}<{\displaystyle \frac{\pi }{2}}.`$ (16)
Here we only give a very brief outline of the calculation of the renormalized Higgs boson self-energies. A detailed description can be found in Ref. .
The mass of the lightest Higgs boson receives contributions from all sectors of the MSSM, but not all are numerically of equal relevance. The dominant corrections arise from the $`t\stackrel{~}{t}`$ sector of the MSSM: The results for the renormalized self-energies $`\widehat{\mathrm{\Sigma }}_s`$ have been derived analytically in the FD approach in Refs. . These results, however, are rather lengthy. In order to derive a compact analytical expression, several approximations have been made as described in Ref. . The main step of the approximation consists of a Taylor expansion in
$$\mathrm{\Delta }_{\stackrel{~}{t}}=\frac{|m_tX_t|}{M_S^2}$$
(17)
of the $`\widehat{\mathrm{\Sigma }}_s(0)`$. For the one-loop correction the expansion has been performed up to $`𝒪(\mathrm{\Delta }_{\stackrel{~}{t}}^8)`$; all three renormalized one-loop Higgs-boson self-energies give a contribution. Concerning the two-loop self-energies, the expansion has been carried out up to $`𝒪(\mathrm{\Delta }_{\stackrel{~}{t}}^4)`$; in the approximation considered here only $`\widehat{\mathrm{\Sigma }}_{\varphi _2}^{(2)}`$ gives a non-zero contribution.
Leading contributions beyond $`𝒪(\alpha \alpha _s)`$ have been taken into account by incorporating the leading two-loop Yukawa correction of $`𝒪(G_F^2m_t^6)`$ . Furthermore the result has been expressed in terms of the $`\overline{\mathrm{MS}}`$ top-quark mass
$$\overline{m}_t=\overline{m}_t(m_t)\frac{m_t}{1+\frac{4}{3\pi }\alpha _s(m_t)}$$
(18)
instead of the pole mass $`m_t`$. This leads to an additional contribution in $`𝒪(\alpha \alpha _s^2)`$.
The analytical expressions arising from the $`t\stackrel{~}{t}`$ sector are given at the one-loop level as follows:
$`\widehat{\mathrm{\Sigma }}_{\varphi _1}^{(1)}(0)`$ $`=`$ $`{\displaystyle \frac{G_F\sqrt{2}}{\pi ^2}}M_Z^4\mathrm{\Lambda }\mathrm{cos}^2\beta \mathrm{log}\left({\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right),`$ (19)
$`\widehat{\mathrm{\Sigma }}_{\varphi _1\varphi _2}^{(1)}(0)`$ $`=`$ $`{\displaystyle \frac{G_F\sqrt{2}}{\pi ^2}}M_Z^2\mathrm{cot}\beta \left[{\displaystyle \frac{3}{8}}\overline{m}_t^2+M_Z^2\mathrm{\Lambda }\mathrm{sin}^2\beta \right]\mathrm{log}\left({\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right),`$ (20)
$`\widehat{\mathrm{\Sigma }}_{\varphi _2}^{(1)}(0)`$ $`=`$ $`{\displaystyle \frac{G_F\sqrt{2}}{\pi ^2}}{\displaystyle \frac{\overline{m}_t^4}{8\mathrm{sin}^2\beta }}\{2{\displaystyle \frac{M_Z^2}{\overline{m}_t^2}}+{\displaystyle \frac{11}{10}}{\displaystyle \frac{M_Z^4}{\overline{m}_t^4}}`$ (21)
$`+\left[126{\displaystyle \frac{M_Z^2}{\overline{m}_t^2}}\mathrm{sin}^2\beta +8{\displaystyle \frac{M_Z^4}{\overline{m}_t^4}}\mathrm{\Lambda }\mathrm{sin}^4\beta \right]\mathrm{log}\left({\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right)`$
$`+{\displaystyle \frac{X_t^2}{M_S^2}}\left[12+4{\displaystyle \frac{M_Z^2}{\overline{m}_t^2}}+6{\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right]+{\displaystyle \frac{X_t^4}{M_S^4}}\left[14{\displaystyle \frac{\overline{m}_t^2}{M_S^2}}+3{\displaystyle \frac{\overline{m}_t^4}{M_S^4}}\right]`$
$`+{\displaystyle \frac{X_t^6}{M_S^6}}\left[{\displaystyle \frac{3}{5}}{\displaystyle \frac{\overline{m}_t^2}{M_S^2}}{\displaystyle \frac{12}{5}}{\displaystyle \frac{\overline{m}_t^4}{M_S^4}}+2{\displaystyle \frac{\overline{m}_t^6}{M_S^6}}\right]`$
$`+{\displaystyle \frac{X_t^8}{M_S^8}}[{\displaystyle \frac{3}{7}}{\displaystyle \frac{\overline{m}_t^4}{M_S^4}}{\displaystyle \frac{12}{7}}{\displaystyle \frac{\overline{m}_t^6}{M_S^6}}+{\displaystyle \frac{3}{2}}{\displaystyle \frac{\overline{m}_t^8}{M_S^8}}]\},`$
with
$$\mathrm{\Lambda }=\left(\frac{1}{8}\frac{1}{3}s_W^2+\frac{4}{9}s_W^4\right),s_W^2=1\frac{M_W^2}{M_Z^2}.$$
(22)
The two-loop contributions in $`𝒪(\alpha \alpha _s)`$ read:
$`\widehat{\mathrm{\Sigma }}_{\varphi _1}^{(2)}(0)`$ $`=`$ $`0,`$
$`\widehat{\mathrm{\Sigma }}_{\varphi _1\varphi _2}^{(2)}(0)`$ $`=`$ $`0,`$
$`\widehat{\mathrm{\Sigma }}_{\varphi _2}^{(2)}(0)`$ $`=`$ $`{\displaystyle \frac{G_F\sqrt{2}}{\pi ^2}}{\displaystyle \frac{\alpha _s}{\pi }}{\displaystyle \frac{\overline{m}_t^4}{\mathrm{sin}^2\beta }}[4+3\mathrm{log}^2\left({\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right)+2\mathrm{log}\left({\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right)6{\displaystyle \frac{X_t}{M_S}}`$ (23)
$`{\displaystyle \frac{X_t^2}{M_S^2}}\{3\mathrm{log}\left({\displaystyle \frac{\overline{m}_t^2}{M_S^2}}\right)+8\}+{\displaystyle \frac{17}{12}}{\displaystyle \frac{X_t^4}{M_S^4}}].`$
The two-loop Yukawa correction in this approximation reads:
$$\widehat{\mathrm{\Sigma }}_{\varphi _2}^{(2)}(0)=\frac{9}{16\pi ^4}\frac{G_F^2\overline{m}_t^6}{\mathrm{sin}^2\beta }\left[\mathrm{log}^2\left(\frac{\overline{m}_t^2}{M_S^2}\right)2\frac{X_t^2}{M_S^2}\mathrm{log}\left(\frac{\overline{m}_t^2}{M_S^2}\right)+\frac{1}{6}\frac{X_t^4}{M_S^2}\mathrm{log}\left(\frac{\overline{m}_t^2}{M_S^2}\right)\right].$$
(24)
$`M_S`$ has to be chosen according to
$$M_S=\{\begin{array}{cc}\sqrt{m_{\stackrel{~}{q}}^2+\overline{m}_t^2}:\hfill & M_{\stackrel{~}{t}_L}=M_{\stackrel{~}{t}_R}=m_{\stackrel{~}{q}}\hfill \\ \left[M_{\stackrel{~}{t}_L}^2M_{\stackrel{~}{t}_R}^2+\overline{m}_t^2(M_{\stackrel{~}{t}_L}^2+M_{\stackrel{~}{t}_R}^2)+\overline{m}_t^4\right]^{\frac{1}{4}}:\hfill & M_{\stackrel{~}{t}_L}M_{\stackrel{~}{t}_R}\hfill \end{array}$$
(25)
For the one-loop corrections from the other sectors of the MSSM the logarithmic approximation given in Ref. has been used for FeynHiggsFast, see also Ref. . In order to obtain the radiatively corrected Higgs boson masses and the mixing angle, the renormalized Higgs boson self-energies have to be inserted into eq. (15), and the corresponding diagonalization has to be performed.
We have also implemented the result of the MSSM contributions to $`\mathrm{\Delta }\rho `$ . Here the corrections arising from $`\stackrel{~}{t}/\stackrel{~}{b}`$-loops up to $`𝒪(\alpha \alpha _s)`$ have been taken into account. The result is valid for arbitrary parameters in the $`\stackrel{~}{t}`$\- and $`\stackrel{~}{b}`$-sector, also taking into account the mixing in the $`\stackrel{~}{b}`$-sector which can have a non-negligible effect in the large $`\mathrm{tan}\beta `$ scenario .
The two-loop result can be separated into the pure gluon-exchange contribution, which can be expressed by a very compact formula, allowing a very fast evaluation, and the pure gluino-exchange contribution, which is given by a rather lengthy expression. The latter correction goes to zero with increasing gluino mass and can thus be discarded for a heavy gluino. Concerning the implementation of the $`\rho `$-parameter into FeynHiggsFast, we have neglected the gluino exchange contribution. The $`\rho `$-parameter can be used as an additional constraint (besides the experimental bounds) on the squark masses. A value of $`\mathrm{\Delta }\rho `$ outside the experimentally preferred region of $`\mathrm{\Delta }\rho ^{\mathrm{SUSY}}10^3`$ indicates experimentally disfavored $`\stackrel{~}{t}`$\- and $`\stackrel{~}{b}`$-masses.
For illustration of the quality of the compact approximation, we compare $`m_h`$ and $`\alpha _{\mathrm{eff}}`$ with the results from the complete Feynman-diagrammatic calculation. The FD calculation contains the full diagrammatic one-loop contribution , the complete leading two-loop corrections in $`𝒪(\alpha \alpha _s)`$ , and the two contributions beyond $`𝒪(\alpha \alpha _s)`$ (see eq. (18) and eq. (24)) without approximation.
For the Standard Model parameters we use $`M_Z=91.187\mathrm{GeV},M_W=80.39\mathrm{GeV},G_F=1.16639\times 10^5\mathrm{GeV}^2,\alpha _s(m_t)=0.1095`$, and $`m_t=175\mathrm{GeV}`$. In the numerical evaluation we have furthermore chosen the trilinear couplings in the scalar top and bottom sector to be $`A_b=A_t`$. This fixes together with the choice of $`\mu `$ (the Higgs mixing parameter) the mixing in the scalar bottom sector. The parameter $`M`$ appearing in the plots is the $`SU(2)`$ gaugino mass parameter, it enters as an independent parameter in the full result only, see Ref. .
In Refs. it has been shown that the lightest Higgs boson mass $`m_h`$ as a function of $`X_t`$ reaches a maximum at around $`|X_t/m_{\stackrel{~}{q}}|2`$. This case we refer to as ’maximal mixing’. A minimum of $`m_h`$ is reached for $`X_t0`$, which we refer to as ’no mixing’.
Fig. 1 displays the dependence of $`m_h`$ on $`m_{\stackrel{~}{q}}`$ for the cases of no mixing and maximal mixing, and we have set $`M_A=500\mathrm{GeV}`$. For $`\mathrm{tan}\beta `$ we have chosen two typical values: $`\mathrm{tan}\beta =1.6`$ as a low and $`\mathrm{tan}\beta =40`$ as a typical high value. Very good agreement is found in the no-mixing scenario as well as in the maximal-mixing scenario, the deviations lie below $`2\mathrm{GeV}`$.
The dependence on $`M_A`$ is shown in Fig. 2. The quality of the approximation is typically better than $`1\mathrm{GeV}`$ for the no-mixing case and better than $`2\mathrm{GeV}`$ for the maximal-mixing case. Only for very small (and experimentally already excluded) values of $`M_A`$ a deviation of $`5\mathrm{GeV}`$ occurs. The peaks in the plot for $`\mathrm{tan}\beta =1.6`$ in the full result are due to the threshold $`M_A=2m_t`$ in the one-loop contribution, originating from the top-loop diagram in the $`A`$ self-energy. This peak does not occur in the approximation formula (where the momentum dependence of the $`A`$ self-energy has been neglected) and can thus lead to a larger deviation around the threshold.
In Fig. 3 we display the quality of the short analytical approximation formula for the effective mixing angle. The relative difference between $`\alpha _{\mathrm{eff}}`$ obtained from the full calculation (i.e. calculating the renormalized Higgs self-energies $`\widehat{\mathrm{\Sigma }}_s(0)`$ without any approximation) and the angle obtained with the help of the approximated Higgs self-energies, denoted as $`\alpha _{\mathrm{eff}}(\mathrm{approx})`$ is shown as a function of $`M_A`$. We use three values of $`\mathrm{tan}\beta `$, $`\mathrm{tan}\beta =3,20,40`$. Sizable deviations occur only in the region $`100\mathrm{GeV}M_A150\mathrm{GeV}`$, especially for large $`\mathrm{tan}\beta `$. In this region of parameter space the values of $`m_h`$ and $`m_H`$ are very close to each other. This results in a high sensitivity to small deviations in the Higgs boson self-energies entering the Higgs-boson mass matrix (15), (16). Otherwise the relative difference stays below 3%.
## 3 The Fortran program FeynHiggsFast
The complete program FeynHiggsFast consists of about 1300 lines Fortran code. The executable file fills about 65 KB disk space. The calculation for one set of parameters, including the $`\mathrm{\Delta }\rho `$ constraint, takes about $`2\times 10^5`$ seconds on a Sigma station (Alpha processor, 600 MHz processing speed, 512 MB RAM). The program can be obtained from the FeynHiggs home page:
http://www-itp.physik.uni-karlsruhe.de/feynhiggs .
Here the code is available, together with a short instruction, information about bug fixes etc.
FeynHiggsFast consists of a front-end (program FeynHiggsFast) and the main part where the evaluation is performed (starting with subroutine feynhiggsfastsub). The front-end can be manipulated by the user at will, whereas the main part should not be changed. In this way FeynHiggsFast can be accommodated as a subroutine to existing programs<sup>2</sup><sup>2</sup>2 This has been carried out, for example, for the program HDECAY , into which FeynHiggsFast has been incorporated as a subroutine. , thus providing an extremely fast and for many purposes sufficiently accurate evaluation for the masses and mixing angles in the MSSM Higgs sector.
The input of FeynHiggsFast are the low energy SUSY parameters, listed in Tab. 1. Concerning the $`\stackrel{~}{t}`$-sector, the user has the option to enter either the physical parameters, i.e. the masses and the mixing angle ($`m_{\stackrel{~}{t}_1}`$, $`m_{\stackrel{~}{t}_2}`$ and $`\mathrm{sin}\theta _{\stackrel{~}{t}}`$) or the unphysical parameters $`M_{\stackrel{~}{t}_L}`$, $`M_{\stackrel{~}{t}_R}`$ and $`X_t`$. From these input parameters FeynHiggsFast calculates the masses and the mixing angle of the neutral $`𝒞𝒫`$-even Higgs sector, as well as the mass of the charged Higgs boson and the $`\rho `$-parameter.
## 4 Conclusions
FeynHiggsFast is a Fortran code for the calculation of masses and mixing angles in the Higgs sector of the MSSM. In addition it evaluates the SUSY corrections to the $`\rho `$-parameter, arising from the scalar top and bottom sector. Concerning the evaluation in the neutral $`𝒞𝒫`$-even Higgs sector, FeynHiggsFast is based on a simple analytic formula, derived from a more complete result obtained in the Feynman-diagrammatic approach. Beyond the one-loop level, FeynHiggsFast contains the dominant corrections in $`𝒪(\alpha \alpha _s)`$ and further subdominant contributions. The accuracy of the approximation compared to the full result is better than $`2\mathrm{GeV}`$ for most parts of the MSSM parameter space.
The program is available via the WWW page
http://www-itp.physik.uni-karlsruhe.de/feynhiggs.
The code consists of a front-end and a subroutine. The front-end can be manipulated at the user’s will. By accommodating the front-end, FeynHiggsFast can serve as a subroutine to existing programs, thus providing an extremely fast and reasonably accurate evaluation of the masses and mixing angles in the MSSM Higgs sector. These can then be used as inputs for further computations.
### Acknowledgements
S.H. thanks the organizers of the Les Houches workshop for the inspiring and relaxed atmosphere. |
warning/0002/cond-mat0002135.html | ar5iv | text | # Non-equilibrium effects in transport through quantum dots
## I Introduction.
Quantum dots are paradigms to study the transition from macroscopic to microscopic physics. At present, the role of single-electron charging is well understood. Processes which otherwise are found in solids at the single-atom level, such as the Kondo effect, are being currently investigated. Other atomic features, like the existence of high spin ground states have also been observed. The existence of many internal degrees of freedom within the dot leads to a variety of effects reminiscent of those found in large molecules. Shake-up processes, associated with the rearrangement of many electronic levels upon the addition of one electron, have been reported.
In the present work, we will show that internal excitations of the dot lead to non-equilibrium effects which can substantially modify the transport properties. In sufficiently small dots, the addition of a single electron may cause significant charge rearrangements, and the resulting change in the electrostatic potential of the dot modifies the electronic level structure. In the limit when the level separation is much smaller than other relevant scales, this process leads to an “orthogonality catastrophe”, first discussed in relation with the sudden switching of a local potential in a bulk metal. In addition, an electron tunneling event changing the charge of the dot is associated with a charge depletion in the leads or in neighboring dots. The attraction between the electron in the dot and the induced positive charge leads to the formation of an excitonic resonance, similar to the well known excitonic effects in X-ray absorption.
The relevance of these effects for tunneling processes in mesoscopic systems was first discussed in refs. (see also). The existence of excitonic-like resonances has indeed been reported in mesoscopic systems. The formation of bonding resonances has been observed in double well systems. Non-equilibrium effects like those studied here have also been discussed in Ref. where transport through a tunnel junction via localized levels due to impurities was analyzed, and have been observed experimentally in Ref. . A different way for modifying the conductance of quantum dots through the formation of excitonic states has been proposed in. In that case the exciton is a real bound state, while the excitonic mechanism discussed here is a dynamical process. In some devices, the charge rearrangement may take place far from the tunneling region. In this case no excitons are formed but the orthogonality catastrophe persists. This process has been discussed in relation to the measurement of the charge in a quantum dot by the current through a neighboring point contact.
We will study the simplest deviations from the standard Coulomb blockade regime. Our analysis is valid for quantum dots where the spacing between electronic levels, $`\mathrm{\Delta }ϵ`$, is much smaller than the charging energy, $`E_C`$, or the temperature $`T`$. The non-equilibrium effects discussed here will be observable if, in addition, the number of electrons in the dot is not too large so that changes in the charge state lead to non-uniform redistributions of the charge. Thus, we will consider an intermediate situation between the Kondo regime and orthodox Coulomb blockade (see below for estimates). The processes that we consider will be present in double-dot devices, but, for definiteness, we study here a single dot.
The main features discussed above can be described by a generalization of the dissipative quantum rotor model, which has been studied widely in connection with conventional Coulomb blockade processes. We present a detailed numerical analysis of the generalized model, along similar lines as previous work by some of us on the dissipative quantum rotor.
The paper is organized as follows: In the next section, we show how to estimate the parameters which characterize non-equilibrium effects. The model is reviewed in section III, with emphasis on details needed for the subsequent calculations. The numerical method is presented in section IV. In Secs. V and VI, we give results for the renormalized capacitance and conductance. In Sec. VII we discuss some possible experimental evidence of the effects studied here.We close with some conclusions.
## II Non-equilibrium effects.
### A Inhomogeneous charge redistribution.
The standard Coulomb blockade model assumes that, upon a change in the charge state of the dot, the electronic levels within a quantum dot are rigidly shifted by the charging energy, $`E_C=e^2/2C`$ with $`C`$ the capacitance of the dot . Deviations from this assumption have been studied by means of an expansion in terms of $`g^1`$, the inverse dimensionless conductance, $`gk_\mathrm{F}l`$, where $`k_\mathrm{F}`$ is the Fermi wave-vector, and $`l`$ is the mean free path. It is also assumed that $`k_\mathrm{F}`$ is small compared to the inverse Thomas-Fermi screening length, $`k_{\mathrm{TF}}=\sqrt{4\pi e^2N(ϵ_\mathrm{F})/ϵ_0}`$. To lowest order beyond standard Coulomb blockade effects, the change in the charge state of the dot leads to an inhomogeneous potential, and induces a term in the Hamiltonian, which can be written as:
$$_{\mathrm{int}}=(QQ_{\mathrm{offset}})\psi ^{}(\stackrel{}{𝐫})U(\stackrel{}{𝐫})\psi (\stackrel{}{𝐫})d^d𝐫.$$
(1)
Here $`Q_{\mathrm{offset}}`$ denotes offset charges in the environment, the operator $`Q`$ measures the total electronic charge in the dot, and $`\psi ^{}(\stackrel{}{𝐫})`$ creates an electron at position $`\stackrel{}{𝐫}`$. The potential $`U(\stackrel{}{𝐫})`$ modifies the constant shift of the energy levels of the dot assumed in the standard Coulomb blockade model. It appears due to the restricted geometry of the dot. After a charge tunnels into the dot, a pile-up of electrons at the surface of the dot is induced. As a result there is a net attraction of electrons towards the surface, besides a constant shift given by $`e^2/2C`$. In general, the potential $`U`$ in eq. (1) can be obtained from the Hartree approximation (in the Thomas-Fermi limit) for dots and leads of arbitrary shape. For a spherical dot of radius $`R`$ the potential $`U(\stackrel{}{𝐫})`$ has the simple form:
$$U(\stackrel{}{𝐫})=\frac{e^2e^{k_{FT}(Rr)}}{ϵ_0k_{\mathrm{TF}}R^2}+K$$
(2)
and, for a two-dimensional circular dot:
$$U(\stackrel{}{𝐫})=\frac{e^2}{2ϵ_0k_{\mathrm{TF}}R\sqrt{R^2r^2}}+K.$$
(3)
K is a constant which ensures $`U(\stackrel{}{𝐫})=0`$. Non-equilibrium effects arise because the potential in eq. (1) is time dependent, as it changes upon the addition of electrons to the dot. Hamiltonians with terms such as eq. (1) were first discussed in (see next section). Note that the potential is localized in the surface region, where the tunneling electron is supposed to land. Finally, the potential is attractive, leading to the localization of the new electron near the surface, giving rise to excitonic effects.
Other non-equilibrium effects can arise if the charge of the dot induces inhomogeneous potentials in other metallic regions of the device. In this case, the only effect expected is the orthogonality catastrophe (see below), due to the shake-up of the electrons both in the dot and in the other regions. We take this possibility into account in the analysis in the following sections.
### B Effective tunneling density of states.
In the absence of non-equilibrium effects, the conductance of a junction between the dot and the leads is
$$g=\frac{2e^2}{h}\underset{\mathrm{channels}}{}|t_i|^2N_{i,\mathrm{lead}}(E_\mathrm{F})N_{i,\mathrm{dot}}(E_\mathrm{F}),$$
(4)
where the summation is over the channels, $`t_i`$ is the hopping matrix element through channel $`i`$, $`N(E_\mathrm{F})`$ is the density of states at the Fermi level, and we use the standard theory of tunneling in the weak transmission limit. Eq. (4) implicitly assumes a constant density of states, as appropriate for a metallic contact.
The non-equilibrium effects to be considered can be taken into account through a modification of the effective tunneling density of states. In this case the electron propagators in eq. (4) are the non-equilibrium ones, in an analogous way to the modifications required in the study of X-ray absorption spectra of core levels in metals, or tunneling between Luttinger liquids. As in those problems, we can distinguish two cases:
i) The analogue of the X-ray absorption process: The charging of the dot leads to an effective potential which modifies the electronic levels. At the same time, an electronic state localized in a region within the range of the potential is filled. The interaction between the electron in this state and the induced potential must be taken into account (the excitonic effects, in the language of the Mahan-Nozières-de Dominicis theory).
ii) The analogue of X-ray photoemission: The charging process leads to a potential which modifies the electronic levels. The tunneling electron appears in a region outside the range of this potential. Only the orthogonality effect caused by the potential needs to be included.
Taking into account the distinctions between these two possibilities, the effective (non-equilibrium) density of states in the lead and the dot becomes (omitting the channel index, $`i`$):
$`D_{\mathrm{eff}}(\omega )`$ $`=`$ $`{\displaystyle _0^\omega }𝑑\omega ^{}N_{\mathrm{dot}}^{\mathrm{empty}}(\omega ^{})N_{\mathrm{lead}}^{\mathrm{occ}}(\omega \omega ^{})`$ (5)
$``$ $`|\omega |^{1ϵ}`$ (6)
with $`ϵ`$ given by
$`ϵ=\{\begin{array}{cc}_{j=1,2}2\frac{\delta _j}{\pi }\left(\frac{\delta _j}{\pi }\right)^2\hfill & (\mathrm{excitonic}\hfill \\ & \mathrm{resonance})\hfill \\ _{j=1,2}\left(\frac{\delta _j}{\pi }\right)^2\hfill & (\mathrm{orthogonality}\hfill \\ & \mathrm{catastrophe})\hfill \end{array}`$ (11)
Here $`\delta _j`$ is the phase shift induced by the new electrostatic potential in the lead states ($`j=1`$) or in the dot states ($`j=2`$). The exponent is positive, $`ϵ>0`$, if excitonic effects prevail, while $`ϵ<0`$ if the leading process is the orthogonality catastrophe.
We can get an accurate estimate for $`ϵ`$ in the simple cases of a spherical or circular quantum dot decoupled from other metallic regions discussed in. We assume that tunneling takes place through a single channel, and the contact is of linear dimensions $`k_\mathrm{F}^1`$. In Born approximation the effective phase shift becomes:
$$\delta N(ϵ_\mathrm{F})_\mathrm{\Omega }U(\stackrel{}{𝐫})d^d𝐫\{\begin{array}{cc}\frac{e^2N(ϵ_\mathrm{F})}{ϵ_0k_{FT}k_\mathrm{F}^3R^2}\hfill & \mathrm{spherical}\mathrm{dot}\hfill \\ \frac{e^2N(ϵ_\mathrm{F})}{ϵ_0k_{FT}k_\mathrm{F}R}\hfill & \mathrm{circular}\mathrm{dot}\hfill \end{array}$$
(12)
where $`\mathrm{\Omega }`$ is the region where tunneling processes to the leads are non negligible, typically of dimensions comparable to $`k_F^1`$. Note that, for a very elongated dot (d=1), the phase shift will not depend on its linear size. As mentioned earlier, the leads can modify significantly these estimates. The tunneling electron can be attracted to the image potential that it induces, enhancing the excitonic effects ($`ϵ>0`$). On the other hand, shake-up processes in metallic regions decoupled from the tunneling processes will increase the orthogonality catastrophe, without contributing to the formation of the excitonic resonance at the Fermi energy.
## III The model.
The shake-up processes mentioned in the preceding section are described by the Hamiltonian:
$``$ $`=`$ $`_Q+_\mathrm{R}+_\mathrm{L}+_\mathrm{T}+_{\mathrm{int}}`$ (13)
$`_Q`$ $`=`$ $`{\displaystyle \frac{(QQ_{\mathrm{offset}})^2}{2C}}`$ (14)
$`_i`$ $`=`$ $`{\displaystyle \underset{k}{}}ϵ_{k,i}c_{k,i}^{}c_{k,i},i=\mathrm{L},\mathrm{R}`$ (15)
$`_\mathrm{T}`$ $`=`$ $`te^{i\varphi }{\displaystyle \underset{k,k^{}}{}}c_{k,\mathrm{R}}^{}c_{k^{},\mathrm{L}}+h.c.`$ (16)
$`_{\mathrm{int}}`$ $`=`$ $`(QQ_{\mathrm{offset}}){\displaystyle \underset{k,k^{}}{}}\left(V_{k,k^{}}^Rc_{k,\mathrm{R}}^{}c_{k^{},\mathrm{R}}V_{k,k^{}}^Lc_{k,\mathrm{L}}^{}c_{k^{},\mathrm{L}}\right)`$ (17)
Here $`[\varphi ,Q]=ie`$. The Hamiltonian separates the junction degrees of freedom into a collective mode, the charge $`Q`$, and the electron degrees of freedom of the electrodes and the dot. This separation is standard in analyzing electron liquids, where collective charge oscillations (the plasmons) are treated separately from the low-energy electron-hole excitations. In our case, this implies that only those states with energies lower than the charging energy are to be included in $`_i`$, $`_T`$ and $`_{int}`$ in eq.(LABEL:Hamiltonian). Higher electronic states contribute to the dynamics of the charge, described by $`_Q`$. The Hamiltonian, eq.(LABEL:Hamiltonian), suffices to describe transport processes at voltages and temperatures smaller than the charging energy.
In the following, we will express the offset charge $`Q_{\mathrm{offset}}`$, introduced in eq. (1), by the dimensionless parameter $`n_e=Q_{offset}/e`$. By $`V_{k,k^{}}`$ we denote the matrix elements of $`U(\stackrel{}{𝐫})`$ in the basis of the eigenfunctions near the Fermi level. We allowed that inhomogeneous potentials can be generated on both sides of the junction. We assume that tunneling can take place through several channels (index channel has been omitted). The transmission through each channel should be small for perturbation theory to apply.
The electrical relaxation associated with the tunneling process takes place in two stages. In the first, the tunneling electron is screened by the excitation of plasmons, forming the screened Coulomb potential. The time scale for this process is of the order of the inverse plasma frequency. Next, the screened Coulomb potential excites electron-hole pairs. As the electrons at the Fermi level have a much longer response time, they feel this change as a sudden and local perturbation.
We will restrict ourselves to the regime where the level spacing is small, $`\mathrm{\Delta }ϵT,E_c`$. Using standard techniques, we can integrate out the electron-hole pairs and describe the system in terms of the phase $`\varphi `$ and charge $`Q`$ only. This procedure leads to retarded interactions which are long-range in time, as the electron-hole pairs have a continuous spectrum down to zero energy. It is best to describe the resulting model within a path-integral formalism. Because of the non-equilibrium effects the effective action is a generalization of that derived for tunnel junctions
$`𝒮[\varphi ]`$ $`=`$ $`{\displaystyle _0^\beta }𝑑\tau {\displaystyle \frac{1}{4E_C}}\left({\displaystyle \frac{\varphi }{\tau }}\right)^2+`$ (20)
$`\alpha {\displaystyle _0^\beta }𝑑\tau {\displaystyle _0^\beta }𝑑\tau ^{}E_C^ϵ\left({\displaystyle \frac{\pi }{\beta }}\right)^{2ϵ}{\displaystyle \frac{1\mathrm{cos}[\varphi (\tau )\varphi (\tau ^{})]}{\mathrm{sin}^{2ϵ}[\pi (\tau \tau ^{})/\beta ]}}.`$
It describes the low energy processes below an upper cutoff of order of the unscreened charging energy, $`E_C`$. The parameter $`\alpha t^2N_R(ϵ_\mathrm{F})N_L(ϵ_\mathrm{F})`$ is a measure of the high temperature conductance, $`g_0`$, in units of $`e^2/h`$:
$$\alpha =\frac{g_0}{4\pi ^2(e^2/h)}.$$
(22)
Note that the definition of the action, eq. (LABEL:action), does not allow us to study temperatures much higher than $`E_C`$. The kernel which describes the retarded interaction is given by the effective tunneling density of states, eq. (6). The value of $`ϵ`$ is the anomalous exponent in the tunneling density of states, given in eq. (6).
The action (LABEL:action) has been studied extensively for $`ϵ=0`$, describing charging effects in the single-electron transistor in the usual limit where the electrodes are assumed to be in equilibrium. If $`ϵ>0`$, the model has a non-trivial phase transition. In this case, for $`\alpha >\alpha _{crit}2/(\pi ^2ϵ)`$, the system develops long-range order when $`T=1/\beta 0`$, leading to phase coherence and a diverging conductance.
## IV Computational method
For a given offset charge $`n_e`$, the grand partition function can be written in terms of the phase $`\varphi `$, as a path integral :
$$Z(n_e)=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{exp}(2\pi imn_e)𝒟\varphi \mathrm{exp}\left(S[\varphi ]\right),$$
(23)
where $`m`$ is the winding number of $`\varphi `$, and the paths $`\varphi (\tau )`$ satisfy in sector $`m`$ the boundary condition $`\varphi (\beta )=\varphi (0)+2\pi m`$.
The effective action and partition function can be rewritten in terms of the phase fluctuations $`\theta (\tau )=\varphi (\tau )2\pi m\tau /\beta `$, with boundary condition $`\theta (\beta )=\theta (0)`$, in the form
$$Z=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\mathrm{exp}(2\pi imn_e)I_m(\alpha ,ϵ,\beta ).$$
(24)
The coefficients $`I_m(\alpha ,ϵ,\beta )=𝒟\theta \mathrm{exp}\left(S_m[\theta ]\right)`$ are to be evaluated with the effective action $`S_m[\theta (\tau )]=S[\theta (\tau )+2\pi m\tau /\beta ]`$. They depend on the winding number $`m`$, the temperature, and the dimensionless parameters $`\alpha `$ and $`ϵ`$, but are independent of the offset charge $`n_e`$. This means that the problem reduces, from a computational point of view, to the calculation of the relative values of $`I_m(\alpha ,ϵ,\beta )`$, which can be obtained from Monte Carlo (MC) simulations apart from an overall normalization constant . The partition function is even and periodic with respect to $`n_e`$, $`Z(n_e)=Z(n_e)=Z(n_e+1)`$, and therefore one can restrict the analysis to the range $`0n_e0.5`$.
The MC simulations have been carried out by the usual discretization of the quantum paths into $`N`$ (Trotter number) imaginary-time slices . In order to keep roughly the same precision in the calculated quantities, as the temperature is reduced, the number of time-slices $`N`$ has to increase as $`1/T`$. We have found that a value $`N=4\beta E_C`$ is sufficient to reach convergence of $`I_m`$. Therefore, the imaginary-time slice employed in the discretization of the paths is $`\mathrm{\Delta }\tau =\beta /N=1/(4E_C)`$.
When discretizing the paths $`\varphi (\tau )`$ into $`N`$ points, it is important to treat correctly the $`|\tau \tau ^{}|0`$ divergence that appears in the tunneling term $`S_t[\varphi ]`$ of the effective action \[second term on the r.h.s. of Eq. (LABEL:action)\]. This divergence can be handled as follows. In the discretization procedure, the double integral in $`S_t[\varphi ]`$ translates into a sum extended to $`N^2`$ two-dimensional plaquettes, each one with area $`(\mathrm{\Delta }\tau )^2`$. The above-mentioned divergence appears in the $`N`$ “diagonal” terms ($`\tau =\tau ^{}`$) and can be dealt with by approximating the integrand close to $`\tau =\tau ^{}`$ by $`E_C^ϵ|\tau \tau ^{}|^ϵ(d\varphi /d\tau )^2/2`$. Thus, by integrating this expression over the “diagonal” plaquette $`(i,i)`$, with $`1iN`$, one finds that its contribution to $`S_t[\varphi ]`$ is given by
$$\mathrm{\Delta }S_t(\tau _i,\tau _i)=2E_C^ϵ\frac{\alpha }{ϵ+1}\left(\frac{\mathrm{\Delta }\tau }{2}\right)^{ϵ+2}\left(\frac{d\varphi }{d\tau }\right)_{\tau =\tau _i}^2,$$
(25)
which is regular for $`ϵ1`$. The error introduced by this replacement in the discretization procedure is of the same order as that introduced by the usual discretization of the “non-diagonal” terms. We have checked that the results of our Monte Carlo simulations obtained by using this procedure converge with the Trotter number $`N`$.
The partition function in Eq.(24) has been sampled by the classical Metropolis method for temperatures down to $`k_\mathrm{B}T=E_C/200`$. A simulation run proceeds via successive MC steps. In each step, all path-coordinates (imaginary-time slices) are updated. For each set of parameters ($`\alpha ,ϵ,T`$), the maximum distance allowed for random moves was fixed in order to obtain an acceptance ratio of about $`50\%`$. Then, we chose a starting configuration for the MC runs after system equilibration during about $`3\times 10^4`$ MC steps. Finally, ensemble-averaged values for the quantities of interest were calculated from samples of $`1\times 10^5`$ quantum paths. More details on this kind of MC simulations can be found elsewhere .
## V Renormalization of the capacitance.
We will study the capacitance renormalization for tunneling conductance $`\alpha >0`$ by calculating the effective charging energy $`E_C^{}(T)=e^2/2C^{}(T)`$, which can be obtained as a second derivative of the free energy $`F=k_\mathrm{B}T\mathrm{ln}Z`$:
$$E_C^{}(T)=\frac{1}{2}\frac{^2F}{n_e^2}|_{n_e=0}.$$
(26)
At high temperatures, the free energy $`F(n_e)`$ depends weakly on $`n_e`$, and the curvature \[i.e., $`E_C^{}(T)`$\] approaches zero. At low $`T`$, and for weak tunneling ($`\alpha 1`$), it coincides with the usual charging energy $`E_C`$. By using Eqs. (23) and (26) this renormalized charging energy can readily be expressed as
$$E_C^{}(T)=2\pi ^2k_\mathrm{B}Tm^2_{n_e=0},$$
(27)
where $`m^2_{n_e=0}`$ is the second moment of the coefficients $`I_m`$.
The correlation function in imaginary time $`G(\tau )`$, that will be used below to calculate the conductance, can be calculated from the MC simulations as $`G(\tau )=\mathrm{cos}[\varphi (\tau )\varphi (0)]`$. This means in our context:
$`G(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{m=\mathrm{}}{\overset{\mathrm{}}{}}}\mathrm{exp}(2\pi imn_e)\times `$ (29)
$`{\displaystyle 𝒟\varphi e^{S[\varphi ]}\mathrm{cos}[\varphi (\tau )\varphi (0)]}.`$
A number of features, directly related to the free energy of the model given by action (LABEL:action), are reasonably well understood for $`ϵ=0`$. The effective charge induced by an arbitrary offset charge, $`n_e`$, has been extensively analyzed. A number of analytical schemes give consistent results in the weak coupling ($`\alpha 1`$) regime. These calculations have been extended to the strong coupling limit by numerical methods.
In fig. (1), we present results for the effective charging energy $`E_C^{}`$ as a function of temperature, and for different values of $`ϵ`$. The value of $`E_C^{}`$ is enhanced for $`ϵ<0`$, where the orthogonality catastrophe dominates the physics. A positive $`ϵ`$ reduces the effective charging energy, and, beyond some critical value (see discussion below), $`E_C^{}`$ scales towards zero as the temperature is decreased, showing non-monotonic behavior. The same trend can be appreciated in fig. (2), where the effective charging energy at low temperatures is plotted as a function of $`\alpha `$. Renormalization group arguments show that $`E_C^{}(\alpha ,ϵ)`$ should go to zero for $`\alpha >\alpha _{crit}(ϵ)`$. We have checked the consistency of this prediction with the numerical results by fitting the values of $`E_C^{}(\alpha )`$ at low temperatures by the expression expected from the scaling analysis near the transition:
$$E_C^{}(\alpha ,ϵ)=\left[1\frac{\alpha }{\alpha _{crit}(ϵ)}\right]^{\frac{1}{ϵ}}$$
(30)
In this expression we use $`\alpha _{crit}(ϵ)`$ as the only adjustable parameter. The results are shown in fig. (3). Note that the same equation (30) can be used to fit the results for $`ϵ<0`$, if one uses a negative value for $`\alpha _{crit}`$. There is no phase transition, however, for $`ϵ<0`$.
The results show that the present calculations are very accurate even for relatively large values of $`\alpha `$, where $`E_C^{}`$ converges at very low temperatures.
## VI Evaluation of the conductance.
### A $`ϵ=0`$
The conductance of the single-electron transistor is notoriously more difficult to calculate than standard thermodynamic averages. It cannot be derived in a simple fashion from the partition function, and requires the analytical continuation of the response functions from imaginary to real times or frequencies. Hence, there are no comprehensive results valid for the whole range of values of $`\alpha ,T/E_C`$ and $`n_e`$. For $`ϵ=0`$, the conductance $`g(T)`$ can be written as:
$$g(T)=2g_0\beta _0^{\mathrm{}}\frac{d\omega }{2\pi }\frac{\omega S(\omega )}{e^{\beta \omega }1}$$
(31)
where $`g_0`$ is the normal state conductance, and $`S(\omega )`$ is related to the correlation function in imaginary time:
$`G(\tau )`$ $`=`$ $`e^{i\varphi (\tau )}e^{i\varphi (0)}`$ (32)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{\mathrm{}}}𝑑\omega {\displaystyle \frac{e^{(\beta \tau )\omega }+e^{\tau \omega }}{e^{\beta \omega }1}}A(\omega )`$ (33)
and
$$A(\omega )=(1e^{\beta \omega })S(\omega )$$
(34)
In the previous expressions, the charging energy is the natural cutoff for the energy integrals.
At high temperatures, $`\beta E_C1`$, we can expand in eq. (31):
$$g(T)2g_0\beta _0^{\mathrm{}}\frac{d\omega }{2\pi }\left[\frac{1}{\beta }\frac{\beta \omega ^2}{24}+\frac{7\beta ^3\omega ^4}{5760}+\mathrm{}\right]\frac{e^{(\beta \omega )/2}A(\omega )}{e^{\beta \omega }1}$$
(35)
so that:
$$g(T)g_0\left[G\left(\frac{\beta }{2}\right)\frac{\beta ^2}{24}G^{\prime \prime }\left(\frac{\beta }{2}\right)+\frac{7\beta ^4}{5760}G^{iv}\left(\frac{\beta }{2}\right)+\mathrm{}\right]$$
(36)
At low temperatures, $`\beta E_C1`$, the conductance is dominated by the low energy behavior of $`A(\omega )`$ or, alternatively, $`S(\omega )`$. To lowest order, we expect an expansion of the form:
$$S(\omega )2\pi \delta (\omega \omega _0)+A|\omega |+\mathrm{}$$
(37)
where $`\omega _0`$ is an energy of the order of the renormalized charging energy (or zero at resonance, $`n_e1/2`$), and $`A`$ is a constant which describes cotunneling processes. Inserting this expression in eq. (31), we obtain:
$$g(T)g_0\frac{2\beta \omega _0}{e^{\beta \omega _0}1}+g_0\frac{A}{\pi \beta ^2}_0^{\mathrm{}}𝑑x\frac{x^2}{e^x1}+\mathrm{}$$
(38)
while, on the other hand:
$$G\left(\frac{\beta }{2}\right)2e^{\frac{1}{2}\beta \omega _0}+\frac{A}{\pi }\left(\frac{2}{\beta }\right)^2+\mathrm{}$$
(39)
If $`\omega _0=0`$, both $`g`$ and $`g_0G(\beta /2)`$ have the same limit, as $`T0`$. When $`\omega _00`$ the leading term goes as $`T^2`$ (cotunneling) in both cases, with prefactors equal to $`2.404A/\pi `$ and $`4A/\pi `$, respectively.
From the above discussion of the relation between the high- and low-temperature behavior of $`g(T)`$ and $`G(\beta /2)`$, we find that the interpolation formula
$$g(T)g_0G\left(\frac{\beta }{2}\right)$$
(40)
should give a reasonable approximation over the entire range of parameters (note that the above discussion is independent of the values of $`\alpha `$ and $`n_e`$). Eq. (40) is consistent with the the main physical features expected both in the high and low temperature limits, at and away from resonance. The advantage of using $`G(\beta /2)`$ is that it can be computed, to a high degree of accuracy, by standard Monte Carlo techniques, as it does not require to continue the results to real times. A similar approximation, used to avoid inaccurate analytical continuations has been applied for bulk systems in Ref. .
We show the adequacy of the approximation, eq. (40), by plotting the conductances estimated in this way, as a function of the bias charge $`n_e`$, in fig. (4).
The minimum ($`n_e=0`$) and maximum conductances ($`n_e=1/2`$) for different values of $`\alpha `$ and temperatures are shown in fig. (5).
### B $`ϵ0`$
We now extend the previous approximation to the conductance, eq. (40) to the case $`ϵ0`$. The main modification in eq. (31) is that a factor $`\omega /(1e^{\beta \omega })`$ within the integral has to be replaced by the effective tunneling density of states, $`D_{\mathrm{eff}}`$, given, at zero temperature, by eq. (6). At finite temperatures, the corresponding expression is approximately
$$D_{\mathrm{eff}}(\omega )\text{max}[T(T/E_C)^ϵ,\omega (\omega /E_C)^ϵ].$$
(41)
The relevant range in the integrand in eq. (31) is from $`\omega =0`$ to $`\omega T`$. In the following, we will factor the $`ϵ`$ dependent part of the effective density of states, and we write the generalization of eq. (6) to finite temperatures as:
$$D_{\mathrm{eff}}(\omega )=\frac{\omega }{1e^{\beta \omega }}D_{res}(ϵ,\omega )$$
(42)
Finally, when inserting this expression into eq. (31), we make the approximation:
$$D_{res}(ϵ,\omega )\left(\frac{T}{E_C}\right)^ϵ$$
(43)
With this approximation, we can perform the same analysis in the high and low temperature regimes as before, to obtain the interpolation formula:
$$g(T)g_0\left(\frac{T}{E_C}\right)^ϵG\left(\frac{\beta }{2}\right)$$
(44)
This expression includes again the relevant physical processes at high and low temperatures.
Results for the maximum and minimum values of the conductances, for different values of $`ϵ`$, are presented in fig. (6). In the non phase-coherent regime, at very low temperatures, $`TE_C^{}`$, the conductance away from resonance should vary as $`gT^{22ϵ}`$. Exactly at resonance, $`n_e=1/2`$, the conductance diverges as $`T^ϵ`$. The most interesting result is the divergence of the conductance, at low temperatures, for $`ϵ=0.5`$, where the excitonic effects are strong enough to drive the system to the phase-coherent phase.
The full conductance, as a function of $`n_e`$, is shown in fig. (7), for $`\alpha =0.25`$ and $`ϵ=0.5`$. As mentioned above, for these parameters the system is already in the phase coherent regime. The conductance behaves in a way similar to that in the usual case ($`ϵ=0`$), and a peaked structure develops. The absolute magnitude, however, increases as the temperature is lowered. Note that the effective charging energy is finite as $`T0`$ (see fig. (1)). It is interesting to note that, in this phase with complete suppression of Coulomb blockade effects at low temperatures (high values of $`ϵ`$ and high conductances), the peak structure appears only for an intermediate range of temperatures, and it is washed out at very low temperatures.
In fig. (8) we show the conductance as a function of $`ϵ`$ for a fixed temperature. It is evident the increase of the conductance as $`ϵ`$ increases.
## VII Discussion.
In the following, we discuss some experimental evidence which can be explained within the model discussed here.
It has been pointed out that the correlations between the conductances for neighboring charge states of a quantum dot are too weak to be explained using standard methods for disordered, non-interacting systems. The experiments reported in are in the cotunneling regime. In the presence of non-equilibrium processes, we expect a behavior of the type $`gT^{22ϵ}`$. Note that $`ϵ`$ is determined by phase shifts, which depend on microscopic details of the contacts. Thus, it can be expected to vary with the charge state of the dot, and to lead to large differences in the conductances of neighboring valleys.
It has been shown that the temperature dependence of the conductance quantum dots, away from resonances, can be opposite to that expected in a system exhibiting Coulomb blockade. This effect could not be attributed to Kondo physics, as the data do not show an even-odd alternation. The reported behavior can be explained within our model, assuming that the value of $`ϵ`$ is sufficiently large, and dependent on the charge state of the dot.
In ref. the inelastic contribution to the conductance in a double dot sytem, where the electronic states in the two dots are separated by an energy $`ϵ`$ is measured. The result is approximately given by $`I_{\mathrm{inel}}(ϵ)ϵ^\lambda `$, where $`\lambda `$ is a negative constant of order unity. Taking into account only one electronic state within each dot, the problem can be reduced to that of a dissipative, biased two level system. The inelastic conductance reflects the nature of the low energy excitations coupled to the two level system. The observed power law decay implies that the spectral strength of the coupling, that is the function $`J(\omega )`$ in the standard literature, should be ohmic, $`J(\omega )|\omega |`$. This has led to the proposal that the excitations coupled to the charges in the double dot system are piezoelectric phonons. It is interesting to note that the excitation of electron-hole pairs leads also to an ohmic spectral function. Thus, at sufficiently low energies, an orthogonality catastrophe due to electron-hole pairs shows a behavior indistinguishable from that arising from piezoelectric phonons. The contributions from the two types of excitations can be distinguished at the natural cutoff scale for phonons, which is the energy of a phonon whose wavelength is of the order of the dimensions of the device. On the other hand, the simplest prediction for the expected behavior of the current induced by the emission of the electron-hole pairs is:
$$I(ϵ)K\frac{1}{1e^{\beta \sqrt{\mathrm{\Delta }_0^2+ϵ_b^2}}}\frac{\mathrm{\Delta }_0^2}{\sqrt{\mathrm{\Delta }_0^2+ϵ_b^2}}$$
(45)
where $`\mathrm{\Delta }_0`$ is the tunneling element between the two dots, $`ϵ_b`$ is the bias, and $`K`$ is the coupling constant (referred to as $`ϵ`$ in other sections of this paper). This expression gives the absorption rate of a dissipative two level system in the weak coupling regime, $`K1`$. The natural cutoff for electron-hole pairs is bounded by the charging energy of the system. It would be interesting to disentangle the relative contributions of electron-hole pairs and piezoelectric phonons to the inelastic current.
The photo-induced conductance in a double dot system has also been measured. The analysis of the contribution of inelastic processes due to electron hole pairs to the measured conductance proceeds in the same way as in the interpretation of the previous experiment. Let us suppose that a photon of energy $`\omega _{ph}`$ excites an electron within one dot. Assuming that the coupling to the environment is weak, the rate at which the electron tunnels to the second dot by losing an energy $`ϵ`$ is given by eq.(45). The experiments in suggest that the number of states within each dot are discrete. Then, the induced conductance should show a series of peaks, related to resonant photon absorption within one dot. The height of each peak is determined by the decay rate to lower excited states in the other dot, and it can be written as a sum of terms with the dependence given in eq.(45), where $`ϵ`$ is the energy difference between the initial and final states. The envelope of the spectrum should look like a power law, in qualitative agreement with the experiments.
It is interesting to note that bunching of energy levels in quantum dots have been reported. The separation between peaks defines the charging energy, which, according to the experiments, vanishes for certain charge states. This behavior can be explained if the excitonic effects drive the quantum dot beyond the transition, and charging effects are totally suppressed. This mechanism can also play some role in the observed transitions in granular wires.
## VIII Conclusions.
We have analyzed the effects of non-equilibrium transients after a tunneling process on the conductance of quantum dots. They are related to the change in the electrostatic potential of the dot upon the addition of a single electron. These effects can enhance or suppress the Coulomb blockade. The most striking effect arise from the formation of an exciton-like resonance at the Fermi level after the charging process and lead to the complete suppression of the Coulomb blockade and a diverging conductance at low temperatures. It appears for sufficiently large values of the conductance, $`\alpha `$ , and the non-equilibrium phase shifts which define $`ϵ`$ in our model. The same potential which leads to this dynamic resonance plays a role in the deviations of the level spacings from the standard Coulomb blockade model.
For simplicity, we have considered the simplest case, a single-electron transistor. The analysis reported here can be extended, in a straightforward fashion to other devices, like double quantum dots, where the effects described here should be easier to observe.
As discussed, the non-equilibrium effects considered in this paper may have been observed already in the conductance of quantum dots.
## IX Acknowledgments
One of us (E. B.) is thankful to the University of Karlsruhe for hospitality. We acknowledge financial support from CICyT (Spain) through grant PB0875/96, CAM (Madrid) through grant 07N/0045/1998 and FPI, and the European Union through grant ERBFMRXCT960042. |
warning/0002/quant-ph0002011.html | ar5iv | text | # Time of arrival in the presence of interactions
## I Introduction
In this paper we work out a theoretical framework to compute the time in which a particle that moves through an interacting medium arrives at a given point. In the construction of this framework we will have to deal with problems of very different kind that we introduce now:
First, there is the nature of time in quantum mechanics. It appears as the external evolution parameter in the Schrödinger and Heisenberg equations, common to both, systems and observers alike. However, time arises in many instances (transitions, decays, arrivals, etc.) as a property of the physical systems. The attempts to promote time to the category of observable run early into the obstruction detected by Pauli : A self-adjoint time operator implies an unbounded energy spectrum. This was soon related to the uncertainty relation for time and energy, whose status and physical meaning produced some controversies , and is still subject of elucidation today (see for instance and ). The question remains unsettled for closed quantum systems, specially in the case of quantum gravity, whose formulation is pervaded by the so called problem of time .
Second, the definition of the time-of-arrival (toa), which is probably the simplest candidate time to become a property of the (arriving) physical system, rather than a mere external parameter. Due to its conceptual simplicity, it has been used in many cases to illustrate different problems related to the role of time in quantum theory. Allcock analyzed extensively the difficulties met by the toa concluding that they were insurmountable. The present situation is ambiguous. On the one side, there are theoretical analysis of the toa suggesting that it can not be precisely defined and measured in quantum mechanics. This contradicts the possibility of devising high efficiency absorbers , that could be used as almost ideal detectors for toa . On the other hand, there are explicit constructions of a self-adjoint (albeit in a pre-Hilbert space) toa operator for the non-relativistic free particle in one space dimension , and for the relativistic free particle in 3-D , both avoiding the Pauli problem. There is also an alternative formulation as a Positive Operator Valued Measure (POVM). Finally, the toa has been measured in high precision experiments on the arrival of two entangled photons produced by parametric down-conversion, one of which has undergone tunneling through a photonic band gap (PBG). The experimental results that show superluminal tunneling, neatly identify the Hartman effect and the Wigner time delay (or phase time) as the physically relevant mechanisms for the tunneling time and toa respectively. Whether these results apply only to photons and are due to the specific properties of the PBG used, or can be extended to other particles and barriers, can not be decided in the lack of a satisfactory theory of the toa at a space point through interacting media.
The third question is thus the tunneling time, for which there are three main proposals. Wigner introduced the phase time in his analysis of the relationship between retardation, interaction range, and scattering phase shifts. Buttiker and Landauer introduced the traversal time in their study of tunneling through a time-dependent barrier. Soon after, Buttiker used the Larmor precession as a clock , identifying the dwell , traversal, and reflection times as three characteristic times describing the interaction of particles with a barrier. Recent reviews that include these and other approaches, discussing toa and tunneling times from a modern, unified perspective, can be found in and . The light shed on these questions by the two photon experiments is revised in and .
The main progress, quoted before, towards the formulation of a quantum toa operator has been its explicit construction for a particle moving freely in one space dimension. In this paper we face the problem of extending this formalism to the case of the presence of an interaction potential affecting a region of the (1-D) space. We extend here our previous unpublished results , taking now proper care of the dependence on the arrival position $`x`$ that we consider placed in front of the interaction region, within it, and behind it (if spatially finite).
The plan of the work is as follows: In Section II we construct a toa formalism of general validity. The starting point is the case of the free particle. There, a suitable canonical transformation, the quantum version of the Jacobi-Lie transformation of classical mechanics, gives the toa in interacting media, (even at points where $`V(x)0`$). In Section III we consider an initial state consisting of a narrow Gaussian wave packet prepared at the left of the interaction region and moving towards it. A quasi-classical study of the toa at a point $`x`$ in the interaction region is first carried out. Then we turn to the full quantum mechanical treatment. We first analyze the arrival in the presence of a step potential. Section IV is devoted to the study of the toa at points behind square barriers. We detect, in diferent instances, stauration of and departures from the Hartman effect. The case of $`x`$ at the left of the interacting region is characterized by the (possibly interfering) contributions coming from the incident and the reflected wave packets. This situation is treated in Section V, where we deal separately with the case of total reflection (very high barriers) which has an analog in classical mechanics, and with the case of partial reflection, a pure quantum phenomenon with very rich structure in the time domain. Finally, we summarize our results in Section VI.
In Appendix A we show how the toa can be treated as a derived quantity in the phase space of classical Hamiltonian systems in the case where these are integrable. A short review of the construction and properties of the quantum toa operator for free particles is presented in Appendix B for completness.
## II Time of arrival formalism
To measure the time of arrival of a free particle at a point $`x`$ one would: a) place a detector at $`x`$, b) prepare the initial state $`|\psi `$ of the particle at $`t=0`$, and then, c) record with a clock the time $`t`$ when the detector clicks. The value of $`t`$ gives the toa of the state $`|\psi `$ at $`x`$. Repeating this procedure with identically prepared initial states, one would get the probability distribution in times of arrival at $`x`$. Of course, the results would depend on the initial state chosen, which stores all the information regarding the initial distribution in positions and momenta of the particle.
We want to determine the effect on these times of a position dependent interaction between the particle and the medium, that we describe by a potential energy $`V(q)`$. For instance, to disclose the effect of climbing (or tunneling through) a potential barrier, one would simply put the barrier in between the detector and the initial state, and then record the new times of arrival. With an initial state identical to that prepared for the free case, any difference in the probability distributions should be an effect of the barrier. Several questions can be investigated by changing the properties of the barrier: its height or width if it is rectangular, even its very form. This has been explicitly done in the two photon experiments at Berkeley, by putting alternatively a mirror and an ordinary glass in the path of one of the photons. Of great interest is the dependence on $`x`$ where $`V(x)>0`$, i.e. with the detector within the range of the interaction, and also the time of arrival for $`E>V`$, in which classically there is no reflection. Some of these questions are studied in last sections of this paper.
In classical mechanics particles move along the trajectories $`H(q,p)=`$const. as $`t`$ increases. This allows to work out $`t_x`$, the time of arrival at the point $`q(t)=x`$, by identifying the point $`(q,p)`$ of phase space where the particle is at (say) $`t=0`$, and then by following the trajectory that passes by it, up to the arrival at $`x`$. The mathematical translation of this procedure is given by the equation of time:
$$t_x(q,p)=\text{sign}(p)\sqrt{\frac{m}{2}}_q^x\frac{dq^{}}{\sqrt{H(q,p)V(q^{})}}$$
(1)
that is discussed at length in many textbooks, and whose existence conditions and characterization as a function of the phase space variables are outlined in the Appendix A. We simply note here that $`t_x(q,p)`$ is canonically conjugate to the hamiltonian $`\{t_x(q,p),H(q,p)\}=1`$.
This equation is a troublesome starting point for quantization. First, it involves a (path) integral of operators and should be treated accordingly. Second, it only applies to values of $`x`$ that are classically within the reach of $`(q,p)`$, while in quantum mechanics all values of $`x`$ are attainable. Classically, the particle propagates without reflection up to the turning point $`q_0`$ ($`V(q_0)=E`$), where it is completely reflected. There is no further penetration beyond this point. The situation is different in quantum mechanics: there may be tunneling beyond $`q_0`$ and partial reflection before reaching it. These phenomena cannot be accounted for by Eq. (1), that gives complex numbers for these cases. Now, note that both, tunneling and partial reflection, are absent from the motion of free particles, whose time of arrival has been successfully quantized as said in the introduction. In addition, all the positions are within the reach of the free particle. Summarizing, everything points to the free time of arrival as a main clue to solve the problem.
In this work we desist from attempting the straightforward quantization of the classical expression (1). Instead, we will construct the solution to the interacting case taking as starting point the well known results that apply to the free case. The aim is to produce the quantum version of the Lie transformation from the actual flow in phase space to the canonically equivalent parallel flow of constant velocity translations. In other words, we shall use the quantum version of the canonical transformation to action-angle variables. The Lie procedure – that we sketch for completness in Appendix A – has a property that will be the crux of the matter in our construction. Namely, it permits to define time as a derived variable in phase space in terms of the free action-angle variables as well as, alternative and equivalently, in terms of the original positions and momenta. Obviously, both definitions give the same result as we show explicitly in Eq. (87). Our use of the Lie procedure in the quantum case can be described as the combination of steps a,b and c below:
* The quantization of the time of arrival $`𝐭_\mathrm{𝟎}`$ of the free particle. This is an old problem in quantum mechanics, whose solution in terms of a Positive Operator Valued Measure we describe in Appendix B.
* The construction of the quantum canonical transformation $`U`$ that connects the free particle dynamics with Hamiltonian $`H_0`$ to the case of interest with Hamiltonian $`H_0+V(q)`$. $`U`$ is given by the Möller wave operator as we show in Sections 2.1 and 2.2.
* The application of the canonical transformation $`U`$ to $`𝐭_\mathrm{𝟎}`$ to get the time of arrival $`𝐭`$ in the presence of the interaction potential $`V(q)`$, that is $`𝐭=U𝐭_\mathrm{𝟎}U^{}`$. This is what we do in Section 2.2, where we also address the interpretation of the resulting formalism.
### A Implicit quantum canonical transformations
Classical canonical transformations $`\overline{q}=\overline{q}(q,p)`$, $`\overline{p}=\overline{p}(q,p)`$ in phase space can be defined implicitly by the use of auxiliary functions $`F,G,\overline{F},\overline{G}`$ in the following way:
$`\overline{F}(\overline{q},\overline{p})`$ $`=`$ $`F(q,p)`$ (2)
$`\overline{G}(\overline{q},\overline{p})`$ $`=`$ $`G(q,p).`$ (3)
It is easy to work out the following relation among Poisson brackets:
$$\{\overline{F},\overline{G}\}_{\overline{q}\overline{p}}\{\overline{q},\overline{p}\}_{qp}=\{F,G\}_{qp}$$
(4)
In these conditions, the transformation is canonical (i.e. $`\{\overline{q},\overline{p}\}_{qp}=1`$) if and only if
$$\{\overline{F},\overline{G}\}_{\overline{q}\overline{p}}=\{F,G\}_{qp}.$$
(5)
This relation has the additional property of fixing one of the four functions $`F,G,\overline{F},\overline{G}`$, once the other three are given. We can choose $`F`$ and $`G`$ as the free particle Hamiltonian and time of arrival respectively. Then, if $`\overline{F}`$ is the complete Hamiltonian $`H`$, $`\overline{G}`$ will be the corresponding toa $`t_x`$ given by (1) along the classical trajectories.
Canonical transformations were introduced by Dirac in quantum mechanics by the use of unitary transformations $`U`$ ($`UU^{}=U^{}U=1`$ ). If the operators $`\overline{𝐪},\overline{𝐩}`$ are canonically transformed from $`𝐪,𝐩`$ , then there is a unitary transformation $`U`$ such that
$`\overline{𝐪}`$ $`=`$ $`U^{}𝐪U`$ (6)
$`\overline{𝐩}`$ $`=`$ $`U^{}𝐩U.`$ (7)
Then one can define implicitly quantum canonical transformations, like the classical ones. This possibility has been thoroughly analyzed and developed. The main results of the method are collected in , where one can also find references to other relevant literature. The transformation $`U`$ is given by
$`\overline{𝐅}(\overline{𝐪},\overline{𝐩})`$ $`=`$ $`U^{}\overline{𝐅}(𝐪,𝐩)U=𝐅(𝐪,𝐩)`$ (8)
$`\overline{𝐆}(\overline{𝐪},\overline{𝐩})`$ $`=`$ $`U^{}\overline{𝐆}(𝐪,𝐩)U=𝐆(𝐪,𝐩)`$ (9)
where the last equality in each row is the definition of the barred operators in terms of the unbarred ones, while the first equality comes from the straight application of (7) to the l.h.s. Being $`U`$ a unitary transformation, the spectra of the canonically transformed operators have to coincide, that is:
$`\sigma (\overline{𝐪})=\sigma (𝐪)=,`$ $`\sigma (\overline{𝐩})=\sigma (𝐩)=`$ (10)
$`\sigma (\overline{𝐅})=\sigma (𝐅),`$ $`\sigma (\overline{𝐆})=\sigma (𝐆)`$ (11)
where the second row stands because $`𝐅`$ and $`\overline{𝐅}`$, $`𝐆`$ and $`\overline{𝐆}`$ are also unitarily related operators.
The above relations permit to build the operator $`\overline{𝐆}`$ once $`𝐅,𝐆`$ and $`\overline{𝐅}`$ are given. We assume that $`𝐅`$ and $`\overline{𝐅}`$ are self-adjoint operators, with the eigenstates corresponding to the same eigenvalue $`\lambda _f`$ given by:
$$\overline{𝐅}|\overline{f}=\lambda _f|\overline{f},𝐅|f=\lambda _f|f$$
(12)
They form orthogonal and complete bases satisfying
$`fs|f^{}s^{}=`$ $`\delta _{ss^{}}\delta (\lambda _f\lambda _f^{}),{\displaystyle \underset{s}{}}{\displaystyle _{\sigma (\lambda )}}𝑑\lambda _f|fsfs|`$ $`=\mathrm{I}\mathrm{I}`$ (13)
$`\overline{f}s|\overline{f}^{}s^{}=`$ $`\delta _{ss^{}}\delta (\lambda _f\lambda _f^{}),{\displaystyle \underset{s}{}}{\displaystyle _{\sigma (\lambda )}}𝑑\lambda _f|\overline{f}s\overline{f}s|`$ $`=\mathrm{I}\mathrm{I}`$ (14)
where we allow for some degeneracy (that has to be the same for both $`𝐅`$ and $`\overline{𝐅}`$) labeled by $`s`$. We have also assumed that $`\lambda `$ is continuous, while $`s`$ is a discrete index. These assumptions could be changed straightforwardly if it were necessary. Now, an operator $`U`$ satisfying the first row of Eq.(9) can be given simply as:
$$U=\underset{s}{}_{\sigma (\lambda )}𝑑\lambda _f|\overline{f}sfs|$$
(15)
It is straightforward to verify that it is unitary. We can now proceed to the sought for result: the definition of $`\overline{𝐆}`$ in terms of $`𝐆`$ using $`U`$, that is $`\overline{𝐆}=U𝐆U^{}`$. The full fledged expression is
$$\overline{𝐆}(𝐪,𝐩)=\underset{ss^{}}{}_{\sigma (\lambda )}𝑑\lambda _f𝑑\lambda _f^{}|\overline{f}sfs|𝐆(𝐪,𝐩)|f^{}s^{}\overline{f}^{}s^{}|$$
(16)
that constitutes our main result in the quantum canonical formalism.
### B Definition of the time of arrival
We will now apply the above to the case where $`𝐅`$ is the free Hamiltonian $`H_0`$, $`\overline{𝐅}`$ the complete Hamiltonian $`H`$ and $`𝐆`$ the time of arrival of the free particle Eq.(88). Then, we have $`H_0=U^{}HU`$ and $`\mathrm{\Pi }_0(x)=U^{}\mathrm{\Pi }(x)U`$. In Appendix B we have summarized the toa formalism for the free particle, given by the positive operator valued measure $`P_0`$ of Eq.(93). Accordingly, the POVM $`P`$ of the interacting case will be given by (cf (9))
$$P(\mathrm{\Pi }(x);t_1,t_2)=UP_0(\mathrm{\Pi }_0(x);t_1,t_2)U^{}.$$
(17)
Finally, the time of arrival operator in the presence of interactions (the $`\overline{𝐆}`$ of our problem) is given by
$$𝐭(H,\mathrm{\Pi }(x))=U𝐭_0(H_0,\mathrm{\Pi }_0(x))U^{}.$$
(18)
Three comments are in order here:
* Fixing $`U`$ by the relation between both hamiltonians leads to two different solutions:
$$U_{(\pm )}=\underset{s}{}_0^{\mathrm{}}𝑑E|Es(\pm )Es0|=\mathrm{\Omega }_{(\pm )}$$
(19)
which are the Möller operators connecting the Hilbert space $`_{in}`$ and $`_{out}`$ of free particle states to the Hilbert space $``$ of the bound and scattering states. These operators are only isometric in the presence of bound states, because the correspondence between states in $``$ and free states can not be one to one. In this paper we will consider only well behaved potentials ($`V(q)0q`$), that vanish at the spatial infinity, for which the Möller operators are unitary because there is one free state for each scattering state. In this case, the intertwining relations $`H\mathrm{\Omega }_\pm =\mathrm{\Omega }_\pm H_0`$ can be put in the usual form $`H=\mathrm{\Omega }_\pm H_0\mathrm{\Omega }_\pm ^{}`$. In addition, we shall adhere to the standard conventions, choosing $`\mathrm{\Omega }_{(+)}`$ (with $`E=lim_{ϵ0^+}(E+iϵ)`$) in (19) that gives signal propagation forward in time. The results that would be obtained with $`\mathrm{\Omega }_{()}`$ would correspond to the time reversal of the actual situation. If $`\tau `$ is the time reversal operator $`P_{()}(\mathrm{\Pi }(x);t_1,t_2)=\tau P_{(+)}(\mathrm{\Pi }(x);t_2,t_1)\tau ^{}`$. For notational simplicity, we will omit this label $`(+)`$ wherever possible.
* The reduction of the problem to a sort of free particle problem by means of a canonical transformation as done in (18) should not be a surprise. On the contrary, this is the quantum counterpart of the classical situation where the trajectories of completely integrable phase space flows can be straightened out to those of a free particle by means of a canonical transformation. To the classical Lie transformation of Appendix A that carries out this stretching corresponds the quantum transformation described above and in the previous Section 2.1. Concretely, Equation (87) is the classic analog to (18).
* $`x`$ is the actual detector position in the interacting case. Therefore, the arguments of $`𝐭`$ in (18) have to be $`\mathrm{\Pi }(x)=|xx|`$ and $`H`$. This gives for the argument of $`𝐭_0`$ an object $`\mathrm{\Pi }_0(x)=\mathrm{\Omega }^{}\mathrm{\Pi }(x)\mathrm{\Omega }`$ which is not a position projector. Instead, it collects all the states of the free particle that add up to produce the position eigenstate $`|x`$ of the interacting case by the canonical transformation. Much of the difference between the classic and quantum cases is hidden here, in particular the quantum capability to undergo classically forbidden jumps in phase space.
Summarizing, in the interacting case we have a toa operator given by
$$𝐭_x=\underset{s}{}_{\mathrm{}}^+\mathrm{}𝑑tt|txstxs|,$$
(20)
where
$$|txs=(\frac{2H}{m})^{1/4}e^{iHt}\mathrm{\Pi }_s|x$$
(21)
Above we have introduced the projector $`\mathrm{\Pi }_s=𝑑E|Es(+)Es(+)|`$, which is obtained from the $`\mathrm{\Pi }_{s0}`$ of Eq. (92) by the canonical transformation (19). We now have the tools necessary for a physical interpretation in terms of a POVM: Given an arbitrary state $`\psi `$ at $`t=0`$, its time of arrival at a position $`x`$ has to be, according to (20),
$$\psi |𝐭_x|\psi =\frac{1}{P(x)}\underset{s}{}_{\mathrm{}}^+\mathrm{}𝑑tt|txs|\psi |^2,$$
(22)
with the standard interpretation of $`_s|txs|\psi |^2`$ as the (yet unnormalized) probability density that the state $`|\psi `$ arrives at $`x`$ in the time $`t`$. The probability of arriving at $`x`$ at any time is then $`P(x)=𝑑t_s|txs|\psi |^2`$, giving a normalized probability density in times of arrival
$$P(t,x)=\frac{1}{P(x)}\underset{s}{}|txs|\psi |^2$$
(23)
normalization that has been used in (22). Note that in the cases where $`P(x)`$ vanishes this conditional probability is devoid of meaning: If there are no arrivals at all, there are no arrivals in any finite (or infinitesimal) interval of time.
The above equations (22,23) can be given a form that is very useful for computation, while throws some light on the physical meaning of the different quantities involved. By using explicitly (21), one gets
$`P(x)`$ $`=`$ $`{\displaystyle \underset{s}{}}\{{\displaystyle }dE({\displaystyle \frac{2E}{m}})^{1/4}x|Es(+)Es(+)|\psi \}^{}\times `$ (25)
$`\{{\displaystyle 𝑑E^{}(\frac{2E^{}}{m})^{1/4}x|E^{}s(+)E^{}s(+)|\psi }\}{\displaystyle 𝑑te^{i(EE^{})t}}`$
$`=`$ $`2\pi {\displaystyle \underset{s}{}}{\displaystyle 𝑑E(\frac{2E}{m})^{1/2}|x|Es(+)Es(+)|\psi |^2}`$ (26)
Using a similar procedure, one gets for (22)
$`\psi |𝐭_x|\psi `$ $`=`$ $`{\displaystyle \frac{i\pi }{P(x)}}{\displaystyle \underset{s}{}}{\displaystyle }dE({\displaystyle \frac{2E}{m}})^{1/2}\times `$ (28)
$`\{x|Es(+)Es(+)|\psi \}^{}\stackrel{}{{\displaystyle \frac{d}{dE}}}\{x|Es(+)Es(+)|\psi \}`$
$`=`$ $`{\displaystyle \frac{2\pi }{P(x)}}{\displaystyle \underset{s}{}}{\displaystyle 𝑑E(\frac{2E}{m})^{1/2}|x|Es(+)Es(+)|\psi |^2}`$ (29)
$`{\displaystyle \frac{d}{dE}}\{\mathrm{arg}x|Es(+)+\mathrm{arg}Es(+)|\psi \}`$ (30)
## III The entrance into the interaction region
We start here to analyze the theoretical predictions of our formalism. To begin with, we consider the simple case of an initial Gaussian state prepared at $`t=0`$ in a zone where $`V(q)=0`$, and directed towards the interaction region. This wave packet $`\psi `$ of width $`\mathrm{\Delta }q=2\delta `$, is centered at $`q_0<0`$ -well to the left of the onset of the interaction- with mean momentum $`p_0>0`$. In configuration and momentum spaces we have:
$`q|\psi `$ $`=`$ $`({\displaystyle \frac{1}{2\pi \delta ^2}})^{1/4}e^{\delta ^2p_0^2}e^{(\frac{qq_0}{2\delta }i\delta p_0)^2}`$ (31)
$`p|\psi `$ $`=`$ $`({\displaystyle \frac{2\delta ^2}{\pi }})^{1/4}e^{\delta ^2(pp_0)^2ipq_0}`$ (32)
respectively. For appropriate values of $`q_0,p_0`$ and $`\delta `$, such that $`p_0\delta >>1`$ and $`|q_0|>>\delta `$, almost all the packet is initially at the left of the origin and moving with positive momentum towards the right. We use this simplifying assumption (the neglect of the Gaussian’s tails with $`q>0,p<0`$) in our qualitative arguments, and in the intuitive descriptions of the processes that we will develope below. This will be indicated explicitly in the formulas by the use of $``$ instead of $`=`$. However, we shall work with the full expressions (31) wherever necessary in the calculations. For simplicity, we consider that the potential vanishes to the left of the origin. Preparing the state $`\psi `$ as said above with $`\psi (q)0`$ for $`q>0`$, and its Fourier transform $`\stackrel{~}{\psi }(p)0`$ for $`p<0`$, we have $`Es(+)|\psi \delta _{sr}(\frac{m}{2E})^{1/4}\stackrel{~}{\psi }(p)`$, so that
$$txs|\psi \delta _{sr}𝑑Ee^{iEt}x|Er(+)\stackrel{~}{\psi }(p)$$
(33)
valid for the full range of values of $`x`$. Now, the initial state contributes to the time of arrival (30) a quantity $`d/dE\mathrm{arg}Es(+)|\psi mq_0/p`$, the same that in the free case.
### A The quasi-classical case
We start with the simple but illustrative case where the potential departs from 0 for positive $`q`$ with $`V(0)=0`$, and is so smooth that the WKB method is valid. Then, for $`E>V(x)`$ and to lowest order, one can neglect the exponentially small reflection that would vanish classically, getting energy eigenstates of the form
$$x|Er(+)\theta (x)\sqrt{\frac{m}{2\pi p}}e^{ipx}+\theta (x)\sqrt{\frac{m}{2\pi p(x)}}e^{i_0^x𝑑qp(q)}$$
(34)
where $`p(q)=\sqrt{2m(EV(q))}`$. To this order and with a properly normalized wave packet as ours, (26) gives
$$P(x)\theta (x)+\theta (x)P_+(x),P_+(x)=_0^{\mathrm{}}𝑑p\frac{p}{p(x)}|\stackrel{~}{\psi }(p)|^2$$
(35)
so that $`\frac{p}{p(x)}|\stackrel{~}{\psi }(p)|^2`$ is the (unnormalized) probability of arrival at the point $`x`$ with momentum $`p(x)`$. For the probability in times of arrival one gets
$`P(t,x)`$ $``$ $`{\displaystyle \frac{\theta (x)}{2\pi }}\left|{\displaystyle _0^{\mathrm{}}}𝑑pe^{iEt}\stackrel{~}{\psi }(p)\right|^2`$ (37)
$`+{\displaystyle \frac{\theta (x)}{2\pi P_+(x)}}\left|{\displaystyle _0^{\mathrm{}}}𝑑E\sqrt{{\displaystyle \frac{m}{p(x)}}}\stackrel{~}{\psi }(p)e^{i(Et_0^x𝑑qp(q))}\right|^2`$
which is the same as that of free particles for $`x<0`$ as corresponds to this order of approximation in which reflection is neglected, so that there is no information about $`V`$ at the left of the origin. Finally,
$`\psi |𝐭_x|\psi `$ $``$ $`\theta (x){\displaystyle _0^{\mathrm{}}}𝑑p|\stackrel{~}{\psi }(p)|^2{\displaystyle \frac{m}{p}}\{xq_0\}`$ (39)
$`+{\displaystyle \frac{\theta (x)}{P_+(x)}}{\displaystyle _0^{\mathrm{}}}𝑑p{\displaystyle \frac{p}{p(x)}}|\stackrel{~}{\psi }(p)|^2\{{\displaystyle \frac{mq_0}{p}}+m{\displaystyle _0^x}{\displaystyle \frac{dq}{p(q)}}\}`$
Therefore, for negative $`x`$ we recover the toa of the free particle. What the above expression gives for $`x>0`$ is nothing else than the classical time of arrival at $`x`$, Eq. (1), for initial conditions $`(q_0,p)`$ weighted by the probability of these conditions.
### B Step potential and Hartman effect
In general, the approximations that led to (34) do not hold. For instance, reflection has to be taken into account, or $`V`$ is such that the semiclassical approximation is no longer valid, etc. In any case, the particle may eventually reach a point $`q`$ where $`E=V(q)`$. Any further penetration beyond that point is a quantum fenomenon worth to investigate in terms of the toa. We address this question by considering a step potential $`V(q)=\theta (q)V`$ intercepting the path of the wave packet $`\psi `$. We will then analyze the fate of the components of the wave packet with $`p>p_V=\sqrt{2mV}`$ and with $`p<p_V`$. Classically, a particle in the first group will arrive with momentum $`p^{}=\sqrt{|p^2P_V^2|}`$ at the points $`x>0`$, while one in the second group will bounce back at $`q=0`$, without penetrating to the right. In the quantum case, one has for $`x<0`$ a superposition of both, reflection and transmission, regardless of $`p/p_V`$, while for $`x>0`$ one has
$`txs|\psi `$ $``$ $`{\displaystyle \frac{\delta _{sr}}{\sqrt{2\pi }}}{\displaystyle }dE({\displaystyle \frac{m}{2E}})^{1/4}e^{iEt}\times `$ (41)
$`\{\theta (EV)T_>e^{ip^{}x}+\theta (VE)T_<e^{p^{}x}\}\stackrel{~}{\psi }(p)`$
where $`T_>=2p/(p+p^{})`$ and $`T_<=2p/(p+ip^{})`$. Then,
$$P(x)_{p_V}^{\mathrm{}}𝑑p|T_>\stackrel{~}{\psi }(p)|^2+_0^{p_V}𝑑pe^{2p^{}x}|T_<\stackrel{~}{\psi }(p)|^2$$
(42)
is the probability of arrival at $`x`$, while
$`P(t,x)`$ $``$ $`{\displaystyle \frac{1}{2\pi P(x)}}|{\displaystyle }dE({\displaystyle \frac{m}{2E}})^{1/4}e^{iEt}\times `$ (44)
$`\{\theta (EV)T_>e^{ip^{}x}+\theta (VE)T_<e^{p^{}x}\}\stackrel{~}{\psi }(p)|^2`$
gives the probability distribution in toa of the particles that arrive at this point. Finally,
$`\psi |𝐭_x|\psi `$ $``$ $`{\displaystyle \frac{1}{P(x)}}[{\displaystyle _{p_V}^{\mathrm{}}}dp|T_>\stackrel{~}{\psi }(p)|^2\{{\displaystyle \frac{mq_0}{p}}+{\displaystyle \frac{mx}{p^{}}}+{\displaystyle \frac{m}{p}}{\displaystyle \frac{d\mathrm{arg}(T_>)}{dp}}\}`$ (46)
$`+{\displaystyle _0^{p_V}}dpe^{2p^{}x}|T_<\stackrel{~}{\psi }(p)|^2\{{\displaystyle \frac{mq_0}{p}}+{\displaystyle \frac{m}{p}}{\displaystyle \frac{d\mathrm{arg}(T_<)}{dp}}\}]`$
In the case of low potential steps $`p_V<<p_0`$ (c.f. Eq. (31)), where one can neglect the integrals over the interval $`[0,p_V]`$, the probability of arrival reduces to the average of the transmision coefficient $`|T_>|^2`$, which is independent of $`x`$ as corresponds to a transmitted free particle. $`T_>`$ is real in this case, so that $`\psi |𝐭_x|\psi `$ is given by averaging over $`p`$ the time spent to go from $`q_0`$ to 0 at momentum $`p`$ plus the time spent to go from 0 to $`x`$ at momentum $`p^{}`$. The only effect of the step is the reduction of the momentum from $`p`$ to $`p^{}`$.
In the opposite case where $`p_V>>p_0`$, only the integrals over $`[0,p_V]`$ give a sizeable contribution. The probability of arrival vanishes (exponentially) beyond the distance $`\mathrm{\Delta }x=\frac{1}{p^{}}`$ associated through the uncertainty principle to the difference $`\mathrm{\Delta }E`$ between the energy of the step and the energy of the particle. One then expects to detect a relative of this fenomenon in the time of arrival. In fact, the time spent from 0 to $`x`$ is given here through $`\frac{m}{p}\frac{d\mathrm{arg}T_<}{dp}=\frac{m}{pp^{}}`$, which is independent of the distance $`x`$, that is replaced by $`\mathrm{\Delta }x`$. This is a case of the Hartman effect that here arises from the change
$$\frac{p}{p+p^{}}e^{ip^{}x}\frac{p}{p+ip^{}}e^{p^{}x}$$
(47)
in the energy eigenstates as $`p`$ crosses $`p_V`$ from above. In short, the effect is a consequence of the fact that the phase is independent of $`x`$ for $`p<p_V`$.
In the general case one should take into account both contributions to (46). The relative importance of the second contribution in the rhs would depend on $`p_0p_V`$ and will always decrease exponentially with increasing $`x`$. However, a proper analysis of this situation calls for a description of particles better that that provided by first quantization and wave packets. We will defer this question to the next section where we discuss tunneling, the instance where the particle may reappear again beyond some point.
## IV Arrival at the other side
In this section we will study the modification of the times of arrival of quantum particles that traverse potential barriers. Our treatment deepens on the current understanding of the tunneling and dwell times. The literature is full of ad hoc heuristic arguments often disconnected from the standard mathematical and interpretative apparatus of quantum mechanics, whose value is therefore difficult to asses, as is their comparison with experiment. Here, we will follow the standard quantum mechanical treatment of Section 2.
The time of arrival at a point $`x`$ will now be given through a probability amplitude
$$txs|\psi =𝑑E(\frac{2E}{m})^{1/4}e^{iEt}x|Es(+)Es(+)|\psi $$
(48)
We prepare the initial state as usual (as a right mover at the left of the barrier, c.f. above Eq. (33)). We again can approximate $`Es(+)|\psi \delta _{rs}(\frac{m}{2E})^{1/4}\stackrel{~}{\psi }(p)`$. The scattering state of relevance in (48) is given by
$`q|Er(+)`$ $`=`$ $`\sqrt{{\displaystyle \frac{m}{2\pi p}}}\theta (q)(e^{ipq}+R(p)e^{ipq})+\theta (q)\theta (aq)A(q,p)`$ (49)
$`+`$ $`\theta (qa)T(p)e^{ipq}`$ (50)
Expression valid for an arbitrary potential barrier contained in the range $`(0,a)`$, where $`A(q,p)`$ solves the appropriate Schrödinger equation with energy $`E=p^2/2m`$. Also, $`T(p)`$ and $`R(p)`$ are the transmission and reflection coefficients of the barrier. For a barrier of infinite range, the first and third terms in the rhs of (50) should be better understood as asymptotic limits.
Finally, in the case where $`x`$ is at the right of the barrier, the amplitude can be approximately given by
$$txs|\psi \frac{\delta _{sr}}{\sqrt{2\pi }}𝑑E(\frac{m}{2E})^{1/4}e^{i(Etpx)}T(p)\stackrel{~}{\psi }(p)$$
(51)
The normalized probability density in times of arrival at $`x`$ counts all the particles eventually recorded at $`x`$ and only them, that is, the transmitted particles. According to (23) it is given by
$`P(t,x)`$ $`=`$ $`{\displaystyle \frac{1}{P(x)}}{\displaystyle \underset{s}{}}|txs|\psi |^2`$ (52)
$``$ $`{\displaystyle \frac{1}{2\pi P(x)}}\left|{\displaystyle 𝑑E(\frac{m}{2E})^{1/4}e^{i(Etpx)}T(p)\stackrel{~}{\psi }(p)}\right|^2`$ (53)
where we have normalized dividing by $`P(x)`$, the total probability of arrival at $`x`$ in whatever time $`t`$
$$P(x)=\underset{s}{}_{\mathrm{}}^+\mathrm{}𝑑t|txs|\psi |^2_0^+\mathrm{}𝑑p|T(p)\stackrel{~}{\psi }(p)|^2$$
(54)
that is independent of $`x`$ in cases like this, where $`x`$ is beyond the range of the potential. In addition, it approximately simplifies to $`|T(p_0)|^2`$ for narrow wave packets with mean momentum $`p_0`$ not too close (by above or by below) to the barrier momentum $`p_V=\sqrt{2mV}`$. After a straightforward calculation we get for the average time of arrival at the other side of the barrier
$`\psi |𝐭_x|\psi `$ $``$ $`{\displaystyle \frac{i}{2P(x)}}\times `$ (56)
$`{\displaystyle 𝑑E\left[(\frac{m}{2E})^{1/4}e^{ipx}T(p)\stackrel{~}{\psi }(p)\right]^{}}\stackrel{}{{\displaystyle \frac{d}{dE}}}\left[({\displaystyle \frac{m}{2E}})^{1/4}e^{ipx}T(p)\stackrel{~}{\psi }(p)\right]`$
that can be written as
$$\psi |𝐭_x|\psi \frac{1}{P(x)}_0^{\mathrm{}}𝑑p|T(p)\stackrel{~}{\psi }(p)|^2\frac{m}{p}\{xq_0+\frac{d\mathrm{arg}(T(p))}{dp}\}$$
(57)
an expression thas has appeared before in the literature sometimes supported by heuristic arguments alone. It can be understood as the average value of the Wigner time over the transmitted state.
We will illustrate the predictions of the formalism for a simple square barrier of height $`V`$ and width $`a`$. The transmission coefficient is in this case:
$$T(p)=\frac{2pp^{}e^{ipa}}{2pp^{}\mathrm{cos}p^{}ai(p^2+p^2)\mathrm{sin}p^{}a}$$
(58)
where $`p^{}=\sqrt{p^2p_V^2}`$, that is imaginary for $`p`$ below $`p_V`$. Note the contribution $`pa`$ to the phase of $`T(p)`$. This will substract a term $`a`$ to the path length $`xq_0`$ that appears in (57). The barrier has effective zero width or, in other words, it is traversed instantaneously. This is the Hartman effect for barriers. To be precise, the effect is not complete, it is compensated by the other dependences in $`p^{}a`$ present in the phase of $`T(p)`$. In fact, it dissapears for $`(p_V/p)0`$, where all the $`a`$ dependences of the phase cancel out, as was to be expected because the barrier effectively vanishes in this limit. In the opposite case $`(p/p_V)0`$ the effect saturates and there is an advance $`\frac{ma}{p}`$ in the time of arrival of transmitted plane waves, that turns into unexpected results for intermediate barrier momenta.
We present our results for the time of arrival of the transmitted particles in Figs. 1 and 2. We consider the same initial state in both cases, namely the Gaussian wave packet of (31) with $`q_0=30,p_0=2,\delta =10`$, and $`m=1`$, (we always use the natural units of the problem with $`\mathrm{}=1`$). We have computed the time of arrival of the wave packet at $`x=50`$ for an assortment of potential heights and widths, and have chosen the contents of those figures to highlight the most important results.
We show the time of arrival at the other side of a barrier of momentum $`p_V`$ in the range $`a=(1.6,2.6)`$ in Fig. 1. For incident plane waves with momentum $`p_0`$, the barrier would be crossed over for $`p_V<p_0`$, and tunneled through for $`p_V>p_0`$. Some retardation would be expected in the first case, just because the travel over the barrier would be slowlier than the free travel. This is clearly seen at the left of $`p_0`$ in the figure. Classically, the delay would grow from zero (time $`t_0`$) to infinity as $`p_V`$ grows from $`0`$ to $`p_0`$. The quantum behaviour is similar, with the oscillations of the phase time swept away by the average that remains finite. To the right of $`p_0`$, there is a dramatic difference between the Wigner result, that inmediately sticks to the Hartman prediction $`t_H`$, and the wave packet result, for which the time continues to increase up to a certain barrier height and then, suddenly, drops to $`t_H`$. This strange behaviour can be explained in the following manner: Not being monoenergetic, the wave packet has momentum components above and below $`p_V`$. The first of these cross above the barrier, get retarded, and are responsible for the high time value for $`p_V`$ just to the right of $`p_0`$. However, as the barrier continues to grow, they become an ever lesser part of the packet. The other parts of the packet (the components with momentum $`p<p_V`$) tunnel through the barrier, and experience the Hartman advance. They would arrive at $`x`$ in a time $`t_H`$. Their relative importance in the wave packet increases steadily as $`p_V`$ continues to grow and, eventually, they overcome the retarded components and the process becomes pure tunneling. Then, the time of arrival drops to $`t_H`$. We have numerically checked this behaviour, that we have analyzed for several values of the barrier width in the range (2,30). All the results are similar: Monotonic grow of the time from $`p_V=0`$ (where $`t=t_0`$), up to $`p_V2.5`$, where $`t`$ drops suddenly to $`t_H`$. The general trend is a slow increase in the value of the barrier momentum $`p_V`$ at which the drop takes place, that shifts from about 2.2 to 2.7 as $`a`$ changes from 10 to 30. The maximum value of the time of arrival $`t_x`$ that is obtained just before the drop also increases; it is around 95 for $`a=10`$ and around 450 for $`a=20`$.
We show in Fig. 2 the average time of arrival and the Wigner (phase) time as a function of the barrier width $`a`$ in the range $`a=(0,15)`$. We display the predictions for different barrier heights $`p_V=0,1.6,1.8,2.2,2.3,2.4`$ and $`2.6`$. For the free case ($`p_V=0`$, or $`a=0`$) all the results converge to $`t_0=40`$. We now disscuss the solid lines $`t_x`$. The oscillatory curves above $`t_0`$ correspond to $`p_V<p_0`$. The get steeper as their momenta approach $`p_0`$ from below. The curves that stand partially below $`t_0`$ correspond to $`p_V>p_0`$ (tunneling). They share a similar behaviour: As the barrier width grows from $`a=0`$ the time of arrival decreases, practically saturating the Hartman time $`t(a)t_H(a)=t_0\frac{ma}{p_0}`$. Then suddenly, at a certain width (that increases with $`p_V`$), the average time jumps dramatically to values that correspond to a long retardation. Note that the jump for $`p_V=2.6`$ lies outside the range of the figure. This behaviour is complementary to that shown in Fig. 1. Here, for $`p_V>p_0`$ and moderate $`a`$, tunneling is the dominant phenomenon and the time average tends to reproduce $`t_H`$. However, as the barrier gets wider, tunneling gets more and more depressed. In comparison, the intensity of the retarded components that pass over the barrier is basically independent of $`a`$. They get relatively more and more important and, eventually overcome tunneling, giving rise to the observed transition. In practice, for wide enough barriers, the probability of tunneling vanishes, and the other side can be reached only by the very improbable and very slow travel over the barrier. This behaviour has been noticed independently in , and explained in the same way. In addition, we have the tools to check these explanations. In particular, the first product of our formalism is $`P(t,x)`$, the probability distribution in times of arrival at $`x`$. Our numerical analysis for $`x=50`$ and the different $`p_V`$’s and $`a`$’s that we are discussing here show similar almost Gaussian shapes for these distributions, as correspond to the initial wave packets chosen, and similar widths for these $`P(t,x)`$, whose maxima are placed close to the corresponding mean values $`t_x`$. As expected, the probabilities get numerically smaller as the corresponding events become more and more unlikely. In short, these distributions give the best support for the validity of the explanation offered here for this striking behaviour, that can be understood only after weighting the obtained time of arrival with the relative probability of the actual event to which it corresponds.
## V Quantum reflections
Having analyzed the modifications introduced by the transmission phenomena in the time of arrival at the other side of potential barriers, we turn to the case of reflection. We divide the analysis into the two seemingly different cases in which there is classical reflection, and in which it is absent. The first case is characterized by the presence of at least one turning point in the path of the particle. The second one, by the absence of any of them. Quantum mechanically there could be some transmission in the first case, and some reflection in the second one. Accordingly, we separate the disscussion that follows into the two main disjoint cases that cover all the possibilities. These are the case where the potential energy grows to infinity somewhere (total reflection), and the case where it is bounded everywhere (with partial reflection and transmission).
### A The case of total reflection
The potential energy could grow unbound, thus reflecting any conceivable incoming state. We consider here a monotonic potential energy that vanishes for $`q\mathrm{}`$ and goes to infinity for $`q\mathrm{}`$ so that $`lim_{q+\mathrm{}}q|E=0`$. This removes the degeneracy of the energy eigenstates. As no state may arrive from the right, $`q|El(+)=0`$. The eigenstates $`|E`$ will contain the same amount of positive and negative momenta, so that their asymptotic form normalized to one traveling particle per unit time is $`lim_q\mathrm{}q|E=\frac{1}{\sqrt{2\pi }}(\frac{m}{2E})^{1/4}\mathrm{cos}(pq+\delta (E))`$, where $`\delta (E)`$ is the phase shift. This also fixes completely the eigenstates for finite values of $`q`$.
The time of arrival at an arbitrary point $`x`$ is now
$`\psi |𝐭_x|\psi =`$ $`{\displaystyle \frac{2\pi }{P(x)}}{\displaystyle }dE\sqrt{{\displaystyle \frac{2E}{m}}}|x|EE|\psi |^2\times `$ (59)
$`{\displaystyle \frac{d}{dE}}\{\mathrm{arg}x|E+\mathrm{arg}E|\psi \}`$ (60)
which is the average of a quantity independent of $`x`$! This comes about because in the present situation the reflection coefficient $`R=\mathrm{exp}(2i\delta )`$ is unimodular. Then, the net current density vanishes, so that $`\mathrm{arg}x|E`$ is independent of $`x`$. This is the quantum version of the classical result that the sum of the times of arrival at $`x`$ of the incoming and returning particles is twice the toa at the turning point and so, independent of $`x`$. Obviously this ceases when $`|R|`$ becomes smaller than 1 (so that the net current density is finite), something that is possible only when $`V`$ is finite everywhere. Even then, the classical result is recovered from the quantum case in the limit $`(E/V)<<1`$ where $`|R|1`$.
The individual times of arrival of the incoming and the returning particles can be obtained straightforwardly by writing the enegy eigenstates as
$$q|E=\frac{1}{\sqrt{2\pi }}(\frac{m}{2E})^{1/4}M(q,E)\mathrm{cos}\varphi (q,E)$$
(61)
where $`M`$ is a real function with $`lim_q\mathrm{}M(q,E)=1`$, that vanishes faster than an exponential for $`q+\mathrm{}`$ to satisfy the asymptotic form of the Schrödinger equation. The state is thus written as the superposition at each point of an incoming and a reflected wave with equal amplitudes, so that the net current vanishes everywhere. The phase $`\varphi `$ is fixed by $`lim_q\mathrm{}\varphi (q,E)=pq+\delta (E)`$ to match the asymptotic form of the eigenstate disscussed above. Its derivative gives the two opposite velocity fields $`v_\pm (q,E)=\pm \frac{d\varphi (q,E)}{mdq}`$ interfering at $`q`$. We recall that this exact expression is valid for all the potentials of the form we are considering here. The probability of ever arriving at $`x`$ and the toa can be given by straightforward application of (26) and (30) by
$`P(x)`$ $`=`$ $`{\displaystyle 𝑑EM^2(x,E)\mathrm{cos}^2\varphi (x,E)|E|\psi |^2}`$ (62)
$`\psi |𝐭_x|\psi `$ $`=`$ $`{\displaystyle \frac{1}{2P(x)}}{\displaystyle }dEM^2(x,E)\mathrm{cos}^2\varphi (x,E)\times `$ (64)
$`|E|\psi |^2\left[t_i(x,E)+t_r(x,E)\right]`$
which is the weighted average over energies of the times of arrival of the incoming and the reflected waves:
$`t_i(x,E)`$ $`=`$ $`{\displaystyle \frac{d}{dE}}\{\varphi (x,E)+\mathrm{arg}E|\psi \}`$ (65)
$`t_r(x,E)`$ $`=`$ $`{\displaystyle \frac{d}{dE}}\{\varphi (x,E)+\mathrm{arg}E|\psi \}`$ (66)
whose sum is explicitly $`x`$ independent.
To illustrate these results we consider now the case of a potential that vanishes at the left of the origin and is linear at the right, i.e. $`V(q)=\theta (q)fq`$, where $`f`$ is the force exerted on the particle. This could be a model for a (charged) particle in a constant electric field, or in the gravity field of the Earth. In this case one gets $`M`$ and $`\varphi `$ in terms of the Airy function Ai and its derivativeAi’.
$$M(q,E)=\{\begin{array}{cc}\hfill 1& \hfill \text{for}q0\\ \hfill \sqrt{\frac{\text{Ai}[z]^2+(\frac{k_f}{p})^2\text{Ai’}[z]^2}{\text{Ai}[z_0]^2+(\frac{k_f}{p})^2\text{Ai’}[z_0]^2}}& \hfill \text{for}q>0\end{array}$$
(67)
where $`z=k_fqp^2/k_f^2,z_0=p^2/k_f^2`$ with $`k_f=(2mf)^{1/3}`$. For the phase one has
$$\varphi (q,E)=\{\begin{array}{cc}\hfill \mathrm{arctan}(\frac{k_f\text{Ai’}[z_0]}{p\text{Ai}[z_0]})& \hfill \text{ for}q0\\ \hfill \mathrm{arctan}(\frac{k_f\text{Ai’}[z]}{p\text{Ai}[z]})& \hfill \text{ for}q>0\end{array}$$
(68)
so the phase shift is given simply by $`\delta (E)=\varphi (0,E)`$.
We present in Figures 3 and 4 our results for the the case of a force of nominal value $`f=100`$, being the parameters of the initial Gaussian wave packet (31) $`q_0=2,p_0=10,\delta =1`$ and $`m=1`$. For the normalized probability distributions in times of arrival (23) we get pairs of peaks of equal heights - as correspond to total reflection - that tend to merge into one as the detector is displaced towards the classical turning point. This behaviour of the peaks is also observed for the averaged times of arrival, that follow the classical times. The small deviations from the parabolic form are negligible in comparison with the widths of the distributions shown in Fig. 3. We have explored numerically the details that change uninterestingly according to the values of $`f,p_0,\delta `$ etc. so, we do not show them here. The general picture is always the same: at the far left ($`|q_0|>>E/f`$) the potential acts as an infinite height wall. The only sizeable consequences of the actual strenght of the force are felt at positions between the origin and the turning point, where they resemble the classical effects. Part of this comes from the fact that here position and energy combine into only a variable $`qE/f`$. But the resemblance arises because total reflection is always present here, quantum as well as classically. This will be more clear in the next section where we consider partial reflection that lacks of classical analog.
### B Partial reflections
In classical mechanics a potential interaction energy speeds up or slows down the particles according to the local value of the force $`F(q)=\frac{V(q)}{q}`$. Accelerated or decelerated, the particles continue to move along the same path without reversing the direction. Only when one of them intercepts a turning point (i.e. a point $`q`$ where $`E=V(q)`$) the particle bounces back or, in other words, is reflected with probability $`P_R=1`$. In the absence of these points, the particle is always transmitted with probability $`P_T=1`$. Thus, most of the time $`P_T=1,P_R=0`$. Only at the turning points $`P_T=0,P_R=1`$.
Quantum dynamics offers a very different perspective of the motion of the particles. The Schrödinger equation implies that at every point where the potential energy is finite, the particle is partially transmitted and partially reflected, that is $`0P_T1,0P_R1`$, with $`P_T+P_R=1`$. The case of total reflection analyzed in the previous section is one of close correspondence between the classical and the quantum results, as we shown there. Interesting departures from the classical behaviour arise when there is no classical reflection. We will analyze this case here.
To fix ideas, we consider a well behaved potential energy $`V(q)0`$ finite $`q`$, that vanishes at the spatial infinity faster than $`q^1`$. In these conditions the energy eigenstates can be written everywhere as a well defined superposition of transmitted $`\mathrm{\Phi }_{tr}(q,E)`$ and reflected $`\mathrm{\Phi }_{ref}(q,E)`$ waves, characterized by the positive or negative value of their currents: $`\frac{i}{2m}(\mathrm{\Phi }_{tr}^{}\stackrel{}{\frac{d}{dq}}\mathrm{\Phi }_{tr})0`$, and $`\frac{i}{2m}(\mathrm{\Phi }_{ref}^{}\stackrel{}{\frac{d}{dq}}\mathrm{\Phi }_{ref})0`$, with different amplitudes $`|\mathrm{\Phi }_{tr}||\mathrm{\Phi }_{ref}|`$ as corresponds to this case of partial reflection. The eigenstates of interest can be written as
$$q|Er=\sqrt{\frac{m}{2\pi p}}\{\mathrm{\Phi }_{tr}(q,E)+\mathrm{\Phi }_{ref}(q,E)\}$$
(69)
These waves are univocally determined by their asymptotic conditions, namely:
$`\underset{q\mathrm{}}{lim}\mathrm{\Phi }_{tr}(q,E)=`$ $`e^{ipq},`$ $`\underset{q+\mathrm{}}{lim}\mathrm{\Phi }_{tr}(q,E)=T(E)e^{ipq}`$ (70)
$`\underset{q\mathrm{}}{lim}\mathrm{\Phi }_{ref}(q,E)=`$ $`R(E)e^{i(pq+2\delta (E))},`$ $`\underset{q+\mathrm{}}{lim}\mathrm{\Phi }_{ref}(q,E)=0`$ (71)
as is the case for an incoming rightmover (69). The results of the previous section are recovered in the limit where $`T(E)0`$ which is the case only if the potential energy grows to infinity somewhere.
If we prepare our initial Gaussian state $`\psi (q)`$ at a point $`q=q_0`$ where the potential energy is smooth enough, and keep the initial momentum $`p_0>0`$ large enough to consider $`\stackrel{~}{\psi }(p)0`$ for $`p<0`$, we can use the approximations
$$Es|\psi \delta _{rs}\sqrt{\frac{m}{p}}\mathrm{\Phi }_{tr}^{}(q_0,E)|\stackrel{~}{\psi }(p)|\delta _{rs}\sqrt{\frac{m}{p}}e^{ipq_0}|\stackrel{~}{\psi }(p)|$$
(72)
We have used the second of these already in Eq.(33). It is valid when $`V(q)0`$ for $`q`$ in the $`q_0`$ neighbourhood where $`\psi (q)`$ is sizeable. We assume this is the case in what follows.
One of the deepest consequences of the superposition of transmitted and reflected components that makes up the eigenstate (69) is that it leads to the inescapable presence of interferences. In fact, the probability of presence at a point $`q`$, and other quantities depending on it, contain the sum $`|q|Er|^2|\mathrm{\Phi }_{tr}|^2+|\mathrm{\Phi }_{ref}|^2+2\mathrm{}(\mathrm{\Phi }_{tr}\mathrm{\Phi }_{ref}^{})`$, whose last term is the interference term. One could say that, everywhere in its motion through the interaction region, the quantum particle will be found in an evolving entangled state of transmitted and reflected components. This can be traced back mathematically to the continuity of the solutions of the Schrödinger equation and of their first derivatives, and to the associated Wronskian theorem. Physically, this may introduce all sorts of interpretative difficulties in the analysis of particle motion.
Summarizing, interferences pervade the realm of quantum motion. They will show up in almost every quantum mechanical situation. Our analysis of the time of arrival is not an exception. We have avoided refering to them till now by focusing on very specific cases. These were: The choice in Sect. 3.1. of a very smooth potential analyzable semiclassically by the WKB method, that neglects reflection. The analysis in Sect. 3.3. of the time of arrival at points located at the other side of the barrier, where $`\mathrm{\Phi }_{ref}=0`$ so that any interference with the transmitted wave vanishes. Finally, the analysis made in the previous section, where we just ignored the effects due to the overlap of incoming and reflected waves in $`P(t,x)`$, and the lack of a clear cut separation between $`t_i`$ and $`t_r`$ in the presence of interferences. To be precise, we dealt with reflection without paying the due attention to these subtleties. We repair the ommission here.
The amplitude in time of arrival at a position $`x`$ within the interaction range can be given by using (69) and (72) in (48)
$`txs|\psi `$ $`=`$ $`\{A_{tr}(t,x)+A_{ref}(t,x)\}`$ (73)
$``$ $`{\displaystyle \frac{\delta _{sr}}{\sqrt{2\pi }}}{\displaystyle _0^{\mathrm{}}}𝑑p\sqrt{{\displaystyle \frac{p}{m}}}|\stackrel{~}{\psi }(p)|e^{i(Et+pq_0)}\{\mathrm{\Phi }_{tr}(x,E)+\mathrm{\Phi }_{ref}(x,E)\}`$ (74)
This gives for the probability of ever arriving at $`x`$ Eq. (26) the sum of three terms: The two separated probabilities $`P_{tr},P_{ref}`$ of arriving with positive or with negative current density, and a quantum interference term, whose presence deprives the previous two of direct physical meaning. We thus get $`P(x)=P_{tr}(x)+P_{ref}(x)+I(x)`$ with
$$P_{_{\genfrac{}{}{0pt}{}{tr}{ref}}}(x)=𝑑t|A_{_{\genfrac{}{}{0pt}{}{tr}{ref}}}(t,x)|^2𝑑p|\stackrel{~}{\psi }(p)|^2|\mathrm{\Phi }_{_{\genfrac{}{}{0pt}{}{tr}{ref}}}(x,E)|^2$$
(75)
and an interference term
$`I(x)`$ $`=`$ $`2{\displaystyle 𝑑t\mathrm{}\{A_{tr}(t,x)A_{ref}^{}(t,x)\}}`$ (76)
$``$ $`2{\displaystyle 𝑑p|\stackrel{~}{\psi }(p)|^2\mathrm{}\{e^{ipq_0}\mathrm{\Phi }_{tr}(x,E)\mathrm{\Phi }_{ref}(x,E)^{}\}}`$ (77)
The above quantities depend on the probabilities of transmission or reflection from the initial position $`q_0`$ to the actual value $`x`$. Consider a bounded potential barrier of finite range, but otherwise arbitrary. Behind the barrier $`P_{ref}`$ vanishes, while $`P_{tr}`$ is given by (54) with a value independent of $`x`$, but strongly dependent of $`p_0,\delta `$ and of the barrier’s height and width. For $`x`$ at the left of the barrier $`\mathrm{\Phi }_{tr}=e^{ipx}`$ (what we are denoting as transmission is here incidence), but $`\mathrm{\Phi }_{ref}=R(E)e^{i(px+2\delta (E))},`$ and only when there is no reflection (no barrier) the intereferences dissappear. For the total reflection case of the previous section, we get $`P_{tr}=P_{ref}`$, while the interference term gives rise to the term $`\mathrm{cos}\{2(px+\delta (E))\}`$ that builds up the factor $`\mathrm{cos}^2\varphi `$ that appears in (62) and (64). However it does not prevent the definition of the quantities (65) and(66) that allowed to split the toa (64) into two positive contributions interpretable as the independent $`t_x`$ of an incoming packet and a reflected one (Fig. 4).
For finite barriers reflection is always present with an energy dependent coefficient $`R(E)<1`$; it is less probable than incidence, and tends to vanish as the barrier does. In Fig. 5 we give the probability distributions of toa $`P(t,x)`$ at a point $`x`$, whose bumps indicate, as in Fig. 3, the arrival of incident and reflected parts of the time evolved initial wave packet. This is the Gaussian one with $`m=1,p_0=2,\delta =10`$, placed at $`q_0=150`$. The arrival position is at $`x=100`$, far from $`q_0`$ to avoid interferences. The two upper figures are for a barrier of width $`a=4`$. At the left is the case where $`p_V=2.2`$, and at the right that with $`p_V=1.9`$. In both cases there is an incidence bump centered at $`t=m(xq_0)/p_0=25`$, and a structure to its right corresponding to reflection. For $`p_V=2.2`$, and for all the cases of total classical reflection ($`p_V>p_0`$), the latter is a Gaussian-like bump shifted from the classical value at $`t=m(xq_0)/p_0=125`$ by an amount $`\frac{m}{p}(\frac{d\varphi }{dp})`$. However, for $`p_V=1.9`$ (in general for $`p_V<p_0`$), the reflected distribution has a multi-bump shape difficult to understand in terms of the phase time or of any other approximation. In particular, neither the number of peaks, nor their positions heights and widths can be approximated by straight stationary phase methods. Two illustrative cases of these shapes are shown in some detail in the two examples of the lower part that correspond to $`p_V=1.9`$ and two close widths $`a=4`$ and $`a=6`$.
## VI Conclusions
We have worked out a formalism for obtaining the time of arrival at a space point of particles that move through interacting media. Our construction follows a circuitous path: we desist from first computing the classical toa of the problem, and then quantizing it, a procedure that leads to a dead end. Instead, we start from the quantum toa of the free moving particle, and then transform it canonically to the interacting case. This is achieved by the use of the appropriate Möller operator that implements the quantum version of the Jacobi-Lie canonical transformation to free translation coordinates in phase space. In the classical case we have the transformation of Eq. (82) whose quantum counterpart is
$$\{,H\}\stackrel{\mathrm{\Omega }^{}}{}\{_0,H_0\}$$
(78)
where $`_0`$ and $``$ are the Hilbert spaces of the free and interacting particles, and $`H_0,H`$ the respective Hamiltonian operators in these spaces. For simplicity, we have only addressed explicitly cases in which the transformations are unitary, which is the case when $`\sigma (H)=\sigma (H_0)`$. More general situations that require of isometric transformations, deserve a separate treatment by their physical relevance.
What we obtained here is a quantum formalism for the toa in terms of a POVM given by
$$P(t_1,t_2;x)=\underset{s=r,l}{}_{t_1}^{t_2}𝑑t|txstxs|$$
(79)
which measures the probability of arrival at $`x`$ during the time interval $`(t_1,t_2)`$. The normalized probability distribution $`P(t,x)`$ was given in the Eq. (23) of Sect. II.B. Our results are thus within the standard formalism of quantum mechanics and can be interpreted in the standard way. There is nothing special that singles out our theoretical predictions as unsuitable for comparison with the experimental results. On the contrary, our formalism predicts the result of actual experiments in the form of numeric values and statistics for the recorded events.
After the definition and theoretical analysis of Sect. II. we have performed explicit and complete calculations for the cases of an unbounded linear potential, of the step potential and of the square barrier. Our analysis of the quasi-classical case shows that in this limit the toa is simply given by the average of the classical time of Eq. (1) over the quasi-classical wave function. In the case of reflection, and for the arrival point placed between the initial position of the wave packet and the turning point ($`x<0`$), the probability distribution $`P(t,x)`$ is governed by the quantum superposition of the incident $`(A_{tr})`$ and the reflected $`(A_{ref})`$ wave packets. In the case of total reflection, where both are equally probable $`P_{tr}(x)=P_{ref}(x)`$, we have obtained separate positive $`t_x`$ even when both amplitudes overlap. These were interpreted as the toa’s of the incident and reflected particles, and compared successfully with the classical prediction. For partial reflection, $`P_{ref}(x)<P_{tr}(x)`$ non overlapping amplitudes are necessary to to get separate average values for these times. This problem is shared with the position and other operators. It is not a defect of the formalism, but an effect of the interferences. Fortunately enough, our formalism provides us with the probability distribution $`P(t,x)`$ whose diverse humpy-bumpy shapes (Figs. 3 and 5) give the most complete information of the posible experimental outcomes.
In the course of our numerical analysis we have detected that the phase time $`\tau _\varphi `$ not always gives a good approximation to the most probable time of arrival. It provides a first estimate of the time spent in the transmission or reflection, after substracting the time of free flight. For transmitted wave packets we have reobtained the advancement (i.e. a decrease in the toa) in the case of pure quantum tunneling. This phenomenon, predicted by Hartman long time ago , has been experimentally evinced by the two photon experiments at Berkeley and the tunneling of optical pulses at Wien . However, our formalism predicts a striking departure from the Hartman bound that we explain in detail in Sect. IV. Our results for square barriers neatly show the expected advancement roughly proportional to the width $`\mathrm{\Delta }t=ma/p`$ (Figs. 1 and 2). However, whatever the mean energy $`(E<V)`$ of the incident wave packet, there is always a width $`a_0`$ such that for $`a>a_0`$ the (very retarded) components of the packet that stand above the barrier dominate over the (probabilistically very depressed) tunneled ones, giving an overall effective strong retardation. In other words, when the barrier is wide enough, its width dominates over the Hartman lenght $`\mathrm{\Delta }x`$ disscussed above Eq. (47), that has a purely quantum origin. This restores the classical expectation of no tunneling and very long delays.
We have also found other unanticipated phenomenon for purely quantum reflection: the multiple bump structure that appears when $`p_V<p_0`$. We have shown in Fig. 5 this structure, that in some sense is a counterpart of the interference pattern that appears in multiple reflection of stationary waves. We think that this feature, even if less spectacular than the superluminal tunneling of photons, deserves experimental confirmation. An appropriate modification of the two photon experiments could serve for this purpose. It would require to place a quantum mirror in the path of one of the entangled photons, and check for the presence (or absence) of the multiple dip structure in the number of coincidence counts predicted by the formalism.
All the examples above show that our construction of a quantum toa operator suitable for the presence of interactions allows the exploration of many physical details in relevant situations. Its extension to higher dimensional cases poses no conceptual difficulties and opens the possibility of treating new questions. Of great theoretical and experimental interest will be the extension of this formalism to the cases in which the Hamiltonian has bound states, where isometric (instead of simply unitary) transformations will be requiered.
## Appendix A
In the modern literature , a classical Hamiltonian system with $`n`$ degrees of freedom is called completely integrable ( a là liouville) when it satisfies the conditions $`a`$ and $`b`$ below:
* There are $`n`$ compatible conservation laws $`\mathrm{\Phi }_i(q_1,\mathrm{},q_n,p_1\mathrm{},p_n;t)=C_i`$, $`i=1,\mathrm{},n`$, that is:
+ $`\dot{\mathrm{\Phi }}_i=\{\mathrm{\Phi }_i,H\}+\frac{\mathrm{\Phi }_i}{t}=0,i=1,\mathrm{},n.`$
+ $`\{\mathrm{\Phi }_i,\mathrm{\Phi }_j\}=0,i,j=1,\mathrm{},n.`$
* The conservation laws define $`n`$ isolating integrals that can be written as:
+ $`\mathrm{\Phi }_i=C_ip_i=\varphi _i(q_1,\mathrm{},q_n,C_1,\mathrm{},C_n;t),i=1,\mathrm{},n.`$
+ $`\frac{\varphi _i}{q_j}=\frac{\varphi _j}{q_i}i,j=1,\mathrm{},n.`$
In these conditions, Hamilton equations define an integrable flow, that is, a system of holonomic coordinates $`(q(t),p(t))`$ in phase space for each instant of time:
$`q_i(t)`$ $`=q_i(q_0,p_0;t),`$ $`i=1,\mathrm{},n.`$ (80)
$`p_i(t)`$ $`=p_i(q_0,p_0;t),`$ $`i=1,\mathrm{},n.`$ (81)
In other words, given a set of initial conditions $`(q_0,p_0)`$ of the system, at each instant of time $`t`$ the system arrives at a point $`(q(t),p(t))`$ in phase space. Conversely, these points define the corresponding times of arrival. In this case, time meets the requirements to qualify as a derived variable in phase space.
As Lie pointed out, for any arbitrary time there is a special choice of coordinates in phase space that mathematically eliminates the effects of interactions from these integrable flows, (the new positions are ignorable coordinates). More simply, integrable systems are canonically equivalent to a set of translations (or circular motions) at constant speed. It is customary to denote the variables that determine these translations as action-angle variables, which strictly is appropriate only in the case of periodic systems, where the (closed) flow lines are topologically equivalent to circles.
For integrable flows, there is a canonical transformation (the Jacobi-Lie transformation)
$$\{q,p;H(q,p)\}\stackrel{W(q,P)}{}\{Q,P;\overline{H}(Q,P)\}$$
(82)
with $`H(q,p)=\overline{H}(Q,P)`$, that gives the free translation coordinates $`P(t)=P`$, and $`Q(t)=\frac{P}{m}t+Q`$ of the translation flow with $`\overline{H}(Q,P)=\frac{P^2}{2m}`$, in terms of the coordinates and momenta $`(q(t),p(t))`$ of the actual flow with $`H(q,p)=\frac{p^2}{2m}+V(q)`$. This transformation is of the form $`W(q,P)`$, that is, a function of the old coordinates and the new momenta, so that
$$Q=\frac{W}{P},p=\frac{W}{q}$$
(83)
Finally, $`W`$ can be obtained explicitly as a complete integral of the Hamilton-Jacobi equation:
$$H(q,\frac{W}{q})=\frac{P^2}{2m}$$
(84)
Now, the canonical relation among the new and the old variables is:
$`P`$ $`=`$ $`\text{sign}(p)\sqrt{2mH(q,p)}`$ (85)
$`Q`$ $`=`$ $`{\displaystyle _0^q}{\displaystyle \frac{dq^{}}{\sqrt{1\frac{V(q^{})}{H(q,p)}}}}+Q_0`$ (86)
where $`Q_0`$ is a constant. As a byproduct, time gets defined in equivalent manner in terms of the old variables, or of the new ones. If the particle arrives at $`q(t)=x`$ in the instant $`t(x)=t`$, then:
$$t(x)=\frac{m}{P}(XQ)=\text{sign}(p)_q^x\frac{mdq^{}}{\sqrt{2m(H(q,p)V(q^{}))}}$$
(87)
where $`X=W(x,P)/P`$ (obviously, $`X=Q(t)`$ by construction). This duality, devoid of practical interest in the classical domain, is at the foundations of the quantum method developed in this paper. Finally, note that for simplicity we have specialized the notation to the case of autonomous Hamiltonian systems with only one degree of freedom, all of them trivially integrable ($`H(q,p)=E`$ being the needed conserved quantity).
## Appendix B
For free particles Eq. (1) gives $`t_{x0}(q,p)=m(xq)/p`$ that, in spite of its simplicity, presents some problems for quantization whose solution we outline here. First of all, it requires symmetrization:
$$𝐭_{x0}(𝐪,𝐩)=m(\frac{x}{𝐩}\frac{1}{2}\{𝐪,\frac{1}{𝐩}\}_+)=e^{i𝐩x}\sqrt{\frac{m}{𝐩}}𝐪\sqrt{\frac{m}{𝐩}}e^{i𝐩x}$$
(88)
As is well known, the eigenstates $`|txs0`$ of this operator in the momentum representation can be given as ($`\mathrm{}=1`$)
$$p|txs0=\theta (sp)\sqrt{\frac{|p|}{m}}\mathrm{exp}(i\frac{p^2}{2m}t)p|x$$
(89)
where we use $`s=r`$ for right-movers ($`p>0`$), and $`s=l`$ for left movers ($`p<0`$.) The label $`0`$ stands for free case. Finally, the argument $`sp`$ of the step function that appears in the momentum representation is $`+p`$ for $`s=r`$, and $`p`$ for $`s=l`$. The degeneracy of the energy with respect to the sign of the moment is explicitly shown by means of the label $`s`$ in the energy representation, where
$$Es^{}0|txs0=\delta _{s^{}s}(\frac{2E}{m})^{1/4}e^{iEt}Es0|x$$
(90)
Summarizing, there is a time (of arrival at $`x`$) representation spanned by the eigenstates
$$|txs0=(\frac{2H_0}{m})^{1/4}e^{iH_0t}\mathrm{\Pi }_{s0}|x$$
(91)
where $`\mathrm{\Pi }_{s0}`$ projects on the subspace of right-movers ($`s=r`$), or left-movers ($`s=l`$), i.e.
$$\mathrm{\Pi }_{s0}=_0^{\mathrm{}}𝑑E|Es0Es0|$$
(92)
These time eigenstates are not orthogonal, which in the past gave rise to serious doubts about their physical meaning. The origin of this problem can be traced back to the fact that (88) is not self-adjoint, that is $`\phi |𝐭_{x0}\psi 𝐭_{x0}\phi |\psi `$. This was proved by Pauli long time ago and is due to the lower bound on the energy spectrum. The problem emerges as soon as one attempts integration by parts in the energy representation. Ref. is a recent illuminating review of these and other related questions.
The measurement problem posed by this not self-adjoint toa operator can be solved by interpreting it in terms of a Positive Operator Valued Measure (POVM), that only requires the hermiticity of $`𝐭_{x0}`$ (i.e. $`𝐭_{x0}=(𝐭_{x0})_{}^{}{}_{}{}^{}`$). Here, instead of a Projector Valued spectral decomposition of the identity operator, one has the POVM
$`P_0(\mathrm{\Pi }(x);t_1,t_2)`$ $`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle _1^2}𝑑t|txs0txs0|`$ (93)
$`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle _1^2}𝑑t({\displaystyle \frac{2H_0}{m}})^{1/4}e^{iH_0t}\mathrm{\Pi }_{s0}\mathrm{\Pi }(x)\mathrm{\Pi }_{s0}e^{iH_0t}({\displaystyle \frac{2H_0}{m}})^{1/4}`$ (94)
where $`\mathrm{\Pi }(x)=|xx|`$ is the projector on $`x`$. Here, $`P_0(1,2)^2P_0(1,2)`$ because $`|txs0txs0|`$ is not a projector, as the states are not orthogonal, but where the limit as $`t\mathrm{}`$ of $`P_0(t,+t)`$ is the identity. The attained time operator is no longer sharp, but is well suited for measurement. This solution has been implemented in , and extensively analyzed in refs. and in the review . In this POVM formulation the toa is given by the spectral decomposition
$$𝐭_0(H_0,\mathrm{\Pi }(x))=_{\mathrm{}}^+\mathrm{}𝑑tt(\frac{2H_0}{m})^{1/4}e^{iH_0t}𝒫_0(x)e^{iH_0t}(\frac{2H_0}{m})^{1/4}$$
(95)
where $`𝒫_0(x)=_s\mathrm{\Pi }_{s0}\mathrm{\Pi }(x)\mathrm{\Pi }_{s0}`$, which is not a projector. |
warning/0002/cond-mat0002248.html | ar5iv | text | # Folding transitions of the square–diagonal two–dimensional latticePACS numbers: 05.50.+q (Ising problems); 64.60.-i (General studies of phase transitions); 82.65.Dp (Thermodynamics of surfaces and interfaces).
## Abstract
The phase diagram of a vertex model introduced by P. Di Francesco (Nucl. Phys. B 525, 507 1998) representing the configurations of a square lattice which can fold with different bending energies along the main axes and the diagonals has been studied by Cluster Variation Method. A very rich structure with partially and completely folded phases, different disordered phases and a flat phase is found. The crumpling transition between a disordered and the flat phase is first–order. The CVM results are confimed by the analysis of the ground states and of the two limits where the model reduces to an Ising model.
1. Introduction
Polymerized membranes are two–dimensional networks of molecules with fixed connectivity . Their configurational properties and the existence of a crumpling transition between a folded and a flat phase are relevant for the behavior of some biological systems . Models of polymerized membranes can be defined on regular lattices. The constraint of fixed connectivity gives rise to the definition of complicate vertex models which offer the advantage that explicit analytical calculations can be performed . In these models the bonds between the vertices of the network have a fixed length and also, as a further simplification, self–avoidance is not considered. In this way the only degrees of freedom of the network are related to its possible states of folding and each state can be weighted by a Boltzmann factor depending on the relative angle between adjacent plaquettes in the network.
The model studied in this paper has been introduced by Di Francesco and concerns the folding properties of a square lattice which can be folded along the main axes and along the diagonals. This model follows previous studies of the folding properties of a triangular network embedded in $`d`$–dimensional lattices. In the case $`d=2`$ each triangle of the network can be only in a “up” or “down” state. The entropy of this problem is the same as in the three–colouring problem of the honeycomb lattice . The introduction of a bending rigidity induces a first–order crumpling transition which has been studied in . The model with the triangular network embedded in the face–centered–cubic lattice was defined in and a first–order transition has been found also in this case .
In this paper we consider the two–dimensional folding of the square–diagonal lattice and study the phase diagram of this system in terms of the bending rigidities $`\kappa _S`$ and $`\kappa _L`$ relative to short and long edges. Actually, the phase behavior of this model was already considered in by means of a transfer matrix analysis, but the ground state was not properly investigated, and as a consequence some phases were missing. Here we consider the complete set of the ground states of the model and study the phase diagram at finite temperatures by the cluster variation method (CVM) . This method has already proven to be useful for studying vertex models and folding problems .
The outline of the paper is the following. In the next Section we define the model; the ground states will be shown in Section 3. In Section 4 the CVM approximation used in this paper will be described and our results will be presented in Section 5.
2. The square–diagonal folding model
We consider the folding configurations of the square–diagonal lattice where any couple of adjacent triangles can be on the same plane or with the two triangles one on the other. Triangles can be folded with respect to the main axes of the square lattice or with respect to the diagonals so that in the definition of the bending energy different couples of triangles with a short edge or a long edge in common have to be considered. We will start the definition of the model by enumerating the folding states in terms of loop configurations.
2.1. Folding and loop configurations
First fix a reference orientation for all the edges of the lattice in such a way that for each triangle the vectorial sum of the oriented edges is zero. There are two choices for this global orientation and one is shown in Fig. 1. In any folding configuration each short edge is mapped on one of the four vectors $`\pm \stackrel{}{e}_1,\pm \stackrel{}{e}_2`$; this mapping defines the state of the network. Then consider two triangles with a long edge in common as the pair shaded in Fig. 1. This pair can be in two folding states that can be represented in the following way: consider another square lattice as in Fig. 2. The lines joining the centers of the edges are dashed or full depending on whether they are dual to the mapped short edge vectors $`\pm \stackrel{}{e}_1`$ or $`\pm \stackrel{}{e}_2`$ respectively. The obvious observation that in each triangle there is a short edge vector in the direction of $`\pm \stackrel{}{e}_1`$ and an edge in the direction of $`\pm \stackrel{}{e}_2`$ implies that for each center of the edges of the new square lattice there will be two dashed and two full lines. As a consequence the dashed lines and the full lines will form two sets of orthogonal closed loops .
Taking into account also the orientation of $`\stackrel{}{e}_1`$ and $`\stackrel{}{e}_2`$, it comes out that each closed loop has its own orientation independently from the others. Therefore the original folding problem is equivalent to a dense two–loop problem where a sign representing the orientation has to be attributed to each closed loop. Further details can be found in .
Consider now the square plaquettes of this new lattice. In each plaquette there are two dashed and two full lines. On each line there can be two possible signs so that there are $`2\times 2^4=32`$ states for each plaquette. Of course, the sign conservation law along the closed loops implies that only some of the $`32^2`$ configurations of two neighboring plaquettes are allowed. The logarithm of the number of all possible states of these plaquette configurations gives the entropy of the folding problem.
2.2. Bending energy
Here we introduce a bending energy weighting the different folding configurations. We attribute a Boltzmann weight $`e^{\kappa _L}`$ or $`e^{\kappa _L}`$ respectively to each unfolded or folded long edge. Similar weights $`e^{\pm \kappa _S}`$ are associated to the two short edge folding states.
In the plaquette formulation long and short edge bending energies appear in an asymmetric way. Long edge bending energy results in an interaction between each couple of adjacent plaquettes in the loop formulation. An energy $`\kappa _L`$ has to be assigned to the network any time a dashed line (or equivalently a full line) crosses the edge between two plaquettes without bending. In the other case an energy $`\kappa _L`$ has to be assigned to any right angle of dashed lines between neighboring plaquettes (see Fig. 2).
Differently, short edge bending energy results in different weigths for the 32 basic plaquette configurations. In the reference orientation state of Fig. 1, in each square the two short edges on the same diagonal always have opposite orientations. These two short edges acquire the same sign for each folding relative to their perpendicular diagonal. In this case two triangles are folded and an energy $`2\kappa _S`$ has to be assigned to the network. In the loop formulation this is equivalent to have a couple of parallel lines in one plaquette with the same sign. Similar considerations hold for an unfolded edge and for the other couple of parallel lines in a plaquette. Therefore, in total, each plaquette can have an energy equal to $`4\kappa _S`$, 0, or $`4\kappa _S`$ .
2.3. The vertex representation
From the discussion above it is clear that our model can be thought as a vertex model defined on a two–dimensional square lattice $`^2`$. A vertex variable $`\sigma _x\{1,\mathrm{},32\}`$ is associated to each site $`x`$ of the lattice; to each value of $`\sigma _x`$ corresponds one of the $`32`$ vertices introduced above (a possible representation is shown in Fig. 3).
Constraints between neighboring vertices are implied by the sign conservation law along the closed loops. By using the labelling of Fig. 3 these constraints can be written in a simple form: we define the two $`32\times 32`$ matrices $`\mu ^h(\sigma ,\sigma ^{})=\delta _{u_2(\sigma ),u_4(\sigma ^{})}`$ and $`\mu ^v(\sigma ,\sigma ^{})=\delta _{u_3(\sigma ),u_1(\sigma ^{})}`$ where $`\sigma ,\sigma ^{}\{1,\mathrm{},32\}`$ and $`u(\sigma )`$ is the vector $`u`$ associated to the vertex $`\sigma `$. A configuration $`\overline{\sigma }=\{\sigma _x\}_^2`$ is allowed if $`\mu ^h(\sigma _x,\sigma _y)=1`$ and $`\mu ^v(\sigma _x,\sigma _y)=1`$ respectively for any horizontal, $`x,y_h`$, and vertical, $`x,y_v`$, pair of nearest neighboring sites. Notice that on each bond there are only $`256`$ allowed configurations out of the $`32^2=1024`$ total ones. We remark that when we write $`x,y_{h,v}`$ we suppose the sites $`x`$ and $`y`$ lexicographically ordered.
The hamiltonian of the model is in the form
$$H(\overline{\sigma })=\underset{x^2}{}H^s(\sigma _x)+\underset{x,y_h}{}H^h(\sigma _x,\sigma _y)+\underset{x,y_v}{}H^v(\sigma _x,\sigma _y),$$
(2.1)
where the two body potential is given by the two $`32\times 32`$ matrices $`H^v(\sigma ^{},\sigma )=H^h(\sigma ,\sigma ^{})=2\kappa _L\delta _{v(\sigma ),1v(\sigma ^{})}\kappa _L`$ for any $`\sigma ,\sigma ^{}\{1,\mathrm{},32\}`$ and the single body potential is given by the vector $`H^s(\sigma )=4\kappa _S`$ if $`\sigma =1,\mathrm{},8`$, $`H^s(\sigma )=\kappa _S`$ if $`\sigma =25,\mathrm{},32`$ and $`H^s(\sigma )=0`$ otherwise. Finally we can write the partition function of the model:
$$Z(\kappa _S,\kappa _L)=\underset{\overline{\sigma }}{}\left(\underset{x,y_h}{}\mu ^h(\sigma _x,\sigma _y)\underset{x,y_v}{}\mu ^v(\sigma _x,\sigma _y)\right)\mathrm{exp}\{H(\overline{\sigma })\},$$
(2.2)
where the inverse temperature $`\beta `$ has been adsorbed in the hamiltonian.
3. Ground states
The ground states of the network can be conveniently discussed in terms of loop configurations. First consider what happens for $`\kappa _S,\kappa _L>0`$ at very low temperatures. A large positive $`\kappa _L`$ selects states with straight lines so that in a square lattice with periodic boundary conditions there are the dashed and full closed loop shown in Fig. 4(a). A positive value of $`\kappa _S`$ favors alternate sign values on these lines. This corresponds to a completely flat state with an energy per plaquette given by $`E_1=4\kappa _S2\kappa _L`$. There are 8 possible degenerate ground states of this kind and each of them can be realized with two basic plaquette states.
In the sector $`\kappa _S<0,\kappa _L>0`$ elementary triangles want to be folded along diagonals. A positive value of $`\kappa _L`$ still select straight lines. All the dashed lines have one sign and the same for the full lines. The network is in a state where two triangles sharing a long edge are on the same plane and all the other triangles are above one of those two. Also in this case the degeneracy is 8 while one plaquette state is enough for building the global state which has an energy $`E_2=4\kappa _S2\kappa _L`$ per plaquette and is reported in Fig. 4(b).
Loop–configuration ground states for the remaining sectors $`\kappa _S<0,\kappa _L<0`$ and $`\kappa _S>0,\kappa _L<0`$ are respectively shown in Fig. 4(c) and Fig. 4(d). They have an energy equal to $`E_3=4\kappa _S+2\kappa _L`$ and $`E_4=4\kappa _S+2\kappa _L`$. They describe a completely folded state with all triangle one above the other, and a state with four triangles with a common vertex on the same plane and all other triangles folded above. There are needed 2 and 4 different plaquettes for building up ground states in sectors 3 and 4 respectively.
4. The CVM approximation
In this section we describe the CVM approximation we used to study the phase diagram of model (2.2). We do not enter into all the details related to the CVM technique, we just refer to .
The CVM is based on the minimization of the free–energy density functional obtained by a truncation of the cluster (cumulant) expansion of the corresponding functional appearing in the exact variational formulation of statistical mechanics. The existence of different horizontal and vertical interactions in the hamiltonian (2.1) and the structure of the ground states of the model suggest that one should consider at least a square (four–vertex) approximation, that is one should consider a square of four vertices as the largest cluster in the expansion of the free–energy functional. Moreover, in order to reproduce the structure of the ground states we are forced to partition our lattice into four square lattices with spacing two (see Fig. 5). We denote by $`A`$, $`B`$, $`C`$ and $`D`$ these four sublattices and we introduce four square density matrices: $`\rho _{ABCD}`$, $`\rho _{BADC}`$, $`\rho _{CDAB}`$ and $`\rho _{DCBA}`$. We also need four pair density matrices $`\rho _{AB}`$, $`\rho _{BA}`$, $`\rho _{CD}`$ and $`\rho _{DC}`$ for the four horizontal different bonds, four pair density matrices $`\rho _{AC}`$, $`\rho _{CA}`$, $`\rho _{BD}`$ and $`\rho _{DB}`$ for the four different vertical bonds and four single site matrices $`\rho _A`$, $`\rho _B`$, $`\rho _C`$ and $`\rho _D`$. Finally, we define the sets $`𝒫=\{ABCD,BADC,CDAB,DCBA\}`$, $`^h=\{AB,BA,CD,DC\}`$, $`^v=\{AC,CA,BD,DB\}`$ and $`𝒮=\{A,B,C,D\}`$, and we write the free–energy functional per plaquette as follows:
$$\begin{array}{cc}f(\{\rho _X:X𝒫\})=\hfill & \frac{1}{4}_{X^h}\mathrm{Tr}_X\left[\rho _XH^h\right]+\frac{1}{4}_{X^v}\mathrm{Tr}_X\left[\rho _XH^v\right]+\frac{1}{4}_{X𝒮}\mathrm{Tr}_X\left[\rho _XH^s\right]\hfill \\ & \\ & +\frac{1}{4}(_{X𝒫}\mathrm{Tr}_X\left[\rho _X\mathrm{log}\rho _X\right]_{X^h}\mathrm{Tr}_X\left[\rho _X\mathrm{log}\rho _X\right]\hfill \\ & \\ & _{X^v}\mathrm{Tr}_X\left[\rho _X\mathrm{log}\rho _X\right]+_{X𝒮}\mathrm{Tr}_X\left[\rho _X\mathrm{log}\rho _X\right])\hfill \\ & \\ & +_{X𝒫}\lambda _X\left(\mathrm{Tr}\left[\rho _X\right]1\right)\hfill \end{array}$$
(4.1)
where $`\lambda _X`$, with $`X𝒫`$, are four Lagrange multipliers ensuring that $`\rho _X`$ with $`X𝒫`$ are correctly normalized and for any $`X𝒫^h^v𝒮`$ we have denoted by $`\mathrm{Tr}_X`$ the sum over all the allowed configurations on the set $`X`$.
Following the recipe of the CVM one should find the densities $`\{\rho _X:X𝒫\}`$ that minimize the functional (4.1). Hence the next step consists in taking derivatives of the free energy with respect to the square densities. In order to do this we must take into account that bond and single site densities can be obtained from square densities via a partial tracing. Namely,
$$\begin{array}{cccc}\rho _{AB}=\mathrm{Tr}_{CD}\rho _{ABCD}\hfill & \rho _{BA}=\mathrm{Tr}_{DC}\rho _{BADC}\hfill & \rho _{CD}=\mathrm{Tr}_{AB}\rho _{CDAB}\hfill & \rho _{DC}=\mathrm{Tr}_{BA}\rho _{DCBA}\hfill \\ & & & \\ \rho _{AC}=\mathrm{Tr}_{BD}\rho _{ABCD}\hfill & \rho _{BD}=\mathrm{Tr}_{AC}\rho _{BADC}\hfill & \rho _{CA}=\mathrm{Tr}_{DB}\rho _{CDAB}\hfill & \rho _{DB}=\mathrm{Tr}_{CA}\rho _{DCBA}\hfill \\ & & & \\ \rho _A=\mathrm{Tr}_B\rho _{AB}\hfill & \rho _B=\mathrm{Tr}_B\rho _{BA}\hfill & \rho _C=\mathrm{Tr}_D\rho _{CD}\hfill & \rho _D=\mathrm{Tr}_D\rho _{DC}\hfill \end{array}$$
(4.2)
Actually, each bond or single site density can be derived via partial tracing of different higher order densities. This means that suitable Lagrange multiplier must be introduced to ensure that different partial tracings lead to the same result. More precisely, each bond belongs to two different squares (see Fig. 5): we have to associate a family of multiplier to each bond to ensure that the same result is obtained by tracing over the two plaquettes sharing the bond itself (for instance, we need $`\mathrm{Tr}_{CD}\rho _{ABCD}=\mathrm{Tr}_{CD}\rho _{CDAB}`$). We get eight different families $`\{\lambda _X:X^h^v\}`$ and each family contains $`256`$ different multipliers, one for each allowed bond configuration. Moreover, there exist four different bonds sharing the same single site: to each site we associate three different multiplier families $`\{\lambda _{X,i}:X𝒮,i=1,2,3\}`$, each of them made of $`32`$ different multipliers. Hence we have $`20`$ different families of multipliers resulting in a total number of $`8\times 256+12\times 32=2432`$ multipliers.
The functional that we have to minimize is no more the free energy (4.1), but the one in which all the Lagrange multipliers are introduced:
$$\begin{array}{cc}g(\{\rho _X:X𝒫\})=\hfill & 4f(\{\rho _X:X𝒫\})+\mathrm{Tr}_A\left[\lambda _{A,1}\mathrm{Tr}_B\left[\rho _{BA}\rho _{AB}\right]\right]+\mathrm{}\mathrm{}\hfill \\ & \\ & +\mathrm{Tr}_{AB}\left[\lambda _{AB}\mathrm{Tr}_{CD}\left[\rho _{CDAB}\rho _{ABCD}\right]\right]+\mathrm{}\mathrm{}\hfill \end{array}$$
(4.3)
where dots stand for other seven similar bond terms and eleven similar single site terms. Now, let us label with $`\alpha \{1,\mathrm{},4608\}`$ the allowed $`4608`$ square states. We denote by $`\sigma _X(\alpha )`$ the vertex associated to the site $`X𝒮`$ corresponding to $`\alpha `$. To obtain the equilibrium densities we have to set equal to zero the derivatives of the functional (4.3) taken with respect to $`\rho _Y(\alpha )`$ with $`Y𝒫`$ and $`\alpha \{1,\mathrm{},4608\}`$. In the case $`Y=ABCD`$ we obtain:
$$\begin{array}{cc}\rho _{ABCD}(\alpha )=\hfill & \mathrm{const}\times \mathrm{exp}\{H^h(\sigma _A(\alpha ),\sigma _B(\alpha ))H^v(\sigma _A(\alpha ),\sigma _C(\alpha ))H^s(\sigma _A(\alpha ))\hfill \\ & \\ & +\lambda _{AB}(\sigma _A(\alpha ),\sigma _B(\alpha ))\lambda _{CD}(\sigma _C(\alpha ),\sigma _D(\alpha ))+\lambda _{AC}(\sigma _A(\alpha ),\sigma _C(\alpha ))\hfill \\ & \\ & \lambda _{BD}(\sigma _B(\alpha ),\sigma _D(\alpha ))+\lambda _{A,1}(\sigma _A(\alpha ))+\lambda _{A,2}(\sigma _A(\alpha ))\lambda _{B,1}(\sigma _B(\alpha ))\hfill \\ & \\ & \lambda _{B,3}(\sigma _B(\alpha ))\lambda _{C,2}(\sigma _C(\alpha ))+\lambda _{C,3}(\sigma _C(\alpha ))\}\hfill \\ & \\ & \times \rho _{AB}(\sigma _A(\alpha ),\sigma _B(\alpha ))\rho _{AC}(\sigma _A(\alpha ),\sigma _C(\alpha ))\rho _A(\sigma _A(\alpha ))^1\hfill \end{array}$$
(4.4)
The equation above is actually a set of $`4608`$ equations. Three similar sets can be found by considering the cases $`X\{BADC,CDAB,DCBA\}`$. The complete set of $`4\times 4608=18432`$ equations, together with the equations for the multipliers, can be solved by means of the natural iteration method . The equations for the multipliers can be obtained by taking derivatives of (4.3) with respect to the multipliers and are of the form
$$\lambda _X(\sigma _X)=\lambda _X(\sigma _X)+\mathrm{const}\mathrm{log}\frac{\phi _1(\sigma _X)}{\phi _2(\sigma _X)}$$
(4.5)
where $`X^h^v𝒮`$, $`\sigma _X`$ is an allowed configuration on $`X`$ and $`\phi _1`$ and $`\phi _2`$ are linear combinations of higher order densities on clusters $`YX`$ traced over $`YX`$.
5. Phase diagram
The phase diagram of the model, as predicted by our approximation, has been reported in Fig. 6, where open symbols denote second order transitions, and full symbols denote first order transitions. The transition lines have been drawn using a limited number of points, since the precise determination of such points is a very computationally demanding task. In particular, the transitions between the phase D1 (Disordered 1) and the phases F (Folded) and LF (L–Folded) have been simply sketched since it was possible to determine their location only with a rough approximation.
Several phases appear in the diagram. First of all we have four long–range ordered phases corresponding to the four possible ground states described in Section Folding transitions of the square–diagonal two–dimensional latticethanks: PACS numbers: 05.50.+q (Ising problems); 64.60.-i (General studies of phase transitions); 82.65.Dp (Thermodynamics of surfaces and interfaces).. Phase Fl (Flat) is stable for large and positive $`\kappa _S`$ and $`\kappa _L`$ and represents a flat phase. Phase SF (S–Folded) is stable for large enough absolute values of $`\kappa _S<0`$ and $`\kappa _L>0`$. It represents a partially folded phase in which folding occurs mainly along short edges. Phase F is stable for large enough absolute values of $`\kappa _S<0`$ and $`\kappa _L<0`$. It represents a completely folded phase. Phase LF is stable for large enough absolute values of $`\kappa _S>0`$ and $`\kappa _L<0`$. It represents a partially folded phase in which folding occurs mainly at long edges. Like the corresponding ground states, all these phases have degeneracy 8.
On the high temperature side of the ordered flat phase Fl we have a small slice of the disordered phase D2 (Disordered 2). In the central part of the phase diagram, we have the disordered phase D1, which has larger entropy and larger energy than D2. It is noteworthy that the entropy of this phase at the infinite temperature point $`\kappa _S=\kappa _L=0`$ is $`S_{\mathrm{}}0.9204`$ in our approximation, while the estimate in , obtained by transfer matrix methods, corresponds, with our definitions, to $`S_{\mathrm{}}0.9196`$.
Between phases SF and F we have the partially ordered phase PO (Partially Ordered). For large negative $`\kappa _S`$ the transitions between this phase and phases SF and F tend asymptotically to $`|\kappa _L|=\frac{1}{2}\mathrm{ln}\frac{5+\sqrt{17}}{4}0.412`$, which is exactly the estimate for the critical coupling of the square lattice Ising model in the present approximation . This is a consequence of a property of the model, which reduces to a square lattice Ising model in the limit $`\kappa _S\mathrm{}`$. In the loop gas formulation, $`\kappa _S\mathrm{}`$ implies that all loops in the same set (that is, all dashed loops, or all full loops) must have the same sign. Looking at Fig. 3 one can verify that to satisfy this condition the state of the system must be a mixture of only 2 out of the 32 plaquette states. For this pair of allowed states there are four different possibilities (and thus the phase PO has degeneracy 4): (1,5) (corresponding to $`+`$ signs on both loop sets), (2,6) ($`+`$ on full loops and $``$ on dashed loops), (3,7) ($``$ on both loop sets) and (4,8) ($``$ on full loops and $`+`$ on dashed loops). Hence we have the breaking of the loop sign inversion symmetry, which is restored only at the transition to the phase 5. Given a pair of allowed states, the hamiltonian reduces to an ordinary Ising hamiltonian, with equal states on adjacent plaquettes giving a contribution $`\kappa _L`$ to the total energy, and different states giving $`+\kappa _L`$. It is therefore an exact result that, in the limit $`\kappa _S\mathrm{}`$, there must exist the three phases SF, F and PO (corresponding respectively to the ferromagnetic, antiferromagnetic and disordered phases of the Ising model), separated by second order phase transitions at $`|\kappa _{Lc}|=\frac{1}{2}\mathrm{ln}\left(1+\sqrt{2}\right)0.441`$.
A similar situation occurs in the limit $`\kappa _L\mathrm{}`$. Looking at Fig. 4 one sees that in this limit both the full and the dashed lines form the smallest possible square loops. It follows that the model reduces to two decoupled Ising models, where the Ising variables are the loop signs and two loops (both full, or both dashed) having parallel edges in the same plaquette interact with an energy $`\pm 2\kappa _S`$. Therefore the phases F and LF (corresponding respectively to the antiferromagnetic and ferromagnetic phases of the limiting Ising model) will undergo second order phase transitions towards the disordered phase D1 at $`|\kappa _{Sc}|=\frac{1}{4}\mathrm{ln}\left(1+\sqrt{2}\right)0.220`$. The estimate for this value in our approximation is of course $`|\kappa _L|=\frac{1}{4}\mathrm{ln}\frac{5+\sqrt{17}}{4}0.206`$.
Comparing our phase diagram with that by Di Francesco one can see several striking differences. First of all, among our ordered phases, only the completely flat and completely folded ones were reported in , and no low temperature transition was found among them. In addition, no intermediate phase (like our phases PO and D2) was found between the ordered and the disordered phases. We note that the CVM results are fully confirmed by the ground states analysis and the limits $`\kappa _S\mathrm{}`$ and $`\kappa _L\mathrm{}`$.
In conclusions, we have shown that the CVM techniques can be adapted to study very complex vertex models as the one considered in this paper representing the folding configurations of a square–diagonal lattice. As in the case of the folding of the triangular lattice, a first–order crumpling transition between the flat phase and a disordered phase has been found. The extension of this analysis to the $`d`$–dimensional folding problem appears not easy from a numerical point of view.
Acknowledgements
G.G. acknowledges support by MURST (COFIN97). E.C. whishes to express his thanks to the European network “Stochastic Analysis and its Applications” ERB–FMRX–CT96–0075 for financial support. |
warning/0002/cond-mat0002147.html | ar5iv | text | # The ground state of Sr3Ru2O7 revisited; Fermi liquid close to a ferromagnetic instability
## Abstract
We show that single-crystalline Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> grown by a floating-zone technique is an isotropic paramagnet and a quasi-two dimensional metal as spin-triplet superconducting Sr<sub>2</sub>RuO<sub>4</sub> is. The ground state is a Fermi liquid with very low residual resistivity ($``$3 $`\mu \mathrm{\Omega }`$cm for in-plane currents) and a nearly ferromagnetic metal with the largest Wilson ratio $`R_\mathrm{W}10`$ among paramagnets so far. This contrasts with the ferromagnetic order at $`T_\mathrm{C}`$=104 K reported on single crystals grown by a flux method \[Cao et al., Phys. Rev. B $`\mathrm{𝟓𝟓}`$, R672 (1997)\]. However, we have found a dramatic changeover from paramagnetism to ferromagnetism under applied pressure. This suggests the existence of a substantial ferromagnetic instability in the Fermi liquid state.
* PACS number(s): PACS numbers: 75.40.-s, 71.27.+a, 75.30.Kz
The discovery of superconductivity in the single-layered perovskite Sr<sub>2</sub>RuO<sub>4</sub> <sup>?</sup> has motivated the search for new superconductors and anomalous metallic materials in Ruddlesden-Popper (R-P) type ruthenates (Sr,Ca)<sub>n+1</sub>Ru<sub>n</sub>O<sub>3n+1</sub>. The recent determination of the spin-triplet pairing in its superconducting state suggests that ferromagnetic (FM) correlations are quite important in Sr<sub>2</sub>RuO<sub>4</sub> ,<sup>?</sup> and the existence of enhanced spin fluctuations has been suggested by nuclear magnetic resonance (NMR) .<sup>??</sup> On the other hand, the recent report has shown that enhanced magnetic excitations around $`𝐪`$=0 is not detected but sizeable excitations have been seen around finite $`𝐪`$ in Sr<sub>2</sub>RuO<sub>4</sub> by inelastic neutron scattering.<sup>?</sup> This has stimulated debate on the mechanism of the spin-triplet superconductivity, which had been naively believed to have a close relation to FM ($`𝐪`$=0) spin excitations. Hence, it is desirable to investigate its related compounds as described below.
The simple perovskite (three dimensional) metallic SrRuO<sub>3</sub> ($`n=\mathrm{}`$) has been well known to order ferromagnetically below 160 K with a magnetic moment $`M=0.81.0\mu _\mathrm{B}/`$Ru .<sup>??</sup> FM perovskite oxides are relatively rare except for metallic manganites. For pure thin film SrRuO<sub>3</sub>, analyses of quantum oscillations in the resistivity have given good evidence for the Fermi liquid behavior .<sup>?</sup>
The double layered perovskite Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> ($`n`$=2) is regarded as having an intermediate dimensionality between the systems with $`n=1`$ and $`n=\mathrm{}`$ .<sup>?</sup> Investigations on polycrystalline Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> showed a magnetic-susceptibility maximum around 15 K with Curie-Weiss-like behavior above 100 K and a metallic temperature dependence of the electrical resistivity .<sup>??</sup>
In the study presented here, we have for the first time succeeded in growing single crystals of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> by a floating-zone (FZ) method. Those single crystals (FZ crystals) do not contain any impurity phases (e.g. SrRuO<sub>3</sub>) which was observed in polycrystals .<sup>?</sup> We report herein that the FZ crystal of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is a nearly FM paramagnet (enhanced paramagnet) and a quasi-two dimensional metal with a strongly-correlated Fermi liquid state. In addition, we have performed magnetization measurements under hydrostatic pressures up to 1.1 GPa in order to confirm whether the FM instability is susceptible to pressure. The results suggest that there is a changeover from paramagnetism to ferromagnetism, indicating a strong FM instability. Essential features of magnetism for FZ crystals as well as polycrystals are inconsistent with the appearance of a FM ordering ($`T_\mathrm{c}=104\mathrm{K}`$) at ambient pressure for single crystals grown by a flux method <sup>?</sup> using SrCl<sub>2</sub> flux and Pt crucibles. We will argue that FZ crystals reflect the intrinsic behavior of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub>.
Details of the FZ crystal growth are explained elsewhere .<sup>?</sup> The crystal structure of the samples at room temperature was characterized by powder x-ray diffraction. Electrical resistivity $`\rho (T)`$ was measured by a standard four terminal dc-technique from 4.2 K to 300 K and by an ac method from 0.3K to 5K. Specific heat $`C_P(T)`$ was measured by a relaxation method from 1.8 K to 35 K (Quantum Design, PPMS). The temperature dependence of magnetic susceptibility $`\chi (T)M/H`$ from 2 K to 320 K was measured using a commercial SQUID magnetometer (Quantum Design, MPMS-5S). For magnetization measurements of FZ crystals at ambient pressure, we performed sample rotation around the horizontal axis, normal to the scan direction, using the rotator in MPMS-5S. We could align the crystal axes exactly parallel to a field direction within 0.2 degree using this technique. For high pressures, we measured magnetization using a long-type hydrostatic pressure micro-cell <sup>?</sup> with the SQUID magnetometer. Loaded pressures around 3 K were determined from the shift of superconducting transition temperature of Sn in the micro-cell in a 5 mT field .
The R-P type structure of $`n`$=2 for FZ crystals of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> was confirmed by the powder x-ray diffraction patterns with crushed crystals, which indicated no impurity peaks. Recently, the crystal structure of polycrystalline Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> has been refined by neutron powder diffraction .<sup>??</sup> Although they have concluded that symmetry of the structure is orthorhombic owing to the rotation of the RuO<sub>6</sub> octahedron about the c-axis by about 7 degrees, we deduced lattice parameters at room temperature by assuming tetragonal $`I4/mmm`$ symmetry as $`a=3.8872(4)\mathrm{\AA }`$, and $`c=20.732(3)\mathrm{\AA }`$. These values are in good agreement with those of polycrystals obtained by neutron diffraction <sup>??</sup> and x-ray diffraction .<sup>?</sup>
The temperature dependence of magnetic susceptibility $`\chi (T)=M/H`$ in a field of 0.3 T is shown in Fig. 1. No hysteresis is observed between zero-field cooling (ZFC) and field cooling (FC) sequences, so we conclude that there is no ferromagnetic ordering. Little magnetic anisotropy is observed in contrast to large anisotropy ($`10^2`$) of flux-grown crystals .<sup>?</sup> The nearly isotropic susceptibility of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is qualitatively similar to that of the enhanced Pauli-paramagnetic susceptibility in Sr<sub>2</sub>RuO<sub>4</sub> .<sup>?</sup> For an applied field of 0.3 T, there is no in-plane anisotropy of the susceptibility for the whole temperature range (2 K $`T`$ 300K), within the precision of our equipment (1$`\%`$).
The susceptibility for both $`H//`$ab and $`H//`$c exhibits Curie-Weiss behavior above 200 K. We have fitted the observed $`\chi (T)`$ from 200 K to 320 K with $`\chi (T)=\chi _0+\chi _{\mathrm{CW}}(T)`$, where $`\chi _0`$ is the temperature independent term and $`\chi _{\mathrm{CW}}(T)=C/(T\Theta _\mathrm{W})`$ is the Curie-Weiss term. The effective Bohr magneton numbers $`p_{\mathrm{eff}}`$ deduced from $`C`$ are $`p_{\mathrm{eff}}=`$ 2.52 (2.99) and $`\Theta _\mathrm{W}`$ =$`-`$ 39 K ($`-`$ 45 K) for $`H//ab(H//c)`$. The negative values of $`\Theta _\mathrm{W}`$ normally indicate antiferromagnetic (AFM) correlations in the case of localized-spin systems. However, we cannot conclude that AFM correlations play an important role solely by the negative $`\Theta _\mathrm{W}`$ in an metallic system like Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> .<sup>?</sup>
Around $`T_{\mathrm{max}}`$ =16 K, $`\chi (T)`$ shows a maximum for both $`H//`$ab and $`H//`$c. The maximum has been also observed in the polycrystals. The results of temperature dependence of specific heat, NMR <sup>?</sup> and elastic neutron scattering <sup>??</sup> for polycrystals indicate that there is no evidence for any long range order with definite moments. The FZ crystal shows nearly isotropic $`\chi (T)`$ for all crystal axes below $`T_{\mathrm{max}}`$. Hence, the maximum cannot be accredited to the long range AFM order. Therefore, we conclude Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> to be a $`paramagnet`$. Concerning $`\chi (T)`$ under higher fields, $`T_{\mathrm{max}}`$ is suppressed down to temperatures below 5 K above 6 T .<sup>?</sup> Such a maximum in $`\chi (T)`$ and a field dependent $`T_{\mathrm{max}}`$ are often observed in a nearly ferromagnetic (enhanced paramagnetic) metal like TiBe<sub>2</sub> <sup>?</sup> or Pd .<sup>?</sup> In addition, a similar behavior in $`\chi (T)`$ has been observed in (Ca,Sr)<sub>2</sub>RuO<sub>4</sub> <sup>?</sup> and MnSi ,<sup>?</sup> which are recognized as examples of a critical behavior by spin fluctuations. Similar critical behavior, originating especially from FM spin fluctuations, is also expected in Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub>. Nevertheless, we cannot rule out the possibility of AFM correlations as observed in Sr<sub>2</sub>RuO<sub>4</sub>, caused by the nesting of its Fermi surfaces with the vector $`𝐐`$ $`=(\pm 0.6\pi /a,\pm 0.6\pi /a,0)`$ .<sup>?</sup>
As shown in Fig.2, the specific heat coefficient of the FZ crystal of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is $`\gamma `$ = 110 mJ/(K<sup>2</sup> Ru mol) somewhat larger compared to other R-P type ruthenates {$`\gamma =`$80 mJ/(K<sup>2</sup> Ru mol) for CaRuO<sub>3</sub>, 30 mJ/(K<sup>2</sup> Ru mol) for SrRuO<sub>3</sub> <sup>?</sup> and 38 mJ/(K<sup>2</sup> Ru mol) for Sr<sub>2</sub>RuO<sub>4</sub> <sup>??</sup>}. This suggests that Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is a strongly-correlated metallic oxide. For polycrystalline Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub>, we obtained the value $`\gamma =63`$ mJ/(K<sup>2</sup> Ru mol) using an adiabatic method .<sup>?</sup>
The temperature dependence of the electrical resistivity $`\rho (T)`$ is shown in Fig. 3 above 0.3 K. Both $`\rho _{\mathrm{ab}}(T)`$ and $`\rho _\mathrm{c}(T)`$ are metallic ($`d\rho /dT>0`$) in the whole region. The ratio of $`\rho _c`$/$`\rho _{ab}`$ is about 300 at 0.3 K and 40 at 300 K. This anisotropic resistivity is consistent with the quasi-two-dimensional Fermi surface sheets obtained from the band-structure calculations .<sup>?</sup> With lowering temperature below 100 K, a remarkable decrease of $`\rho _\mathrm{c}(T)`$ is observed around 50 K. This is probably due to the suppression of the thermal scattering with decreasing temperature between quasi-particles and phonons as observed in Sr<sub>2</sub>RuO<sub>4</sub> .<sup>???</sup> Thus, below 50K, interlayer hopping propagations of the quasi-particle overcome the thermal scattering with phonons. This hopping picture for $`\rho _\mathrm{c}(T)`$ is well consistent with the large value of $`\rho _\mathrm{c}(T)`$ and nearly cylindrical Fermi surfaces. On the other hand, $`\rho _{\mathrm{ab}}(T)`$ shows a change of the slope around 20 K. Such a change in $`\rho _{\mathrm{ab}}(T)`$ has also been reported for Sr<sub>2</sub>RuO<sub>4</sub> under hydrostatic pressure ($``$ 3 GPa). That might be possibly due to the enhancement of ferromagnetic spin fluctuations .<sup>?</sup>
As shown in the inset of Fig. 3, the resistivity yields a quadratic temperature dependence below 6 K for both $`\rho _{\mathrm{ab}}(T)`$ and $`\rho _\mathrm{c}(T)`$ , characteristic of a Fermi liquid as observed in Sr<sub>2</sub>RuO<sub>4</sub> .<sup>?</sup> We fitted $`\rho _{\mathrm{ab}}(T)`$ by the formula $`\rho _{\mathrm{ab}}(T)`$=$`\rho _0`$+$`AT^2`$ below 6 K and obtained $`\rho _0`$=2.8 $`\mu \mathrm{\Omega }`$ cm and $`A`$=0.075 $`\mu \mathrm{\Omega }`$cm/K<sup>2</sup>. Since the susceptibility is quite isotropic and temperature independent below 6 K, the ground state of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is ascribable as a Fermi liquid. We now can estimate Kadowaki-Woods ratio $`A/\gamma ^2`$. Assuming that electronic specific heat $`\gamma `$ = 110 mJ/(K<sup>2</sup> Ru mol) is mainly due to the ab-plane component, we obtain $`A/\gamma ^2A_{\mathrm{ab}}/\gamma ^2=0.6\times 10^{-5}\mu \mathrm{\Omega }`$ cm/(mJ/K<sup>2</sup> Ru mol)<sup>2</sup> close to that observed in heavy fermion compounds.
Regarding to $`\chi (T)`$ again, it is important to note that even at temperatures much lower than $`T_{\mathrm{max}}`$, $`\chi (T)`$ remains quite large. It appears that the ground state maintains a highly enhanced value of $`1.5\times 10^2`$ emu/Ru mol, comparable to that obtained for typical heavy fermion compounds. Considering that the observed $`\chi `$ is dominated by the renormalized quasi-particles, we can estimate the Wilson ratio $`R_\mathrm{W}=7.3\times 10^4\times \chi (\mathrm{emu}/\mathrm{mol})/\gamma (\mathrm{mJ}/(\mathrm{K}^2\mathrm{mol})`$) in the ground state. If we regard the observed values at $`T`$ = 2 K as that at $`T`$ = 0 K, we have $`R_\mathrm{W}=10(18)`$ using $`\gamma `$ for single crystals (polycrystals). Despite the difference in the $`\gamma `$ value between polycrystals and single crystals, $`R_\mathrm{W}`$ is much greater than unity. This large value implies that FM correlations are strongly enhanced in this compound, especially when compared with the values of 12 for TiBe<sub>2</sub> and 6 for Pd .<sup>?</sup> Therefore, the ground state of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is characterized by strongly-correlated Fermi liquid behavior with enhanced FM spin fluctuations, i.e. Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is a strongly-correlated nearly FM metal.
Concerning this FM correlations, it should be noted that, using single crystals grown by a chlorine flux method with Pt crucibles ,<sup>?</sup> Cao et al. have investigated remarkable magnetic and transport properties of R-P ruthenates <sup>?</sup> prior to our crystal growth. The ground state of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> was concluded to be an itinerant ferromagnet with $`T_\mathrm{c}=104`$ $`\mathrm{K}`$ and an ordered moment $`M=1.2`$ $`\mu _\mathrm{B}/`$Ru. The flux-grown crystals were reported to have a residual resistivity ($`\rho _0=3`$ m$`\mathrm{\Omega }`$ cm ) $`10^3`$ times greater than that of FZ crystals ($`\rho _0=3\mu \mathrm{\Omega }`$ cm ) for in-plane transport. In addition, FZ crystals reveal $`T`$-square dependent resistivities at low temperatures as already shown, which was not observed in flux-grown crystals. In general, the FZ method with great care can be impurity-free crystal growth, while the flux method tends to contaminate crystals due to impurity elements from both the flux and the crucible. This might be a main reason why the resistivity is much higher for flux-grown crystals. Thus, we suppose with assurance that the data from FZ crystals reflect the intrinsic nature of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> better than those from flux-grown crystals.
In order to acquire the information of the magnetic instability in the FZ crystal of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub>, we have measured magnetization under hydrostatic pressure up to 1.1 GPa. The temperature dependence of magnetization $`M(T)`$ is shown for several pressures under a 0.1 T field along c-axis in Fig. 4. Around 1 GPa, substantial increase is recognized below around 70 K with a clear FM component indicated by the difference between ZFC and FC sequences. Although the remanent moment at 2 K ($`M0.08`$ $`\mu _\mathrm{B}/`$Ru) is much smaller than that expected for $`S`$=1 of Ru<sup>4+</sup>, its susceptibility is quite large (0.4 emu/Ru mol). We infer that this transition is a FM ordering of itinerant Ru<sup>4+</sup> spins. In Fig. 4, we also show the field dependence of magnetization $`M(H`$//c) at 2 K for $`P=0.1`$ MPa and $`P=1`$ GPa. Obvious ferromagnetic component appears at lower fields for $`P=1`$ GPa. Even at higher fields, increase in magnetization by pressure is also present as at lower fields. This feature endorses the drastic changeover from paramagnetism to ferromagnetism induced by applied pressure. To the best of our knowledge, this is the first example of the pressure- induced changeover from Fermi liquid to ferromagnetism.
For the purpose of understanding the observed behavior, we should begin with Stoner theory. In the metallic state with correlated electrons, the ferromagnetic order is driven by the Stoner criterion $`U_{\mathrm{eff}}N(E_\mathrm{F})1`$, where $`U_{\mathrm{eff}}`$ is an effective Coulomb repulsion energy. The systematics of band-width $`W`$ and the density of states $`N(E_\mathrm{F}`$) in the R-P ruthenates is summarized by Maeno $`et`$ $`al`$ .<sup>?</sup> In this system, increasing $`n`$ from 1 to $`\mathrm{}`$ causes enhancement of $`N(E_\mathrm{F}`$) as well as $`W`$. This is opposite to the single band picture, i.e. increasing $`N(E_\mathrm{F})`$ naively means decreasing $`W`$. In the case of R-P ruthenates, the anomalous variation might be due to the modifications of the degeneracy of three $`t_{2g}`$ orbitals for Ru-$`4d`$ electrons. According to the summary ,<sup>?</sup> ferromagnetic SrRuO<sub>3</sub> is characterized by the highest $`N(E_\mathrm{F}`$) and $`W`$ among them, satisfying the Stoner criterion. This implies that the enlargement of $`N(E_\mathrm{F}`$) and $`W`$ reflects stronger three dimensionality in the R-P ruthenates. Hence, applying pressure probably makes Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> closer to SrRuO<sub>3</sub>, leading to FM order. For further investigations, it is required that structural study, resistivity and specific heat under pressures will be performed.
In conclusion, by using the floating-zone method we have succeeded for the first time in growing single crystals of Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> with very low residual resistivity in comparison with that of flux-grown crystals reported previously. The results of magnetization, resistivity and specific heat measurements suggest that Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub> is a strongly-correlated Fermi liquid with a nearly ferromagnetic ground state, consistent with the observation of ferromagnetic ordering below 70 K under applied pressure ($`P`$1 GPa). As far as we know, this is the first example of the pressure-induced changeover from Fermi liquid to ferromagnetism. This ferromagnetic ordering may guarantee the existence of the ferromagnetic spin fluctuations in Sr<sub>3</sub>Ru<sub>2</sub>O<sub>7</sub>.
Authors are very grateful to A. P. Mackenzie for his fruitful advice and critical reading of this manuscript. They thank T. Ishiguro, T. Fujita and K. Matsushige for their helpful supports. They thank S. R. Julian, G.G. Lonzarich, G. Mori, D.M. Forsythe, R.S. Perry, K. Yamada, Y. Takahashi, and M. Sigrist for their useful discussions and technical supports. They also thank N. Shirakawa for his careful reading of this manuscript. |
warning/0002/quant-ph0002069.html | ar5iv | text | # Untitled Document
Quantum oscillator as 1D anyon
Ye. Hakobyan<sup>a)</sup>, V. Ter-Antonyan<sup>b)</sup>
Bogoliubov Laboratory of Theoretical Physics,
Joint Institute for Nuclear Research,
Dubna, Moscow Region, 141980, Russia
It is shown that in one spatial dimension the quantum oscillator is dual to the charged particle situated in the field described by the superposition of Coulomb and Calogero–Sutherland potentials.
Keywords: 1D quantum oscillator, reduction, $`1/x^2`$ interaction, duality, 1D anyon.
I. INTRODUCTION
In one spatial dimension a particle moving in the Calogero–Sutherland potential $`V_{cs}=\mathrm{}^2\nu (1\nu )/2\mu x^2`$ has a very unusual property. Unlike the potential $`V_{cs}`$, the wave function is not invariant under the replacement $`\nu (1\nu )`$. It describes a boson for even $`\nu `$ and a fermion for odd $`\nu `$. Statistics corresponding to the other values of $`\nu `$ is called the fractional statistics<sup>1</sup>, and the system influenced along with $`V_{cs}`$ by a potential binding the particle to the center is called the 1D anyon<sup>2-4</sup>. Nobody has observed a 1D anyon yet, but nevertheless it is of both theoretical<sup>5</sup> and experimental<sup>6</sup> interest. The purpose of the present note is to prove that such an extraordinary object can be constructed from a 1D quantum oscillator.
II. ANYON–OSCILLATOR DUALITY
Consider the Schrödinger equation
$$_u^2\mathrm{\Psi }+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)\mathrm{\Psi }=0,$$
(1)
which describes the 1D quantum oscillator. Introduce the quantum number $`s=0,1/2`$ and write $`N=2n+2s`$, with $`N`$ numerating the energy levels $`E=\mathrm{}\omega (N+1/2)`$ and $`n`$ being integer and nonnegative. Without loss of information we can assume $`u`$ to belong to the region $`0u<\mathrm{}`$. We interpret $`s`$ as a spin of the reduced oscillator. The corresponding wave function is denoted by $`\mathrm{\Psi }_n^{(s)}`$.
Let us look for the function $`\mathrm{\Psi }_n^{(s)}`$ in the form
$$\mathrm{\Psi }_n^{(s)}(u)=Cu^{2s}\overline{\mathrm{\Psi }}_n,$$
(2)
where $`\overline{\mathrm{\Psi }}_n`$ is subordinate to the condition $`\overline{\mathrm{\Psi }}_n(0)0`$, and $`C`$ is a normalization constant. Eq. (1) is easily seen to take the form
$$_u^2\overline{\mathrm{\Psi }}_n+\frac{4s}{u}_u\overline{\mathrm{\Psi }}_n+\frac{2\mu }{\mathrm{}^2}\left(E\frac{\mu \omega ^2u^2}{2}\right)\overline{\mathrm{\Psi }}_n=0.$$
(3)
After change of the variable $`x=u^2`$, we arrive at the equation $`(2\nu =2s+1/2)`$
$$_x^2\overline{\mathrm{\Psi }}_n+\frac{2\nu }{x}_x\overline{\mathrm{\Psi }}_n+\frac{2\mu }{\mathrm{}^2}\left(\frac{\mu \omega ^2}{8}+\frac{E}{4x}\right)\overline{\mathrm{\Psi }}_n=0.$$
(4)
Now we set
$$\overline{\mathrm{\Psi }}_n=x^\nu \mathrm{\Phi }_n^{(\nu )},$$
(5)
then cancel the undesirable term with first derivative in (4) and obtain
$$_x^2\mathrm{\Phi }_n^{(\nu )}+\frac{2\mu }{\mathrm{}^2}\left(\epsilon V_cV_{cs}\right)\mathrm{\Phi }_n^{(\nu )}=0,$$
(6)
where $`V_c=\alpha /x`$, $`V_{cs}`$ is the Calogero–Sutherland potential with $`\nu =1/4`$ or $`3/4`$ and
$$\epsilon =\frac{\mu \omega ^2}{8},\alpha =\frac{E}{4}.$$
(7)
Eq. (6) describes a system which we call the 1D Coulomb anyon.
This equation realizes a special case of a more general equation that has a relation to $`(2+1)`$-dimensional anyons<sup>7</sup>.
Comparing Eq. (1) with Eqs. (6) and (7), we summarize that there are two alternative possibilities connected with Eq. (1) – explicit and hidden. In the first case, the parameter $`\omega `$ is fixed ($`\omega =fix.>0`$) and plays a role of coupling constant, the parameter $`E`$ is quantized and has a meaning of energy, and the system is a 1D quantum oscillator. For the hidden possibility, the parameter $`E`$ is fixed ($`E=fix.>0`$), the coupling constant is equal to $`E/4`$, $`\omega `$ is quantized, the meaning of energy takes the quantity $`\epsilon =\mu \omega ^2/8`$, and the system is the 1D Coulomb anyon. In the above-mentioned sense, the 1D quantum oscillator is dual to the 1D Coulomb anyon.
III. ENERGY LEVELS AND WAVE FUNCTIONS
Let us return to Eq. (6) and make the substitution
$$\mathrm{\Phi }_n^{(\nu )}=y^\nu e^{y/2}Q(y),$$
(8)
where $`y=x(8\mu \epsilon /\mathrm{}^2)^{1/2}`$ and $`Q(0)0`$ and is finite. The function $`Q(y)`$ can diverge at infinity but not higher than the finite power of $`y`$. Using (8) and (6) we come to the equation
$$y_y^2Q+(2\nu y)_yQ(\nu \lambda )Q=0,$$
(9)
with $`\lambda =(\mu \alpha ^2/2\mathrm{}^2\epsilon )^{1/2}`$. Eq. (9) is the equation for a confluent hypergeometric function. It has a general solution<sup>8</sup>
$$Q(y)=C_1F(\nu \lambda ,2\nu ,y)+C_2y^{12\nu }F(1\lambda \nu ,22\nu ,y),$$
(10)
where $`F(a,b,y)`$ is given by the series
$$F(a,b,y)=1+\frac{a}{b}\frac{y}{1!}+\frac{a(a+1)}{b(b+1)}\frac{y^2}{2!}+\mathrm{}$$
convergent for all finite $`y`$. For large $`y`$ the asymptotic formula<sup>8</sup> is valid
$$F(a,b,y)\frac{\mathrm{\Gamma }(b)}{\mathrm{\Gamma }(ba)}(y)^a+\frac{\mathrm{\Gamma }(b)}{\mathrm{\Gamma }(a)}e^y(y)^{ab}.$$
(11)
The second term in (10) for $`\nu =3/4`$ is singular at $`y=0`$, and hence $`C_2`$ has to be taken zero. The first term in (10), as it is evident from (11), is “well-behaved” at infinity under the condition $`3/4\lambda =n`$, where $`n`$ is an integer number greater or equal to zero. For $`\nu =1/4`$ both the terms in (10) are regular at $`y=0`$, but the satisfactory behavior at infinity needs the simultaneous requirements $`1/4\lambda =n`$, $`3/4\lambda =m`$, or $`nm=1/2`$, which is impossible. Hence, either $`C_1=0`$ or $`C_2=0`$. But for $`C_1=0`$ the function $`Q(y)`$ will become zero at $`y=0`$. This contradicts the condition $`Q(0)0`$, and, therefore, we put $`C_2=0`$ and $`1/4\lambda =n`$. Thus, we conclude that $`\nu \lambda =n`$, i.e.,
$$\epsilon _n^{(\nu )}=\frac{\mu \alpha ^2}{2\mathrm{}^2(n+\nu )^2},n=0,1,2,\mathrm{}$$
(12)
Returning to the corresponding eigenfunctions, we put
$$\mathrm{\Phi }_n^{(\nu )}=C_n^{(\nu )}y^\nu e^{y/2}F(n,2\nu ,y).$$
(13)
It is known<sup>9</sup> that
$$F(n,2\nu ,y)=\frac{n!\mathrm{\Gamma }(2\nu )}{[\mathrm{\Gamma }(n+2\nu )]^2}L_n^{2\nu 1}(y),$$
and
$$\underset{0}{\overset{\mathrm{}}{}}e^yy^{2\nu }[L_n^{2\nu 1}(y)]^2𝑑y=2(n+\nu )\frac{[\mathrm{\Gamma }(n+2\nu )]^3}{n!},$$
where $`L_n^{2\nu 1}(y)`$ is an associated Laguerre polynomial. Using this results and taking into account the relation
$$\left(\frac{8\mu \epsilon }{\mathrm{}^2}\right)^{1/4}=\frac{1}{\mathrm{}}\left(\frac{2\mu \alpha }{n+\nu }\right)^{1/2},$$
we find
$$C^{(\nu )}=\frac{\sqrt{\mu \alpha }}{\mathrm{}}\frac{1}{n+\nu }\frac{1}{\mathrm{\Gamma }(2\nu )}\sqrt{\frac{\mathrm{\Gamma }(n+2\nu )}{n!}}.$$
Summarizing, we write
$$\mathrm{\Phi }_n^{(\nu )}=\frac{\sqrt{\mu \alpha }}{\mathrm{}}\frac{1}{n+\nu }\frac{1}{\mathrm{\Gamma }(2\nu )}\sqrt{\frac{\mathrm{\Gamma }(n+2\nu )}{n!}}y^\nu e^{y/2}F(n,2\nu ,y).$$
(14)
So, we have two types of the 1D Coulomb anyons with $`\nu =1/4`$ and $`\nu =3/4`$. They are dual to reduced oscillators with $`s=0`$ and $`s=1/2`$, respectively.
IV. DUALITY FOR SOLUTIONS
Now we will calculate the energy levels $`\epsilon _n`$ and wave functions $`\mathrm{\Phi }_n^{(\nu )}`$ in another, more straightforward, way. For energy levels we have
$$\epsilon =\frac{\mu \omega ^2}{8}=\frac{\mu }{8}\left[\frac{E}{2\mathrm{}(n+\nu )}\right]^2=\frac{\mu }{8}\left[\frac{2\alpha }{\mathrm{}(n+\nu )}\right]^2=\frac{\mu \alpha ^2}{2\mathrm{}^2(n+\nu )^2}.$$
It follows from Eqs. (2) and (5) that
$$\mathrm{\Phi }_n^{(\nu )}=\frac{1}{C}x^{1/4}\mathrm{\Psi }_n^{(\nu )}$$
and, therefore,
$$\underset{0}{\overset{\mathrm{}}{}}|\mathrm{\Phi }_n^{(\nu )}|^2𝑑x=\frac{1}{|C|^2}_0^{\mathrm{}}x^{1/2}|\mathrm{\Psi }_n^{(s)}|^2𝑑x.$$
The integral in the left-hand side is equal to 1, from which it follows that
$$|C|^2=2\underset{\mathrm{}}{\overset{\mathrm{}}{}}u^2|\mathrm{\Psi }_N(u)|^2𝑑u=2\overline{u^2}=\frac{4(n+\nu )\mathrm{}}{\mu \omega },$$
where $`\mathrm{\Psi }_N`$ is the normalized wave function of a 1D quantum oscillator. Thus,
$$\mathrm{\Phi }_n^{(\nu )}=\frac{(1)^n}{2}\sqrt{\frac{\mu \omega }{\mathrm{}(n+\nu )}}x^{1/4}\mathrm{\Psi }_n^{(s)}.$$
(15)
Remind that according to the theory of quantum oscillator<sup>9</sup>,
$$\mathrm{\Psi }_n^{(s)}=\sqrt{2}\left(\frac{\mu \omega }{\pi \mathrm{}}\right)^{1/4}\frac{1}{2^NN!}e^{\mu \omega u^2/2}H_N\left(u\sqrt{\frac{\mu \omega }{\mathrm{}}}\right).$$
(16)
Further, it is known<sup>10</sup> that Hermite polynomials could be expressed in terms of confluent hypergeometric functions. For our case
$$H_{2n+2s}(\sqrt{y})=(1)^n\frac{(2n+2s)!}{n!}(2\sqrt{y})^{2s}F(n,2s+1/2,y).$$
(17)
Using the identification $`y=x\mu \omega /\mathrm{}`$ and the relations $`2s+1/2=2\nu `$ and $`\mu \omega /\mathrm{}=2\mu \alpha /\mathrm{}^2(n+\nu )`$, and taking into account Eqs. (15)-(17) we get
$$\mathrm{\Phi }_n^{(\nu )}=\stackrel{~}{C}_n^{(\nu )}y^\nu e^{y/2}F(n,2\nu ,y),$$
(18)
where
$$\stackrel{~}{C}_n^{(\nu )}=\sqrt{\frac{\mu \alpha }{\mathrm{}^2}}\frac{1}{2^{n\nu +1/4}}\frac{\sqrt{\mathrm{\Gamma }(2n+2\nu +1/2)}}{\pi ^{1/4}n!(n+\nu )},$$
(19)
or more explicitly
$$\stackrel{~}{C}_n^{(1/4)}=\frac{\sqrt{\mu \alpha }}{\mathrm{}}\frac{1}{2^n}\frac{\sqrt{\mathrm{\Gamma }(2n+1)}}{\pi ^{1/4}n!(n+1/4)},$$
$$\stackrel{~}{C}_n^{(3/4)}=\frac{\sqrt{\mu \alpha }}{\mathrm{}}\frac{1}{2^{n1/2}}\frac{\sqrt{\mathrm{\Gamma }(2n+2)}}{\pi ^{1/4}n!(n+3/4)}.$$
From the duplication formula for a gamma-function
$$\mathrm{\Gamma }(2z)=2^{2z1}\pi ^{1/2}\mathrm{\Gamma }(z)\mathrm{\Gamma }(z+1/2)$$
and taking into account that $`\mathrm{\Gamma }(1/2)=\pi ^{1/2}`$, $`\mathrm{\Gamma }(3/2)=\frac{1}{2}\pi ^{1/2}`$, we conclude that $`\stackrel{~}{C}_n^{(\nu )}=C_n^{(\nu )}`$ and, consequently, Eqs. (18) and (14) are identical.
V. CONCLUSIONS
a) The 1D oscillator has only a discrete energy spectrum and, therefore, is a model provided by the property which is known in QCD as confinement. A particle situated in the confinement potential cannot be removed from the center and transferred to infinity. On the other hand, the 1D Coulomb anyon is a system possessing both the discrete and continuous part in the energy spectrum. At the same time, it includes $`1/x^2`$ interaction and, therefore, pretends to be a magnetic monopole in one spatial dimension. All these ideas confirm that our result can be interpreted in the spirit of the Seiberg–Witten duality<sup>11</sup>: The theories with strong coupling (i.e., including confinement) are equivalent to the theories with weak coupling (i.e., without confinement) accompanied by magnetic monopoles. We conclude that the Seiberg–Witten duality has its prototype in 1D quantum mechanics.
b) The anyon–oscillator duality is a simple example of a more complicated dyon–oscillator duality<sup>12-23</sup>. The latter connects the isotropic oscillator with charge–dyon bound system (dyon is a hypothetical object which has both the electric and magnetic charge<sup>24</sup>). The passage from an oscillator to a charge–dyon system is realized by non-bijective bilinear transformations<sup>25</sup> (for the mapping of the 1D Coulomb system into the oscillator refer to<sup>26</sup>).
c) The wave function (13) of 1D Coulomb anyon can formally be extended to the region $`\mathrm{}<y<0`$. Such a continuation is an arbitrary-rule operation and we choose the following one. First, still being in the region $`0<y<\mathrm{}`$, we change $`y`$ in the exponent and confluent hypergeometric function by $`|y|`$ and remain unchanged the factor $`y^\nu `$. Then, we extend the expression to the region $`\mathrm{}<y<0`$. These steps allow us to get rid of divergence in the exponent for large negative values of $`y`$ and conserve the normalization condition in $`\mathrm{}<y<\mathrm{}`$ by multiplying the function $`\mathrm{\Phi }_n^{(\nu )}`$ by the factor $`1/\sqrt{2}`$. The obtained wave function $`\overline{\mathrm{\Phi }}_n^{(\nu )}(y)`$ satisfies Eq. (6) in the region $`\mathrm{}<y<\mathrm{}`$ and has the parity $`(1)^\nu `$, i.e. describes the 1D anyon<sup>4</sup>.
d) Eq. (6) for $`\mathrm{}<x<\mathrm{}`$ and $`\nu =0`$ corresponds to the so-called 1D hydrogen atom<sup>27</sup> (for later references see<sup>28</sup>) which has some mysterious properties. For example, the ground state corresponds to an infinite negative value of the energy and the exited levels are double degenerated. The reason is that the potential $`(1/|x|)`$ is singular in 1D space and the system is provided by hidden symmetry<sup>29-31</sup> and supersymmetry<sup>32,33</sup>. As it is follows from (6) and (12), the Calogero–Sutherland potential transforms the 1D hydrogen atom into two modified atoms with the statistical parameter $`\nu =1/4`$ and $`\nu =3/4`$. This transformation leads to the formation of the ground states with a finite energy level and remove the problem of degeneracy (replacement $`nn+\nu `$).
e) It is easily to be convinced that Eq. (4) is identical to the Schrödinger equation with the Hamiltonian
$$\widehat{H}=\frac{1}{2\mu }\left(i\mathrm{}_x\frac{e}{c}A\right)^2\frac{\alpha }{x}\frac{\mathrm{}^2}{2\mu }\frac{\nu (1\nu )}{x^2}$$
(20)
where $`\alpha =e^2`$, $`A=g/x`$, $`g=i\nu \mathrm{}c/e`$. So, we deal with a charged particle moving in the field created by the 1D Coulomb dyon of the electric charge $`e`$ and purely imaginary magnetic charge $`g`$. The Calogero–Sutherland potential gains the meaning of the Goldhaber term typical of theory of magnetic monopoles<sup>34,35</sup>.
Note that the Hamiltonian in (20) is not Hermitian, but it could be transformed into the Hermitian one if we do the following: 1) consider instead of the semiaxis $`x(0,\mathrm{})`$ the axis $`x(\mathrm{},\mathrm{})`$; 2) substitute $`\alpha /x`$ by $`\alpha /|x|`$; 3) introduce the Yang–Dunkl operator<sup>36</sup> $`\widehat{D}=i\mathrm{}_xeA\widehat{R}/c`$ for the Calogero model, where $`\widehat{R}`$ is the reflection operator.
Exactly the same “covariant derivative” appears in the Calogero-like models with exchange interaction<sup>37,38</sup>.
ACKNOWLEDGMENT
We thank our collaborators L. Mardoyan, A. Nersessian, G. Pogosyan and A. Sissakian. Their evident interest in the subject leads us to the writing of present short note.
Also we are grateful to M.S. Plyushchay for valuable remarks on $`(2+1)`$-dimensional anyons and the Yang–Dunkl operator for Calogero models.
The work of Ye. Hakobyan was supported in part by the Russian Foundation for Basic Research, project no. 98-01-00330.
<sup>a)</sup> Electronic mail: yera@thsun1.jinr.ru
<sup>b)</sup> Electronic mail: terant@thsun1.jinr.ru
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warning/0002/physics0002012.html | ar5iv | text | # REFERENCES
Rupture of multiple parallel molecular bonds under dynamic loading
Udo Seifert
Max-Planck-Institut für Kolloid- und Grenzflächenforschung,
Am Mühlenberg 2, 14476 Golm, Germany
## Abstract
Biological adhesion often involves several pairs of specific receptor-ligand molecules. Using rate equations, we study theoretically the rupture of such multiple parallel bonds under dynamic loading assisted by thermal activation. For a simple generic type of cooperativity, both the rupture time and force exhibit several different scaling regimes. The dependence of the rupture force on the number of bonds is predicted to be either linear, like a square root or logarithmic.
PACS: 87.15 By, 82.20 Mj
Introduction. Single molecule force spectroscopy has made it possible to measure the binding strength of a pair of receptor-ligand (“lock-key”) molecules using vesicles , the atomic force apparatus , or optical tweezers as transducers. Thus, the essential constituents mediating biological adhesion have become accessible to quantitative physical experiments . This experimental progress has fostered theoretical studies of the rupture of such pairs under dynamic loading. Thermal activation being a main contributing factor, Kramers-like descriptions of the rupture process with time-dependent potentials show that the rupture strength of such bonds depends on the loading rate . Such behavior has been found experimentally indeed . While unspecific theoretical models of the rupture process reveal generic features, molecular dynamic studies can address the details of the dynamics of the rupture of specific pairs .
Adhesive contact and the rupture thereof often involves not just one but several molecular pairs of the same or different species . The equilibrium properties of the cooperative effects of such specific interactions are well studied both in theory and in experiments . Concerning the dynamics of rupture of such a contact under loading, detailed models for specific problems such as the peeling of a membrane or the rolling of leucocytes in shear flow have been solved numerically to extract a critical tension or shear rate for rupture. However, it is inherently difficult to separate generic dynamical properties from specific ones using such intricate models. As an example for a generic property consider the following question: How does the time and force necessary to break an adhesive contact under dynamic loading depend on the number of bonds initially present?
The present study addresses this question within a simple model that extends work on the dynamic failure of a single bond to that of a whole patch involving several bonds of the same type. Quite generally, two different limiting cases must be distinguished. If the load is primarily concentrated on one bond at a time with relaxation of the load when the first bonds fails and subsequent loading of the next one, the rupture process basically is a sequence of similar single molecule events. The $`N_0`$ bonds initially present then act in series. The rupture time will be $`N_0`$ whereas the force will exhibit a saw tooth-like pattern with a peak given by the rupture strength of a single bond. Such a behavior has been found and modeled in the related case of unfolding of proteins with several identical domains like titin . The main purpose of the present paper is to analyze the other case where the load is distributed (almost) uniformly among several bonds such that these bonds act in parallel. As more and more bonds rupture, the force on the remaining ones increases. This simple type of cooperativity leads to different scaling regimes for the rupture time and rupture force.
FIG.1: Model geometry for the rupture of parallel bonds. Symbols are explained in the main text.
Model. We model the rupture geometry generically as shown in Fig.1. One partner of the bond (“receptor”) is confined to a substrate. The other (“ligand”) is connected by a polymer to a transducer which is connected by an elastic element to a sled being pulled at velocity $`v`$. For simplicity, we model both the elasticity of the transducer and the polymers as Hookean springs with zero rest length and spring constants $`K_t`$ and $`K_p`$, respectively. As long as a bond is intact the corresponding polymer is stretched to an extension $`x_p`$ which we assume to be the same for all intact bonds. The elongation of the transducer from its resting position is $`x_t`$. Force balance on the transducer becomes $`NK_px_p=K_tx_t`$ where $`N`$ is the number of intact bonds. Geometry dictates the time dependence $`x_p(t)+x_t(t)=vt.`$ From these two relations, we find the time-dependent force on an intact bond as
$$F_b(t)=K_px_p(t)=\frac{K_pK_t}{N(t)K_p+K_t}vt.$$
(1)
Following Bell , we assume that the main effect of such a force is to introduce an instantaneous, time-dependent dissociation rate $`k_0(t)`$ according to
$$k_0(t)=k_0\mathrm{exp}[F_b(t)x_b/k_BT],$$
(2)
where $`k_0`$ is the dissociation rate in the absence of a force. The quantity $`x_b`$ is of the order of the distance between the minimum of the binding potential and the barrier and $`k_BT`$ is the product of Boltzmann’s constant and temperature.
We are mainly interested in the case of a soft transducer defined as $`K_t\stackrel{<}{}K_p`$. In this case, eq. (1) shows that the force on a bond is inversely proportional to the number of intact bonds for all $`N(t)`$. Hence, when a bond ruptures, the force on the remaining ones increases. We now discuss two different cases, irreversible and reversible bonds. In the former case, a bond, once ruptured, cannot rebind. Reversible bonds have a non-zero rebinding rate.
Irreversible bonds. Initially $`N(t=0)N_0`$ bonds are present. The rate equation for their time-dependent decrease is
$$_tN=N(t)k_0\mathrm{exp}[F_b(t)x_b/k_BT].$$
(3)
We scale time with the dissociation rate in equilibrium $`k_0`$ according to $`\tau tk_0.`$ The rate equation in the case of a soft transducer then becomes
$$_\tau N=N\mathrm{exp}[(\mu \tau /N]$$
(4)
with the loading parameter
$$\mu K_tx_bv/k_BTk_0.$$
(5)
This simple rate equation seems not to have an analytical solution. However, its scaling behavior can be extracted by the following analysis. With the substitution $`u(\tau )\tau /N`$ one obtains
$$_\tau u=u(1/\tau +1)+u(\mathrm{exp}[\mu u]1)$$
(6)
For small $`\tau `$, $`u(\tau )\tau /N_0`$ and the second term in (6) can therefore be neglected. The solution $`u_1(\tau )`$ of the corresponding equation becomes $`u_1(\tau )=\tau e^\tau /N_0`$ and hence a purely exponential decay for the number of intact bonds, $`N(\tau )=N_0e^\tau `$. This approximation breaks down for $`\tau \stackrel{>}{}\tau _1`$ with $`\tau _1`$ implicitly defined by
$$\mathrm{max}(1/\tau _1,1)=\mathrm{exp}[\mu u_1(\tau _1)]1.$$
(7)
For $`\tau >\tau _1`$, we can then ignore both the first term and the “-1” in the second term of (6). The corresponding equation $`_\tau u=u\mathrm{exp}[\mu u]`$ is solved by
$$E(\mu u_1)E(\mu u)=\tau \tau _1$$
(8)
where $`E(x)_x^{\mathrm{}}𝑑x^{}e^x^{}/x^{}`$ is the exponential integral and $`u_1u_1(\tau _1)`$ is the cross over value of the first solution at the matching point $`\tau _1`$. Hence the time necessary for complete rupture, $`\tau ^{}`$, can be estimated by setting $`u(\tau ^{})=\tau ^{}/N=\mathrm{}`$ which leads to
$$\tau ^{}=\tau _1+\tau _2=\tau _1+E(\mu u_1).$$
(9)
Based on this approximative solution of (6), three sub-regimes can be identified:
(i) $`\mu \stackrel{<}{}1`$: In this case, the exponential decay holds till $`N(\tau )1`$. Physically, the rupture is then complete. In this trivial regime, where the loading is too small to affect the rupture process at all, the time required for rupture is
$$\tau ^{}\mathrm{ln}N_0.$$
(10)
Note that the same result could have been obtained by analyzing the mean time required for the irreversible decay of $`N_0`$ independent bonds under no force.
(ii) $`1\stackrel{<}{}\mu \stackrel{<}{}N_0`$: In this regime, the exponential decay persists till $`\tau _1\mathrm{ln}(N_0/\mu )`$. At this time the number of bonds has reached $`N(\tau _1)\mu `$. The remaining bonds decay according to (8) which leads to an additional time $`\tau _2`$ of order 1 which is small compared to $`\tau _1`$. Hence the whole rupture time in this regime is of order
$$\tau ^{}\mathrm{ln}(N_0/\mu ).$$
(11)
(iii) $`N_0\stackrel{<}{}\mu `$. In this case, the exponential decay applies till $`\tau _1(N_0/\mu )\mathrm{ln}(\mu /N_0)`$. According to (8) the remaining time $`\tau _2(N_0/\mu )/\mathrm{ln}(\mu /N_0)`$ is smaller than $`\tau _1`$. Hence the total rupture time is
$$\tau ^{}(N_0/\mu )\mathrm{ln}(\mu /N_0).$$
(12)
Thus we find for small loading rates that the rupture time is logarithmic in the number of bonds initially present whereas for large loading, this time becomes linear in $`N_0`$. For fixed $`N_0`$ and increasing $`\mu `$, the rupture time first is independent of $`\mu `$. It then decays logarithmically in $`\mu `$ and finally becomes inversely proportional to $`\mu `$.
The force measured by the transducer is given by
$$F_tN(t)K_px_p(t)(k_BT/x_b)\mu \tau .$$
(13)
Thus, the total force experienced by the soft transducer is independent of the number of intact bonds and increases linearly in time. The dimensionless rupture force $`f^{}=\mu \tau ^{}`$ is thus given by
$`f^{}`$ $``$ $`\mu \mathrm{ln}N_0\mathrm{for}\mu \stackrel{<}{}1,`$ (14)
$``$ $`\mu \mathrm{ln}(N_0/\mu )\mathrm{for}1\stackrel{<}{}\mu \stackrel{<}{}N_0,`$ (15)
$``$ $`N_0\mathrm{ln}(\mu /N_0)\mathrm{for}N_0\stackrel{<}{}\mu .`$ (16)
in the three regimes, respectively.
Reversible bonds. So far, we have neglected the possibility that broken bonds can reform. Hence, rupture from a genuine equilibrium situation where bonds form, break, and rebind requires a refined description where we add a term for rebinding. We assume that one species of the receptor/ligand couple is limited to a total number $`N_1`$ with $`N(t)`$ molecules bound and $`N_1N(t)`$ unbound whereas the other species is available in excess. The rate equation becomes
$$_tN=N(t)k_0\mathrm{exp}[\mu \tau /N(t)]+k_f(N_1N(t)),$$
(17)
where we assume for simplicity that the rate $`k_f`$ for bond formation is not affected by the force. Without loading, the equilibrium number of bonds is
$$N_{eq}=\gamma N_1/(1+\gamma )$$
(18)
where $`\gamma k_f/k_0`$. As loading starts, the number of bonds decreases from this equilibrium value. With $`u(\tau )\tau /N`$ as before, we get
$$_\tau u=u(1/\tau +1+\gamma \gamma N_1u/\tau )+u(\mathrm{exp}[\mu u]1)$$
(19)
For $`\mu =0`$, this equation is solved by
$$u_0(\tau )\tau /N_{eq},$$
(20)
which corresponds to the stationary equilibrium distribution. The loading term becomes relevant at a time $`\tau =\tau _1`$ for which
$$(\mathrm{exp}[\mu u_0(\tau _1)/N_1]1)u_0_\tau u_0=1/N_{eq}.$$
(21)
Two cases must then be distinguished:
(i) For $`\mu \stackrel{<}{}N_{eq}`$, $`\tau _1(N_{eq}/\mu )^{1/2}`$. Up to this time, the loading has not significantly affected the number of bonds. The remaining time till all bonds are ruptured can be estimated to be of the same order as $`\tau _1`$ using (19). Hence,
$$\tau ^{}(N_{eq}/\mu )^{1/2}.$$
(22)
In this case, the rupture time increases as a square root of the equilibrium bonds present and decreases as a square root of the loading parameter.
(ii) For $`N_{eq}\stackrel{<}{}\mu `$, $`\tau _1(N_{eq}/\mu )\mathrm{ln}(\mu /N_{eq})`$, with a remaining time of the same order. Hence in this case, we recover the irreversible result (12) with $`N_0`$ replaced by $`N_{eq}`$.
Since in both cases the rupture time $`\tau ^{}\tau _1`$, we get easily for the rupture force
$`f^{}=\mu \tau ^{}`$ $``$ $`(\mu N_{eq})^{1/2}\mathrm{for}\mu \stackrel{<}{}N_{eq},`$ (23)
$``$ $`N_{eq}\mathrm{ln}(\mu /N_{eq})\mathrm{for}N_{eq}\stackrel{<}{}\mu .`$ (24)
Stiff transducer. So far, we have considered the case of a soft transducer $`(K_t\stackrel{<}{}K_p)`$ for which the force on a bond depends on the number of bonds. Another limiting case is a stiff transducer with $`K_t\stackrel{>}{}N_0K_p`$ and $`K_t\stackrel{>}{}N_{eq}K_p`$ for the case of irreversible and reversible rupture, respectively. According to eq. (1), the force on a bond then is (almost) independent of the number of bonds. Hence, the rupture time is only weakly dependent on the number of bonds. An analysis of the corresponding rate equations along similar lines as above shows for the irreversible case two subregimes with
$`\tau ^{}`$ $``$ $`\mathrm{ln}N_0\mathrm{for}\overline{\mu }\stackrel{<}{}1,`$ (25)
$``$ $`\mathrm{ln}\overline{\mu }/\overline{\mu }\mathrm{for}\overline{\mu }\stackrel{>}{}1`$ (26)
with a loading parameter
$$\overline{\mu }K_px_bv/k_BTk_0$$
(27)
dominated by the polymeric stiffness. For the dimensionless maximal force experienced by the transducer during the rupture process one finds
$`f^{}`$ $``$ $`\overline{\mu }N_0\mathrm{for}\overline{\mu }\stackrel{<}{}1,`$ (28)
$``$ $`N_0\mathrm{ln}\overline{\mu }\mathrm{for}\overline{\mu }\stackrel{>}{}1`$ (29)
in the two cases.
Similarly, for a stiff transducer and reversible bonds, one gets
$`\tau ^{}`$ $``$ $`(\mathrm{ln}N_{eq})^{1/2}/\overline{\mu }^{1/2}\mathrm{for}\overline{\mu }\stackrel{<}{}1,`$ (30)
$``$ $`\mathrm{ln}\overline{\mu }/\overline{\mu }\mathrm{for}\overline{\mu }\stackrel{>}{}1`$ (31)
and for the dimensionless maximal force experienced by the transducer
$`f^{}`$ $``$ $`\overline{\mu }^{1/2}N_{eq}\mathrm{for}\overline{\mu }\stackrel{<}{}1,`$ (32)
$``$ $`N_{eq}\mathrm{ln}\overline{\mu }\mathrm{for}\overline{\mu }\stackrel{>}{}1.`$ (33)
Finally, there is a crossover regime for $`N_{0,eq}K_p\stackrel{>}{}K_t\stackrel{>}{}K_p`$, where the pulling starts as in the soft case. As the number of intact bonds decreases towards the value $`\stackrel{~}{N}K_t/K_p`$, the denominator in (1) becomes dominated by $`K_t`$ and the rupture process proceeds as for a stiff transducer. For the reversible case, it turns out that both the rupture time and the rupture force are dominated by the soft part. Hence the results (22,23) apply for all $`K_t\stackrel{<}{}N_{eq}K_p`$. For the irreversible case, analysis of the crossover regime is slightly more involved. The different scaling regimes for rupture time and force are shown in Fig. 2 without explicit derivation.
Concluding perspective. Based on an analysis of rate equations, the comprehensive scaling analysis presented in this paper has revealed several different regimes for the rupture time and force of parallel molecular bonds under dynamic loading. The most distinctive regime is presumably the square root dependence of rupture time and force (22,23) on loading rate and number of bonds derived for reversible bonds under small loading. Such a square root behavior on the loading rate is different from both the irreversible case and the dependence on loading rate for rupture of a single bond or bonds in series. An experimental result showing such an exponent could therefore be taken as a signature of breaking multiple parallel reversible bonds. Of course, it will be important to work with a model system where the number or density of bonds of at least one partner can be controlled in order to extract the dependence of rupture time and force on this crucial quantity.
An obvious theoretical refinement of the present model would be to include fluctuations of the rupture time for individual bonds. Other ramifications can include allowing lateral interactions between the bonds, combining the simplistic Hookean transducer with a membrane patch with its own elasticity, or modeling the rupture process more delicately than done here to name just a few possibilities. It will be interesting to see how robust the scaling regimes derived in this paper will be under such modifications which can effectively lead to scenarios somewhere between the present “in parallel” case and the “in series” case described briefly in the introduction. Finally, it should be clear that in spite of – or rather because of – the progress made in understanding the single bond behavior, the cooperative effects of several bonds under dynamic loading deserve further attention both in theory and in experiment.
Acknowledgments: I thank E. Sackmann for a stimulating discussion and J. Shillcock for a critical reading of the manuscript.
FIG.2: Dynamical phase diagram for (a) the dimensionless rupture time $`\tau ^{}`$ and (b) the dimensionless rupture force $`f^{}`$ as a function of the two loading parameters $`\mu `$ (5) and $`\overline{\mu }`$ (27) in the case of irreversible rupture. In the region $`w`$, the rupture force is given by $`f^{}\mu \mathrm{ln}(N_0\overline{\mu }/\mu )`$. |
warning/0002/astro-ph0002238.html | ar5iv | text | # Large-Scale Sunyaev-Zel’dovich Effect: Measuring Statistical Properties with Multifrequency Maps
## 1. Introduction
It is by now well established that the precision measurements of the cosmic microwave background expected from upcoming experiments, especially MAP and Planck satellite missions, will provide a gold mine of information about the early universe and the fundamental cosmological parameters (e.g., Jungman et al. 1996). These experiments can in fact do so much more. With all-sky maps across the wide range of uncharted frequencies from $`20`$GHz-$`900`$GHz, the secondary science from these missions will arguably be as interesting as the primary science.
In this paper, we examine the prospects for extracting the large-scale properties of the hot intergalactic gas from multifrequency observations of the CMB. Inverse-Compton scattering of CMB photons by hot gas, known as the Sunyaev-Zel’dovich (SZ; Sunyaev & Zel’dovich 1980) effect, leaves a characteristic distortion in the spectrum of the CMB, which fluctuates in the sky with the gas density and temperature. In the Rayleigh-Jeans (RJ) regime, it produces a constant decrement and with only low frequency measurements, the much larger primary anisotropies in the CMB itself obscure the fluctuations on scales greater than a few arcminutes (e.g., Goldberg & Spergel 1999). The upscattering in frequency implies an increment at high frequencies and a null around $`217`$GHz. This behavior provides a potential tool for the separation of SZ effect from other temperature anisotropy contributors.
Since both the SZ spectrum and the CMB spectrum are accurately known, one can expect that foreground removal techniques developed to isolate the primary anisotropies can be reversed to recover the SZ signal in the presence of noise from the primary anisotropies. Galactic and extragalactic foregrounds will be more challenging to remove. Here we use the latest foreground models from Tegmark et al. (1999) that takes into account the fact that imperfect correlations in the foregrounds between frequency channels inhibits our ability to remove them. Using foreground information together with the expected noise properties of individual experiments, one can determine the minimal detectable signal in each experiment and the upper limit achievable in the absence of detection. Experiments with sufficient signal-to-noise can extract precision measurements for the power spectrum and higher order statistics such as the skewness. Ultimately, they can provide detailed maps of the large-angle SZ effect.
To assess the prospects for an actual detection, we must model the SZ signal itself. The SZ effect is now routinely imaged in massive galaxy clusters (e.g., Carlstrom et al 1996; Jones et al 1993), where the temperature of the scattering medium can reach as high as 10 keV, producing temperature changes in the CMB of order 1 mK at RJ wavelengths. The possibility for detection of massive clusters in CMB satellite data has already been discussed in several studies (e.g., Aghanim et al. , Haehnelt & Tegmark 1996, Pointecouteau et al 1998). Here, however, we are interested in the SZ effect produced by large-scale structure in the general intergalactic medium (IGM) where the gas is expected to be at $`1`$keV in mild overdensities, leading to CMB contributions in the $`\mu `$K range.
It is now widely believed that at least $``$ 50% of the present day baryons, when compared to the total baryon budget from big bang nucleosynthesis, are present in gas associated with hot large-scale structure which has remained undetected given its temperature and clustering properties (e.g., Fukugita et al 1998; Cen & Ostriker 1999; Pen 1999). Recently, Scharf et al (2000) has provided a tentative detection of X-ray emission from a large-scale filament in one of the deep ROSAT PSPC fields; previous attempts yielding upper limits are described in Kull & Böhringer (1999) and Briel & Henry (1995). These results are consistent with current predictions for the X-ray surface brightness based on numerical simulations (e.g., Cen et al 1995). Pen (1999) argued that non-gravitational heating of the gas to $`1`$keV is required to evade bounds from the soft X-ray background. These results suggest that the X-ray emission from this gas may be detectable in the near future with wide-field observations with Chandra X-ray Observatory<sup>1</sup><sup>1</sup>1http://asc.harvard.edu and X-ray Multiple Mirror Mission<sup>2</sup><sup>2</sup>2http://astro.estec.esa.nl/XMM.
On the theoretical front, hydrodynamic simulations of the SZ effect continue to improve (da Silva et al. 1999; Refregier et al. 1999; Seljak et al. 2000). As a consensus from these simulations of basic properties such as the opacity weighted gas temperature and average Compton distortion is still lacking, we will base our assessment of the detectability of the large-scale SZ effect on a simple parameterization of the effect, based on a gas pressure bias model (Refregier et al. 1999), crudely calibrated with the recent hydrodynamic simulations. We employ perturbation theory, non-linear scaling relations, and N-body simulations for the dark matter to assess the statistical properties of the signal. Properly calibrated, these techniques can complement hydrodynamic simulations by extending their dynamic range and sampling volume. Currently, they should simply be taken as order of magnitude estimates of the potential signal.
Throughout this paper, we will take an adiabatic cold dark matter (CDM) model as our fiducial cosmology. We assume cosmological parameters $`\mathrm{\Omega }_c=0.30`$ for the cold dark matter density, $`\mathrm{\Omega }_b=0.05`$ for the baryon density, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.65`$ for the cosmological constant, $`h=0.65`$ for the dimensionless Hubble constant and a COBE-normalized scale invariant spectrum of primordial fluctuations (Bunn & White 1997).
The layout of the paper is as follows. In § 2, we describe the foreground and primary anisotropy removal method and assess their efficacy for upcoming CMB experiments. In § 3, we detail the bias model for the SZ effect and calculate through perturbation theory, analytic approximations and numerical simulations, the low order statistics of the SZ effect: its power spectrum, skewness and bispectrum. In § 4, having estimated the noise and the signal, we assess the prospects for measuring these low order statistics in upcoming experiments. We conclude in § 5 with a discussion of our results.
## 2. Modeling the CMB and Foreground Noise
The main obstacle for the detection of the SZ effect from large-scale structure for angular scales above a few arcminutes is the CMB itself. Here the primary anisotropies dominate the SZ effect for frequencies near and below the peak in the CMB spectrum (see Fig. 1). Fortunately, the known frequency dependence and statistical properties of primary anisotropies allows for extremely effective subtraction of their contribution (e.g., Hobson et al 1998; Bouchet & Gispert 1999). In particular, primary anisotropies obey Gaussian statistics and follow the blackbody spectrum precisely.
Perhaps more worrying are the galactic and extragalactic foregrounds, some of which are expected to to be at least comparable to the SZ signal in each frequency band. These foregrounds typically have spatial and/or temporal variations in their frequency dependence leading to imperfect correlations between their contributions in different frequency bands. We attempt here to provide as realistic an estimate as possible of the prospects for CMB and foreground removal, given our incomplete understanding of the foregrounds.
### 2.1. Foreground Model and Removal
We use the “MID” foreground model of Tegmark et al. (1999) and adapt the subtraction techniques found there for the purpose of extracting the SZ signal. The assumed level of the foreground signal in the power spectrum for three fiducial frequencies is shown in Fig. 1.
The foreground model is defined in terms of the covariance between the multipole moments at different frequency bands<sup>3</sup><sup>3</sup>3A potential caveat for this type of modeling is that it assumes the foregrounds are statistically isotropic whereas we know that the presence of the Galaxy violates this assumption at least for the low order multipoles. We assume that $`1f_{\mathrm{sky}}0.35`$ of the sky is lost to this assumption even with an all-sky experiment.
$$a_{l^{}m^{}}^\mathrm{f}(\nu ^{})a_{lm}^\mathrm{f}(\nu )=C_l^\mathrm{f}(\nu ^{},\nu )\delta _{ll^{}}\delta _{mm^{}},$$
(1)
in thermodynamic temperature units as set by the CMB blackbody. In this section, we will speak of the primary anisotropies and detector noise simply as other foregrounds with very special properties:
$`C_l^{\mathrm{CMB}}(\nu ^{},\nu )`$ $`=`$ $`C_l,`$
$`C_l^{\mathrm{noise}}(\nu ^{},\nu )`$ $`=`$ $`8\mathrm{ln}2\theta (\nu )^2e^{\theta ^2(\nu )l(l+1)}\left({\displaystyle \frac{\mathrm{\Delta }T}{T}}\right)^2|_{\mathrm{noise}}\delta _{\nu ,\nu ^{}}.`$
The FWHM$`=\sqrt{8\mathrm{ln}2}\theta `$ and noise specifications of the Boomerang, MAP and Planck frequency channels are given in Tab. 1. True foregrounds generally fall in between these extremes of perfect and no frequency correlation.
The difference between extracting the SZ signal and the primary signal is simply that one performs the subtraction referenced to the SZ frequency dependence
$$s(\nu )=2\frac{x}{2}\mathrm{coth}\frac{x}{2},$$
(3)
where $`x=h\nu /kT_{\mathrm{cmb}}\nu /56.8`$GHz. Note that in the RJ limit $`s(\nu )1`$ such that
$$C_l^{\mathrm{SZ}}(\nu ,\nu ^{})=s(\nu )s(\nu ^{})C_l^{\mathrm{SZ}}$$
(4)
where $`C_l^{\mathrm{SZ}}`$ is the SZ power spectrum in the RJ limit.
Consider an arbitrary linear combination of the channels,
$$b=\underset{\nu _i}{}\frac{1}{s(\nu _i)}w(\nu _i)a(\nu _i).$$
(5)
Since the subtraction is done multipole by multipole, we have temporarily suppressed the multipole index. The covariance of this quantity is
$$b^2=C^{\mathrm{SZ}}[\underset{\nu _i}{}w(\nu _i)]^2+\underset{\nu _i,\nu _j}{}w(\nu _i)w(\nu _j)\underset{\mathrm{f}}{}\frac{C^\mathrm{f}(\nu _i,\nu _j)}{s(\nu _i)s(\nu _j)}.$$
(6)
Minimizing the variance contributed by the foregrounds subject to the constraint that the SZ estimation be unbiased, we obtain
$$\underset{\nu _i}{}w(\nu _j)\underset{\mathrm{f}}{}\frac{C^\mathrm{f}(\nu _i,\nu _j)}{s(\nu _i)s(\nu _j)}=\mathrm{const}.$$
(7)
Defining the scaled foreground covariance matrix as
$`\stackrel{~}{C}(\nu _i,\nu _j)`$ $`=`$ $`{\displaystyle \underset{\mathrm{f}}{}}{\displaystyle \frac{C^\mathrm{f}(\nu _i,\nu _j)}{s(\nu _i)s(\nu _j)}},`$ (8)
$``$ $`{\displaystyle \underset{\mathrm{f}}{}}\stackrel{~}{C}^\mathrm{f}(\nu _i,\nu _j),`$
we solve for the weights that minimize the noise variance
$$𝐰\stackrel{~}{𝐂}^1𝐞,$$
(9)
where $`𝐞`$ is the vector of all ones $`e(\nu _i)=1`$. Finally we normalize the sum of the weights to unity $`w(\nu _i)=1`$ to obtain an unbiased estimator. Our approach is same as minimizing the foreground variance subject to the constraint that the recovered multipole is an unbiased estimate of the true SZ signal. As each channel is rescaled such that SZ signal corresponds to the RJ level, the weights sum to unity.
The total residual noise variance in the map from the foregrounds per multipole is then
$$N_l=𝐰_l^t\stackrel{~}{𝐂}_l𝐰_l,$$
(10)
and from each foreground component
$$N_l^\mathrm{f}=𝐰_l^t\stackrel{~}{𝐂}_l^\mathrm{f}𝐰_l.$$
(11)
Note that the residual noise in the map is independent of assumptions about the SZ signal including whether it is Gaussian or not. However if the foregrounds themselves are non-Gaussian, then this technique only minimizes the variance and may not be optimal for recovery of non-Gaussian features in the SZ map itself. Bouchet et al (1995) have shown that similar techniques are quite effective even when confronted with non-Gaussian foregrounds. This is a potential caveat especially for cases in which the residual noise is not dominated by the primary anisotropies or detector noise. We shall discuss methods to alleviate this concern in the next section.
### 2.2. Detection Threshold
The residual noise sets the detection threshold for the SZ effect for a given experiment. In Fig. 1, we show the rms of the residual noise after foreground subtraction for the Boomerang, MAP and Planck experiments assuming the “MID” foreground model from Tegmark et al. (1999). With the Boomerang and Planck channels, elimination of the primary anisotropies is excellent up to the beam scale where detector noise dominates. As expected, the MAP channels, which are all on the RJ side of the spectrum, do not allow good elimination of the primary anisotropies.
It is important not to assume that the foregrounds are perfectly correlated in frequency, which is the usual assumption in the literature (Hobson et al 1998; Bouchet & Gispert 1999). There are two types of errors incurred by doing so. The first is that one underpredicts the amount of residual noise in the SZ map (see Fig. 2). The second is that if one calculates the optimal weights in equation (9) based on this assumption the actual residual noise increases. For Planck it can actually increase the noise beyond the level in which it appears in the $`100`$GHz maps with no foreground subtraction at all. The reason is that the cleaning algorithm then erroneously uses the contaminated high and low frequency channels to subtract out the small foreground contamination in the central channels. In Planck, the difference between the predicted and actual rms noise from falsely assuming perfect frequency coherence can be more than two orders of magnitude.
For Boomerang and Planck, the largest residual noise component, aside from detector noise, is dust emission and is sufficiently large that one might worry that current uncertainties in our knowledge of the foreground model may affect the implications for the detection of the SZ effect. It is therefore important to explore variations on our fiducial foreground model.
Multiplying the foreground rms amplitudes uniformly by a factor of 2 (and hence the power by a factor of 4), produces less than a factor of 2 increase in the residual noise rms as shown in Fig. 2. Likewise, as discussed in Tegmark et al. (1999), minor variations in the frequency coherence do not effect the residual noise much in spite of the fact that it is crucial not to assume perfect correlation. We conclude that uncertainties in the properties of currently known foregrounds are unlikely to change our conclusions qualitatively. There is however always the possibility that some foreground that does not appear in the currently-measured frequency bands will affect our results.
The fact that the residual dust contributions are comparable to those of the detector noise for Boomerang and Planck is problematic for another reason. Since the algorithm minimizes to total residual variance, it attempts to keep these two main contributors roughly comparable. However the dust will clearly be non-Gaussian to some extent and one may prefer instead to trade more residual detector noise for dust contamination. One can modify the subtraction algorithm to account for this by artificially increasing the rms amplitude of the dust when calculating the weights in equation (9) while using the real amplitude in calculating the residual noise in equation (11). For example we have explored increasing the amplitude by a factor of 4 (power by 16) for the weights. The result is an almost negligible increase in total residual noise rms but an improvement in dust rejection by a factor of 3-4 in rms. For Planck this brings the ratio of dust to total rms to $`10\%`$ and recall that the noise adds in quadrature so that the total dust contribution is really $`1\%`$ of the total. This more conservative approach is thus advisable but since it leaves the total residual noise rms essentially unchanged, we will adopt the minimum variance noise to estimate the detection threshold.
Fig. 1 directly tells us the detection threshold per $`(l,m)`$ multipole moment. Since the SZ signal is likely to have a smooth power spectrum in $`l`$, one can average over bands in $`l`$ to beat down the residual noise. Assuming Gaussian-statistics, the residual noise variance $`2N_l^2`$ for the power spectrum estimate is then given by
$$N_l^2|_{\mathrm{band}}=f_{\mathrm{sky}}\underset{l_{\mathrm{band}}}{}(2l+1)N_l^2,$$
(12)
where $`f_{\mathrm{sky}}`$ accounts for the reduction of the number of independent modes due to the fraction of sky covered. The result for the three experiments is shown in Fig. 3. In the absence of a detection, they can be interpreted as the optimal 1 $`\sigma `$ upper limits on SZ bandpowers achievable by the experiment. Boomerang and MAP can place upper limits on the SZ signal in the interesting $`\mu `$K regime whereas Planck can detect signals well below a $`\mu `$K.
This noise averaging procedure in principle implicitly assumes that the statistical properties of the residual noise, and by implication the full covariance matrix of the other foregrounds, is precisely known. In reality, they too must be estimated from the multifrequency data itself through either through the subtraction techniques discussed here or by direct modeling of the foregrounds in the maps. Tegmark et al. (1999) found that direct modeling of the foregrounds with hundreds of fitted parameters did not appreciably degrade our ability to extract the properties of the primary anisotropies. The main source of variance there was the cosmic variance of the primary anisotropies themselves whose properties are precisely known. Similarly here the main source of residual variance is either the primary anisotropies (for MAP) or detector noise (for Boomerang and Planck) and their statistical properties may safely be considered known.
## 3. Modeling the SZ Signal
In order to estimate how well the statistical properties of the SZ effect might be recovered with multifrequency CMB maps, we need to model the large-angle SZ effect itself. The current state-of-the-art in hydrodynamic simulations (da Silva et al. 1999; Refregier et al. 1999; Seljak et al. 2000) has reached a qualitative but not quantitative consensus on the statistical properties of the SZ effect. In addition, questions as to the heating of the gas from non-gravitational sources may even change the results qualitatively (Pen 1999). Hydrodynamic simulations are also severely limited in the dynamic range and volume sampled.
Given the current state of affairs, we believe it is useful to explore a parameterized model of the effect whose consequences are simple to calculate and which may be calibrated against hydrodynamic simulations as they continue to improve.
### 3.1. Bias Prescription
In general, the SZ temperature fluctuation $`\mathrm{\Theta }=\mathrm{\Delta }T/T`$ is given by the opacity weighted integrated pressure fluctuation along the line of sight:
$`\mathrm{\Theta }^{\mathrm{SZ}}(\widehat{𝐧},\nu )`$ $`=`$ $`2s(\nu ){\displaystyle _0^{r_0}}𝑑r\dot{\tau }\pi (r,\widehat{𝐧}r),`$ (13)
$`r`$ is the the comoving distance, $`\tau `$ is the Thomson optical depth, overdots are derivatives with respect to $`r`$ and the dimensionless electron pressure fluctuation is
$$\pi =\delta p_e/\rho _e.$$
(14)
One needs to model the statistical properties of $`\pi `$, in particular its power spectrum and bispectrum
$`\pi (𝐤)^{}\pi (𝐤^{})`$ $`=`$ $`(2\pi )^3\delta ^\mathrm{D}(𝐤𝐤^{})P_\pi (k),`$ (15)
$`\pi (𝐤)\pi (𝐤^{})\pi (𝐤^{\prime \prime })`$ $`=`$ $`(2\pi )^3\delta ^\mathrm{D}(𝐤+𝐤^{}+𝐤^{\prime \prime })B_\pi (k,k^{},k^{\prime \prime }),`$
as a function of lookback time or distance $`r`$. In principle we also need the unequal time correlators, but in practice these do not play a role as we shall see.
By analogy to the familiar case of galaxy clustering, it is reasonable to suppose that the pressure fluctuations depend locally on the dark matter density and hence are biased tracers of the dark matter density in the linear regime (Goldberg & Spergel 1999). Hence the statistical properties follow from those of the dark matter distribution
$`P_\pi (k;r)`$ $``$ $`b_\pi (r)^2P_\delta (k;r),`$
$`B_\pi (k,k^{},k^{\prime \prime };r)`$ $``$ $`b_\pi (r)^3B_\delta (k,k^{},k^{\prime \prime };r).`$ (16)
We have restored the time dependence since the bias will be time dependent even in the linear regime and must be extracted from simulations. In general, the bias parameter for the power spectrum and bispectrum need not be the same even in the linear regime since the bispectrum automatically involves higher order corrections (Fry & Gaztanaga 1993). For estimation purposes here we will take them to be equal.
Following Goldberg & Spergel (1999), we chose the form
$`b_\pi (r)=b_\pi (0)/(1+z),`$ (17)
as motivated by findings that the average gas temperature drops off roughly by this factor. We normalize the value of the bias parameter today by comparison with recent hydrodynamic simulations. It is conceptually useful to separate the bias into two factors:
$$b_\pi (0)=\frac{k_BT_e(0)}{m_ec^2}b_\delta ,$$
(18)
i.e. an opacity-weighted average temperature and a bias parameter for the gas density at that temperature. In Refregier et al. (1999), for our fiducial $`\mathrm{\Lambda }`$CDM cosmology, the bias $`b_\delta `$ was found to be $``$ 8 to 9, while in Seljak et al. (2000) it was found to be in the range $``$ of $`3`$ to $`4`$. In both these papers, $`T_e(0)`$ 0.3 to 0.4 keV; these values are lower than the $``$ 1 keV found by Cen & Ostriker (1999) using hydrodynamical simulations with feedback effects. As a compromise between these results, we take $`T_e(0)=0.5`$keV and $`b_\delta =4`$, which corresponds to
$$b_\pi (0)=0.0039.$$
(19)
Note that this is a factor of 2 lower than used in Goldberg & Spergel (1999) and Cooray & Hu (1999).
Needless to say, the resulting predictions should be taken as order-of-magnitude estimates only. As simulations improve, one can expect better values for the bias today and a more detailed modeling of its redshift and perhaps even scale dependence.
### 3.2. Multipole Moments
The multipole moments of the SZ effect under these simplifying assumptions can then be expressed as a weighted projection of the density field (Cooray & Hu 1999):
$`a_{lm}^{\mathrm{SZ}}(0)`$ $``$ $`{\displaystyle 𝑑\widehat{𝐧}Y_l^m(\widehat{𝐧})\mathrm{\Theta }^{\mathrm{SZ}}(\widehat{𝐧},0)}`$ (20)
$``$ $`i^l{\displaystyle }{\displaystyle \frac{d^3𝐤}{2\pi ^2}}\delta (𝐤,r_l)I_l^{\mathrm{SZ}}(k)Y_l^m{}_{}{}^{}(\widehat{𝐤}),`$
where
$`I_l^{\mathrm{SZ}}(k)`$ $``$ $`W^{\mathrm{SZ}}(r_l)\sqrt{{\displaystyle \frac{\pi }{2l}}}{\displaystyle \frac{1}{k}}F_l(k),`$
$`W^{\mathrm{SZ}}(r)`$ $`=`$ $`2b_\pi (r)\dot{\tau },`$ (21)
in the Limber approximation and (Hu 2000a)
$`r_l`$ $`=`$ $`\mathrm{\Omega }_K^{1/2}H_0^1\mathrm{sinh}^1(\mathrm{\Omega }_K^{1/2}H_0l/k),`$
$`F_l`$ $`=`$ $`(1+\mathrm{\Omega }_KH_0^2l^2/k^2)^{1/4}.`$ (22)
The quantities take on a simple forms for a flat universe: $`r_ll/k`$ and $`F_l(k)1`$. The Limber approximation breaks down for $`l\text{ }<\text{ }50`$ but is sufficient for our purposes.
### 3.3. Power Spectrum
The power spectrum of the SZ effect in this simplified model follows from equation (20),
$`C_l^{\mathrm{SZ}}`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle \frac{dk}{k}k^3P_\delta (k;r_l)[I_l^{\mathrm{SZ}}(k)]^2},`$ (23)
$``$ $`{\displaystyle _0^{r_0}}𝑑r{\displaystyle \frac{[W^{\mathrm{SZ}}(r)]^2}{d_A^2}}P_\delta (l/d_A;r),`$
In the second line we have transformed the integration variable under the Limber correspondence: $`k=l/d_A`$ and
$$\frac{dk}{k}F_l^2\mathrm{}=\frac{dr}{d_A}\mathrm{}.$$
(24)
We see that to go from the flat to curved cosmologies in the Limber approximation one simply replaces the radial distance with the angular diameter distance in the integrand.
In evaluating the SZ power spectrum, we have extended the SZ model to the non-linear regime by using the scaling formulae for the nonlinear dark matter power spectrum of Peacock & Dodds (1996). However, modeling the SZ effect with a scale-independent bias factor will clearly break down deep in the non-linear regime. Refregier et al. (1999) have shown that it is a reasonable approximation in the weakly non-linear regime (overdensities $`\text{ }<\text{ }10`$) for $`z\text{ }<\text{ }1`$ but can be in serious error outside of this range. As the weakly non-linear regime is the one of interest for anisotropies at $`l\text{ }<\text{ }1000`$, we will use this approximation to test the effects of non-linearities. The predicted power spectrum in our fiducial model is shown in Fig. 4.
### 3.4. Bispectrum
The bispectrum of the SZ effect also follows from expression (20)
$`B_{l_1l_2l_3}^{m_1m_2m_3}`$ $``$ $`a_{l_1m_1}a_{l_2m_2}a_{l_3m_3}`$
$`=`$ $`\left[{\displaystyle \underset{j=1}{\overset{3}{}}}i^{l_j}{\displaystyle \frac{d^3k_j}{2\pi ^2}I_{l_j}^{\mathrm{SZ}}(k_j)Y_l^m(\widehat{𝐤}_j)}\right]`$
$`(2\pi )^3\delta ^\mathrm{D}(𝐤_1+𝐤_2+𝐤_3)B_\delta (k_1,k_2,k_3).`$
Here the density bispectrum should be understood as arising from the full unequal time correlator
$$\delta (𝐤_1;r_1)\delta (𝐤_2;r_2)\delta (𝐤_3;r_3),$$
(25)
where the temporal coordinate, which we temporarily suppress, is evaluated in the Limber approximation (22).
To further simplify this expression, we expand the delta function
$$\delta ^\mathrm{D}(𝐤_1+𝐤_2+𝐤_3)=\frac{1}{(2\pi )^3}e^{i(𝐤_1+𝐤_2+𝐤_3)\widehat{𝐧}r}d^3x,$$
(26)
and employ the Rayleigh expansion
$$e^{i𝐤\widehat{𝐧}r}=4\pi \underset{lm}{}i^lj_l(kr)Y_l^m(\widehat{𝐤})Y_l^m(\widehat{𝐧}).$$
(27)
We have assumed here a flat universe to simplify the derivation; as we have seen in the last section, we can promote the final result to a curved universe by replacing radial distances with angular diameter distances.
With these relations, the angular integral over the directions of $`𝐤_j`$ collapse to give
$`B_{l_1l_2l_3}^{m_1m_2m_3}`$ $`=`$ $`{\displaystyle r^2𝑑r\left[\underset{j=1}{\overset{3}{}}\frac{2}{\pi }k_j^2𝑑k_jI_{l_j}^{\mathrm{SZ}}(k_j)j_{l_j}(k_jr)\right]}`$ (28)
$`\times B(k_1,k_2,k_3)G_{l_1l_2l_3}^{m_1m_2m_3},`$
where the Gaunt integral is
$`G_{l_1l_2l_3}^{m_1m_2m_3}`$ $``$ $`{\displaystyle 𝑑\widehat{𝐧}Y_{l_1}^{m_1}Y_{l_2}^{m_2}Y_{l_3}^{m_3}}`$ (34)
$`=`$ $`\sqrt{{\displaystyle \frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}}`$
$`\times \left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right).`$
Here, the quantities in parentheses are the Wigner-3$`j`$ symbols whose properties are described in Appendix A of Cooray & Hu (1999). The integrals over the Bessel functions can again be done in the Limber approximation leaving
$`B_{l_1l_2l_3}^{m_1m_2m_3}`$ $`=`$ $`G_{l_1l_2l_3}^{m_1m_2m_3}{\displaystyle 𝑑r\frac{[W^{\mathrm{SZ}}(r)]^3}{r^4}B_\delta (\frac{l_1}{r},\frac{l_2}{r},\frac{l_3}{r};r)},`$
Note that only equal time contributions contribute in the Limber approximation.
We can promote this result to a curved universe by replacing radial distances with angular diameter distances
$`B_{l_1l_2l_3}^{m_1m_2m_3}`$ $`=`$ $`G_{l_1l_2l_3}^{m_1m_2m_3}{\displaystyle 𝑑r\frac{[W^{\mathrm{SZ}}(r)]^3}{d_A^4}B_\delta (\frac{l_1}{d_A},\frac{l_2}{d_A},\frac{l_3}{d_A};r)}.`$
Finally, we can introduce the angular averaged bispectrum as
$`B_{l_1l_2l_3}={\displaystyle \underset{m_1m_2m_3}{}}\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right)B_{l_1l_2l_3}^{m_1m_2m_3},`$ (37)
to obtain the final result
$`B_{l_1l_2l_3}`$ $`=`$ $`\sqrt{{\displaystyle \frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)`$ (41)
$`\times {\displaystyle }dr{\displaystyle \frac{[W^{\mathrm{SZ}}(r)]^3}{d_A^4}}B_\delta ({\displaystyle \frac{l_1}{d_A}},{\displaystyle \frac{l_2}{d_A}},{\displaystyle \frac{l_3}{d_A}};r).`$
One can alternately derive this relation by taking a flat-sky approach and using the general relation between the flat-sky and all-sky bispectra (see Appendix C, Hu 2000b).
Equation (41) gives the SZ angular bispectrum in terms of the underlying density bispectrum. In second order perturbation theory, the density bispectrum is in turn given by
$`B_\delta (k_1,k_2,k_3;r)`$ $`=`$ $`F_2(𝐤_1,𝐤_2)P_\delta (k_1;r)P_\delta (k_2;r)`$ (42)
$`+5\mathrm{perm}.,`$
where
$`F_2(𝐤_1,𝐤_2)={\displaystyle \frac{5}{7}}+{\displaystyle \frac{𝐤_1𝐤_2}{k_2^2}}+{\displaystyle \frac{2}{7}}{\displaystyle \frac{(𝐤_1𝐤_2)^2}{k_1^2k_2^2}}.`$ (43)
Unfortunately, there exists no accurate fitting formula for the bispectrum of the density field in the mildly non-linear regime; we will employ simulations in §3.6 to address this regime. In the deeply non-linear regime, the density field obeys the hierarchical ansatz
$`B_\delta (k_1,k_2,k_3;r)={\displaystyle \frac{Q_3}{2}}[P(k_1;r)P(k_2;r)+5\mathrm{perm}.],`$ (44)
where the power spectra are given by the non-linear scaling of Peacock & Dodds (1996). Scoccimarro & Frieman (1999) find that for power law power spectrum
$`Q_3(n)=[42^n]/[1+2^{n+1}].`$ (45)
Hui (1999) suggests that for a general power spectrum one should replace $`n`$ with the local linear power spectral index at $`(k_1+k_2+k_3)/3`$.
### 3.5. Skewness
The simplest aspect of the bispectrum that can be measured is the third moment of the map smoothed on some scale with a window $`W(\sigma )`$
$`\mathrm{\Theta }^3(\widehat{𝐧};\sigma )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{l_1l_2l_3}{}}\sqrt{{\displaystyle \frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}}`$ (48)
$`\times \left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)B_{l_1l_2l_3}W_{l_1}(\sigma )W_{l_2}(\sigma )W_{l_3}(\sigma ),`$
where $`W_l`$ are the multipole moments (or Fourier transform in a flat-sky approximation) of the window. For simplicity, we will choose windows which are either top hats in real or multipole space.
It is useful to define the skewness parameter
$$S_3(\sigma )=\frac{\mathrm{\Theta }^3(\widehat{𝐧};\sigma )}{\mathrm{\Theta }^2(\widehat{𝐧};\sigma )^2},$$
(49)
where the second moment is that of the SZ signal
$$\mathrm{\Theta }^2(\widehat{𝐧};\sigma )=\frac{1}{4\pi }\underset{l}{}(2l+1)C_l^{\mathrm{SZ}}W_l^2(\sigma ).$$
(50)
The skewness in our fiducial model is shown for both the perturbation theory and HEPT predictions in Fig. 5.
Since the density bispectrum in both the perturbative and non-linear regime scale as $`[P_\delta (k)]^2`$, the amplitude of the underlying density fluctuations roughly scale out of $`S_3`$. However, the pressure bias $`b_\pi `$ does not: $`S_3b_\pi ^1`$. $`S_3`$ thus provides an observable handle on the bias. This general point applies even if the bias is non-linear although its interpretation will be not be as straightforward (see Fry & Gaztanaga 1993 and Mo, Jing & White 1997 for its application in galaxy biasing).
### 3.6. Numerical Simulations
Since we are interested in the properties of the SZ effect in the weakly-nonlinear regime, cosmological simulations are required to recover the complete statistical properties of the signal and calibrate semi-analytic approaches for its low-order statistics. The simplified SZ model employed in this paper has the virtue that it is easy to simulate as it requires only dark matter and not the gas to model. Its main drawback of course is that results must be taken with a grain of salt due to missing physics.
The realism of the basic approach can be improved by better calibrating the bias model against hydrodynamic simulations. One can envision going beyond the simple redshift dependent bias approach taken here to include scale dependence and stochasticity. Even accounting for these additional complications, simple dark matter simulations can continue to complement full hydrodynamic simulations. Hydrodynamic simulations will always be more limited in dynamic range and sampling volume. Indeed, the current state of the art is limited a handful of realizations across one order of magnitude in physical scale (Refregier et al. 1999; Seljak et al. 2000). A single simulation is then “stacked” on the line-of-sight. Given the range of redshifts at which the SZ effect contributes, the simulation volume is traced many times for each line-of-sight. Moreover, the angular resolution decreases monotonically as one approaches the origin at $`z=0`$.
The reduction in dynamic range due to the angular projection is a serious but not unfamiliar problem in cosmology. It occurs whenever the kernel for the projection spans cosmological distances. White & Hu (1999) introduced a technique of tiling multiple particle-mesh simulations which telescope along the line of sight to maintain a fixed angular resolution for the analogous problem in weak lensing. This also avoids the problem of over-representing the filamentary structure of the map noted by Refregier et al. (1999).
We refer the reader to White & Hu (1999) for details of the approach and tests of the method. The simulation all have a $`256^3`$ mesh with $`256^2`$ lines of sight for the ray tracing on a $`6^{}\times 6^{}`$ field. Other relevant parameters are given in Tab. 2: the box size $`L_{\mathrm{box}}`$, the number of particles $`N_{\mathrm{part}}`$, the number of simulations run $`N_{\mathrm{sim}}`$, the number of tiles of the given box size used $`N_{\mathrm{tile}}`$, the maximum redshift to which a given box is used, and the angular resolution of the mesh the maximum and minimum redshift used $`\theta _{\mathrm{mesh}}`$. Note that we cannot shrink the box size along the line-of-sight indefinitely since the fundamental mode of the box must be in the linear regime to provide accurate evolution. This implies that we lose angular resolution near the origin where a fixed physical scale subtends a large angle on the sky. Furthermore at the higher redshift the number of particles must be increased to eliminate shot noise from the initial conditions. Nonetheless, the tiling technique does a good job of maintaining angular resolution at all but the lowest redshifts.
We construct 500 SZ maps from random combinations of the tiles in Tab. 1 for our statistical analysis; one realization is shown in Fig. 6. The average power spectrum is shown in Fig. 4 (top panel) and compared with the linear perturbation theory prediction and the non-linear scaling relation of Peacock & Dodds (1996). We have tested that the deficit of power at the low multipoles is an artifact of the finite field-of-view through monte-carlo realizations of the predicted power spectrum. The roll-off at high multipoles is due to the spatial resolution in the simulations. This also explains the $`10\%`$ deficit at intermediate scales which comes from highly non-linear structure close to the origin. Agreement is restored if one eliminates contributions from overdensities $`>10`$ in the predictions. Since our SZ model is at best valid in the weakly non-linear regime, these contributions should not be included anyway.
Fig. 5 (top panel) shows the results for the skewness in the simulations compared with the second order perturbation theory and HEPT predictions. The agreement here is worse, but is still sufficient for our purposes, given the crudeness of the underlying model for the SZ effect itself.
We can address sample variance questions from the scatter of the results in the individual realizations. Sampling errors for the power spectrum and skewness are shown in the bottom panels of Fig. 4 and 5 respectively. Since these are for individual $`6^{}\times 6^{}`$ planes, they should be scaled by $`0.03f_{\mathrm{sky}}^{1/2}`$ for a given experiment. Sampling errors are one source of noise that we will include in the signal-to-noise calculations in the next section.
## 4. Estimating the Signal-to-Noise
With the SZ signal estimated from the simple bias model of §3 and residual noise calculated from the foreground model and subtraction techniques of §2, we can now estimate the signal-to-noise for the detection of the SZ effect. In Fig. 7, we illustrate the foreground subtraction technique on simulated Planck maps. The signal-to-noise in the maps is of order one for features spanning tens of arcminutes. We shall here show that this level of signal-to-noise is more than sufficient for the purpose of extracting measurements of the low order statistics of the SZ signal.
### 4.1. Power Spectrum
The signal-to-noise in the power spectrum per multipole $`(l,m)`$ mode is simply
$$\left(\frac{S}{N}\right)_{lm}^2=\frac{1}{2}\left(\frac{C_l^{\mathrm{SZ}}}{C_l^{\mathrm{tot}}}\right)^2.$$
(51)
Here, $`C_l^{\mathrm{tot}}`$ is the power spectrum of all contributions in the SZ map,
$`C_l^{\mathrm{tot}}=C_l^{\mathrm{SZ}}+N_l,`$ (52)
where recall that the residual noise $`N_l`$ was defined in equation (10) and includes contributions from detector noise.
Assuming Gaussian statistics for the signal and noise, each mode is independent so that the total signal-to-noise is the quadrature sum
$$\left(\frac{S}{N}\right)^2=\frac{f_{\mathrm{sky}}}{2}\underset{l}{}(2l+1)\left(\frac{C_l^{\mathrm{SZ}}}{C_l^{\mathrm{tot}}}\right)^2.$$
(53)
This quantity gives the variance of the total power measurement in the SZ effect, including sample variance. $`\mathrm{S}/\mathrm{N}1`$ means that one has a precise measurement of the power spectrum not simply a highly significant detection. In Fig. 8, for the Boomerang, MAP and Planck experiments as a function of the maximum $`l`$ mode included in the sum. We also show the ultimate limit of perfect foreground and noise removal where $`C_l^{\mathrm{tot}}=C_l^{\mathrm{SZ}}`$ and $`f_{\mathrm{sky}}=1`$. We will refer to this case here and below as a “perfect experiment”.
With our fiducial choice of the gas bias, Planck should have a highly significant detection of the total signal. One should bear in mind that the bias parameter $`b_\pi `$ is still highly uncertain and that the $`S/N`$ scales as $`b_\pi ^2`$. Nevertheless even a relatively large reduction in the average gas temperature or density bias will not make the signal undetectable in principle. In practice, however remember that one is then relying on a precise subtraction of the noise bias in the measurement of $`C_l^{\mathrm{tot}}`$, which in turn requires that the power spectrum of the dust and other residual foregrounds lurking at least at the $`10\%`$ level in rms (1% in power) are determined comparably precisely.
If the fiducial SZ bias is close to correct, the high total single-to-noise in Planck can be used to break the measurement into bands in $`l`$ and recover the band power with errors
$$\left(\frac{\mathrm{\Delta }C_l^{\mathrm{SZ}}}{C_l^{\mathrm{SZ}}}\right)^2=\frac{f_{\mathrm{sky}}}{2}\underset{l_{\mathrm{band}}}{}(2l+1)\left(\frac{C_l^{\mathrm{SZ}}}{C_l^{\mathrm{tot}}}\right)^2.$$
(54)
We give an example from monte carlo realizations of the Gaussian noise and sample variance from the simulations in Fig. 4. Note that these are errors for a $`6^{}\times 6^{}`$ section of the sky and should be scaled by $`0.03f_{\mathrm{sky}}^{1/2}0.04`$ for Planck.
These signal-to-noise estimates assume that both the signal and noise are Gaussian. Of course in reality the SZ signal is non-Gaussian. In general, gravitational collapse correlates the amount of power in density fluctuations across all scales in the non-linear regime. However since the SZ effect probes many independent density fluctuations along the line-of-sight, the central limit theorem ensures that the SZ signal is far more Gaussian than the density field. We can test how much this affects the signal-to-noise with our simulations. Shown in Fig. 4 is the sampling errors on the band powers from the simulations themselves as compared with those from Gaussian realizations of the same power spectrum. The excess variance over the Gaussian limit is small on the relevant scales given detector noise limitations from Planck.
### 4.2. Skewness
The overall signal-to-noise for the measurement of the third moment of SZ effect is
$$\left(\frac{S}{N}\right)^2=f_{\mathrm{sky}}\frac{\mathrm{\Theta }^3(\widehat{𝐧};\sigma )^2}{\mathrm{Var}}$$
(55)
where the variance is given by
$`\mathrm{Var}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle \underset{l_1l_2l_3}{}}{\displaystyle \frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)^2`$ (59)
$`\times W_{l_1}^2(\sigma )W_{l_2}^2(\sigma )W_{l_3}^2(\sigma )6C_{l_1}^{\mathrm{tot}}C_{l_2}^{\mathrm{tot}}C_{l_3}^{\mathrm{tot}}.`$
In Fig. 9, we show the signal-to-noise for a measurement of the third moment as calculated under the HEPT. We compare the signal-to-noise in Planck with the ideal case of perfect removal of foregrounds and detector noise, and full sky coverage. We use here a tophat window in multipole space out to $`l_{\mathrm{max}}`$ to conform with other signal-to-noise considerations. Cosmic variance and Planck detector noise reduces the signal-to-noise values both at the low and high end for $`l_{\mathrm{max}}`$ values respectively. For Planck, the $`l`$ values in the range of few hundred to $``$ 1000 provides the maximal signal-to-noise for a measurement of the skewness. This corresponds to smoothing scales $`\sigma `$ 10’-30’ for tophat windows in angular space (c.f. Fig. 5). For MAP and Boomerang, the signal-to-noise values are $``$ 0.1, suggesting that a detection of SZ skewness is not likely to be possible in these two experiments.
Again equation (59) assumes Gaussian statistics for the variance and ignores the sample variance of the third moment itself. We test this approximation in Fig. 5 and find that it is reasonable given the level of residual noise for Planck. In constructing an estimator for $`S_3`$, it is important to remove the noise bias since noise variance will always reduce the skewness in the map. We do this by multiplying the estimator by $`(\mathrm{\Theta }_{\mathrm{tot}}^2/\mathrm{\Theta }_{\mathrm{SZ}}^2)^2`$.
Finally, note that in the noise-dominated regime the signal-to-noise in $`S_3`$ scales strongly with the gas bias $`S/Nb_\pi ^3`$, so that the detectability of this effect depends strongly on currently uncertain assumptions.
### 4.3. Bispectrum
The full bispectrum of the SZ effect contains all of the information about its three-point correlations induced by the growth of structure beyond the linear approximation. The skewness is simply one, easily measured, aspect of the bispectrum. The full signal-to-noise ratio of the bispectrum is
$$\left(\frac{S}{N}\right)^2=f_{\mathrm{sky}}\underset{l_1,l_2,l_3}{}\frac{B_{l_1l_2l_3}^2}{6C_{l_1}^{\mathrm{tot}}C_{l_2}^{\mathrm{tot}}C_{l_3}^{\mathrm{tot}}},$$
(60)
where $`C_l^{\mathrm{tot}}`$ follows Eq. (52). We plot the bispectrum cumulative signal-to-noise as a function of signal $`l_3`$, summed over $`l_1`$ and $`l_2`$. We refer the reader to Cooray & Hu (1999) for a detailed discussion on the bispectrum, its variance and the calculation of signal-to-noise ratio.
In Fig. 10, we show the expected cumulative signal-to-noise for the SZ bispectrum in Boomerang, MAP and Planck data and a perfect experiment. The signal-to-noise is calculated under the HEPT approximation for the underlying density field. As shown, MAP and Boomerang allow reasonable limits to be placed on any non-Gaussian signal in the SZ effect while Planck allows a strong possibility for a detection.
Again the same caveats as to the sensitivity of the $`S/N`$ estimate to the underlying assumptions that applied for the skewness also apply here. Moreover, measuring all the configurations of the bispectrum will be a formidable computational challenge as will control over systematic effects in the experiments.
### 4.4. Lensing Correlation
The SZ effect and weak gravitational lensing of the CMB both trace large-scale structure in the underlying density field. By measuring the correlation, one can directly test the manner in which gas pressure fluctuations trace the dark matter density fluctuations. The correlation vanishes in the two-point functions since the lensing does not affect an isotropic CMB due to conservation of surface brightness.
The correlation manifests itself as a non-vanishing bispectrum in the CMB at RJ frequencies (Goldberg & Spergel 1999; Cooray & Hu 1999). Again the cosmic variance from the primary anisotropies presents an obstacle for detection of the effect above the several arcminute scale ($`l2000`$). With the multifrequency cleaning of the SZ map presented here one can enhance the detectability of the effect.
Consider the bispectrum composed of one $`a_{lm}`$ from the cleaned SZ map and the other two from the CMB maps. Call this the SZ-CMB-CMB bispectrum. The noise variance of this term will be reduced by a factor of $`C_l^{\mathrm{tot}}/C_l^{\mathrm{CMB}}`$ compared with the CMB-CMB-CMB bispectrum. As one can see from Fig. 1 this can be up to a factor of $`10^3`$ in the variance. Details for the calculation of the CMB-CMB-CMB bispectrum are given in Cooray & Hu (1999). Here we have updated the normalization for SZ effect, taken $`f_{\mathrm{sky}}=0.65`$ for Planck’s useful sky coverage, and compared the $`S/N`$ of the two bispectra. As shown, the measurement using foreground cleaned Planck SZ and CMB maps has a substantially higher signal-to-noise than that from using the Planck CMB map alone for multipoles $`l10^210^3`$.
Beyond the improvement in signal-to-noise, however, there is an important advantage in constructing the SZ-lensing bispectrum using SZ and CMB maps. A mere measurement of the bispectrum in CMB data can lead to simultaneous detection of non-Gaussianities through processes other than just SZ-lensing cross-correlation. As discussed in Goldberg & Spergel (1999) and extended in Cooray & Hu (1999), gravitational lensing also correlates with other late time secondary anisotropy contributors such as integrated Sachs-Wolfe (ISW; Sachs & Wolfe 1967) effect and the reionized Doppler effect. In addition to lensing correlations, non-Gaussianities can also be generated through reionization and non-linear growth of perturbations (Spergel & Goldberg 1999; Goldberg & Spergel 1999; Cooray & Hu 1999). Bispectrum measurements at a signle frequency can result in a confusion as to the relative contribution from each of these scenarios. In Cooray & Hu (1999), we suggested the possibility of using differences in individual bispectra as a function of multipoles, however, such a separation can be problematic given that these differences are subtle (e.g., Fig 6 of Cooray & Hu 1999).
The construction of the SZ-lensing bispectrum using SZ and CMB maps has the advantage that one eliminates all possibilities, other than SZ, that result in a bispectrum. For effects related to SZ, the cross-correlation of lensing and SZ should produce the dominant signal; as shown in Cooray & Hu (1999), bispectra signal through SZ and reionization effects, such as Ostriker-Vishniac (OV; Ostriker & Vishniac 1986), are considerably smaller.
Conversely, multifrequency cleaning also eliminates the SZ contribution from the CMB maps and hence a main contaminant of the CMB-CMB-CMB bispectrum. This assists in the detection of smaller signals such as the ISW-lensing correlation, Doppler-lensing correlation or the non-Gaussianity of the initial conditions. Such an approach is highly desirable and Planck will allow such detailed studies to be carried out.
A potential caveat is that as noted above, the full bispectrum in an all-sky satellite experiment will be difficult to measure. Zaldarriaga & Seljak (1999) have developed a reduced set of three-point statistics optimized for lensing studies, based on a two point reconstruction of the lensing-convergence maps from temperature gradient information. They show that most of the information is retained in these statistics. Multifrequency cleaning improves the signal-to-noise for these statistics by exactly the same factor as for the full bispectrum.
## 5. Discussion
We have studied the prospects for extracting the statistical properties of the Sunyaev-Zel’dovich (SZ) effect associated with hot gas in large-scale structure using upcoming multifrequency CMB experiments. This gas currently remains undetected but may comprise a substantial fraction of the present day baryons. The SZ effect has a distinct spectral dependence with a null at a frequency of $``$ 217 GHz compared with true temperature anisotropies. This frequency dependence is what allows for effective separation of the SZ contribution with multifrequency CMB measurements.
As examples, we have employed the frequency and noise specifications of the Boomerang, MAP, Planck experiments. The MAP satellite only covers frequencies at RJ part of the frequency spectrum. Consequently, only Boomerang and Planck can take full advantage of multifrequency separation of the SZ and primary anisotropies. We have evaluated the detection threshold for SZ power spectrum measurements (see Fig. 3). Boomerang and MAP should provide limits on the degree scale fluctuations at the several $`\mu `$K level in rms; Planck should be able to detect sub $`\mu `$K signals.
The expected level of the SZ signal in the fiducial $`\mathrm{\Lambda }`$CDM model is still somewhat uncertain. We have employed a simple bias model for the pressure fluctuations, roughly normalized to recent hydrodynamic simulations (Refregier et al. 1999; Seljak et al. 2000), and calculated the resulting signal using analytic scaling relations and particle-mesh dark matter simulations. As hydrodynamic simulations improve, these techniques can be extended with more sophisticated modeling of the bias. They complement hydrodynamic simulations by extending the dynamic range and simulated volume, the latter being important for questions of sample variance.
Assuming this simplified model of the SZ signal, Planck should have signal-to-noise per multipole of order unity for $`l<1000`$. Although the recovered maps are then somewhat noisy, they are sufficient for precise determinations of low order statistics such as the SZ power spectrum, bispectrum and skewness (see Figs. 4-10). The skewness in principle can be used to separate the pressure bias from the underlying amplitude of the density fluctuations. The full bispectrum contains significantly more information but will be difficult to extract in its entirety. Current methods for measuring the bispectrum, tested with the COBE data, have concentrated at measuring specific modes such as $`l_1=l_2=l_3=l`$ (Ferreira et al. 1998). More work will clearly be required, especially in understanding the systematic errors at a sufficient level, but the wealth of information potentially present in the bispectrum should motivate efforts.
Note however that the non-Gaussianity in the SZ signal is not very strong due to the fact that it is constructed from many independent pressure fluctuations along the line of sight. As a consequence, we expect that signal-to-noise ratios can be estimated by Gaussian approximations, but that techniques that try to improve the SZ-primary separation based on non-Gaussianity (Hobson et al 1998; Aghanim & Forni 1999) may not be particularly effective for this signal.
We caution the reader that our oversimplification of the SZ signal can cause problems for a naive interpretation of future detections. For example, Seljak et al. (2000) find that the SZ power spectrum in their simulations is dominated by shot noise from the rare hot clusters not included in our modeling. Fortunately since these contributions are highly non-Gaussian, they can can readily be identified and removed. At the very least, $`X`$-ray bright clusters can be externally identified and removed; this has been shown to substantially reduce the shot noise contribution (Komatsu & Kitayama 1999). The effect we are modeling should be understood as the signal in fields without such clusters.
Another means of separating the SZ signal from large-scale structure from that of massive clusters is to cross correlate it with other tracers of large-scale structure that are less sensitive to highly overdense regions. An added benefit is that such a cross-correlation will also empirically measure the extent to which pressure fluctuations follow mass fluctuations. The CMB anisotropies themselves carry one such tracer in the form of the convergence from weak lensing. It manifests itself as a three-point correlation or bispectrum (Goldberg & Spergel 1999) but without frequency information it is severely sample-variance limited due to confusion noise from primary anisotropies. Measuring the SZ-lensing correlation using the cleaned SZ maps improves the signal-to-noise for the detection by over an order of magnitude at degree scales. Furthermore, the techniques introduced by Zaldarriaga & Seljak (1999) provide a concrete algorithm for extracting most of the three-point signal without recourse to measuring all the configurations of the bispectrum. Conversely, SZ removal from the CMB maps themselves can assist in the detection of other smaller bispectrum signals by eliminating one source of confusion noise.
The cross-correlation coefficient between the SZ effect and CMB weak lensing is relatively modest ($``$ 0.5, see Seljak et al. 2000). This is due to the fact that the SZ effect is a tracer of the nearby universe while CMB lensing is maximally sensitive to structure at $`z3`$. A higher correlation is expected if SZ is cross-correlated with an external probe of low redshift structure. Peiris & Spergel (2000) suggested the cross-correlation of MAP CMB data and Sloan<sup>4</sup><sup>4</sup>4http://www.sdss.org galaxy data. An improved approach would be to use the Planck derived SZ map rather than a CMB map. Using a SZ map reduces noise from the primary anisotropies and guarantees that any detection is due to correlations with the SZ effect. Extending the calculations in Peiris & Spergel (2000) with the Planck generated SZ map, we find signal-to-noise ratios which are on average greater by a factor of $``$ 10 when compared to signal-to-noise values using MAP CMB map. In fact with redshifts for galaxies, Planck SZ map can be cross-correlated in redshifts bins to study the redshift evolution of the gas. Other promising possibilities include cross correlation with soft X-ray background measurements, as well as ultraviolet and soft X-ray absorption line studies.
All these considerations imply a bright future for SZ studies of the hot gas associated with large-scale structure with wide-field multifrequency CMB observations. Its detailed properties should be revealed in its non-Gaussianity and correlation with other tracers of large-scale structure.
We thank Martin White for permission to adapt his PM ray tracing code for these purposes. We acknowledge useful discussions with Lloyd Knox, Joe Mohr, Roman Scoccimarro, Ned Wright and Matias Zaldarriaga. ARC is grateful to John Carlstrom, Michael Turner and Don York for helpful advice and financial support. WH is supported by the Keck Foundation, a Sloan Fellowship, and NSF-9513835. MT acknowledges NASA grant NAG5-6034 and Hubble Fellowship HF-01084.01-96A from STScI, operate by AURA, Inc. under NASA contract NAS5-26555. We acknowledge the use of CMBFAST (Seljak & Zaldarriaga 1996). |
warning/0002/hep-ph0002270.html | ar5iv | text | # Anomalous processes at high temperature and density in a 2-dimensional linear 𝜎 model
## 1 Introduction
In the past years, many works have been devoted to the study of anomalous processes at finite temperature and density, and in particular near the chiral symmetry restoration. In particular, Pisarski concluded that the neutral pion decay amplitude $`\pi ^o2\gamma `$ vanish when chiral symmetry is restored<sup>1</sup><sup>1</sup>1This does not contradict the well established fact that the coefficient of the axial anomaly is temperature independent (see for instance). Indeed, it has been shown in that the relationship between the anomaly and the amplitude can be modified by the existence of an additional 4-vector $`U_\mu `$ (the 4-velocity of the plasma in the observer’s frame) that can enter the general form of thermal amplitudes.. This conclusion is based on a direct calculation of the corresponding diagrams at finite temperature (and zero density) in the imaginary time formalism, and on symmetry considerations that forbid certain couplings in the symmetric phase. An additional conclusion was that this decay might be replaced by $`\pi ^o\sigma 2\gamma `$ in the chirally symmetric phase (which is allowed by the same symmetry argument), as indicated by a calculation of the box diagram at finite temperature.
This result has been confirmed by Baier, Dirks and Kober and by Salcedo using functional approaches in which one integrates out the fermions at the level of the the generating functional. Salcedo in gives also general arguments according to which $`\pi ^o2\gamma `$ (or $`\pi ^o\gamma `$ in two dimensions) should vanish in the chiral phase, and be replaced by amplitudes involving the $`\sigma `$ meson. From a technical perspective, the common point of all these studies is the use at some stage of the imaginary time formalism, in the limit of vanishing external momenta.
In another study, one of us studied how the neutral pion decay amplitude depends on the kinematical configuration of the external legs in order to explain the discrepancies found between and calculations performed in the real time formalism by . To that purpose, the $`\pi ^o\gamma \gamma `$ amplitude has been calculated at finite $`T`$ in the real time formalism, in the limit of small external momenta. It appeared that this limit cannot be uniquely defined (it depends on the path followed to reach the zero momenta point) and that the results of concerning this amplitude do not correspond to its on-shell value, but to a different way of reaching the limit<sup>2</sup><sup>2</sup>2The reason for this is easy to understand: since the energy variables are discrete in the imaginary time formalism, the only way one can consider the “zero momenta limit” in this formalism is to set first the discrete bosonic energies to zero, and then take the limit of zero three momenta. This way, the external momenta are forced to be space-like.. Additionally, the conclusion according to which the pion decay amplitude vanish above the critical point appeared to be questionable since the physical (on-shell) amplitude has a non vanishing limit at the critical point.
To accommodate this result with the general arguments provided in , one can notice that both Pisarski’s symmetry argument and Salcedo’s argument amount to the fact that one power of the quark mass (the mass the quarks acquire through the spontaneous breakdown of chiral symmetry, via their coupling to the average value of the $`\sigma `$ field) appears in the numerator when evaluating the Dirac trace associated to anomalous amplitudes. Therefore, since the average value $`\sigma `$ goes to zero in the chiral phase, the numerator vanishes above the critical point. Implicit in the argument is the fact that the correct dimension is provided by inverse powers of the temperature (as opposed to powers of the quark mass), as is the case in the imaginary time formalism at the static point. In other words, this argument is valid only if the denominator does not vanish when the mass goes to zero. This is precisely what fails when the amplitude is calculated on shell. One gets as expected one power of $`m=g\sigma `$ in the numerator, but the denominator turns out to be $`1/mT`$.
The remaining power of $`m`$ in the denominator indicates that the infrared or collinear behavior of the triangle diagram worsens when $`m0`$. In fact, as noted in and , the constituent quark mass $`m`$ ceases to be the relevant infrared regulator when $`m`$ is smaller than $`gT`$, and should be replaced by a thermal mass of order $`gT`$ that do not vanish in the chiral limit. Since an additional property of fermionic thermal masses is that they respect chiral symmetry, this thermal mass cannot appear in the Dirac’s trace. As a consequence the result $`m/mT`$ obtained for the on-shell amplitude in the bare theory becomes $`m/m_{\mathrm{th}}T`$ after one has resummed the quark thermal mass $`m_{\mathrm{th}}gT`$ (if $`mm_{\mathrm{th}}`$). The consequence of this regularization is that the resummed on-shell decay amplitude vanish in the chiral limit. In other words, Pisarski’s symmetry arguments holds for the physical amplitude only after a proper regularization (in order to get rid of all potential infrared or collinear singularities) has been issued.
Essential in this discussion is the influence of the kinematical conditions for the external legs on the infrared behavior of a thermal amplitude, since it can dramatically alter one’s conclusions. The second important point is that a calculation in the imaginary time formalism at the static point does not give a physical amplitude. Therefore, it would be interesting to test the second half of Pisarski’s conclusions, related to the $`\pi ^o\sigma 2\gamma `$ amplitude, by calculating this amplitude in the real time formalism and studying how it depends on the kinematics (up to now, this amplitude has only be considered at the static point). Since this amplitude is given by a four point function, it is a very complicated task to extract this behavior in its full generality. There is however a toy model in which this kind of study can be done quite simply: the 2-dimensional linear $`\sigma `$ model. Indeed, in this model the neutral pion decays into a single photon, and the analogous of the 4-point amplitude suggested by Pisarski would be $`\pi ^o\sigma \gamma `$, which, being a 3-point function, is rather easy to calculate.
The present paper is devoted to an analysis of the anomalous amplitudes in the 2-dimensional $`\sigma `$ model at finite temperature and chemical potential. We consider both the pion decay $`\pi ^o\gamma `$ 2-point function, and the $`\pi ^o\sigma \gamma `$ 3-point function. Emphasis is put on studying how these functions depend on the kinematics in the limit of small external momenta, near the chiral limit ($`m`$ small compared to $`\mu `$ and $`T`$).
The structure of the paper is as follows. Section 2 defines the model, as well as some notations and shorthands that will be used extensively later. In section 3, we calculate the amplitude for the $`\pi ^0\gamma `$ decay, and reduce it to a very compact form. In section 4, we study the amplitude of $`\pi ^0\sigma \gamma `$. Although a priori much more involved, this amplitude can also be reduced to a simple expression. All our results are expressed in terms of some function $`I(K)`$ defined by an integral. Basic properties and limits of this function are derived in appendix A. Finally, appendix B derives some relations between a few integrals that appear in intermediate stages of section 4.
## 2 Conventions and notations
We consider the 2-dimensional linear $`\sigma `$ model with two quark flavors, in which the mesons are coupled to quark fields as indicated by the following Lagrangian:
$$i\overline{\mathrm{\Psi }}/D\mathrm{\Psi }2g\overline{\mathrm{\Psi }}(\sigma t_0+i𝝅t\gamma ^5)\mathrm{\Psi },$$
(1)
where $`t_0=1/2`$ and $`\mathrm{Tr}(t^at^b)=\delta ^{ab}/2`$. We recall that in two dimensions the Dirac algebra is defined by the following set of relations:
$`\{\gamma ^\mu ,\gamma ^\nu \}=2g^{\mu \nu },`$
$`\gamma ^5={\displaystyle \frac{1}{2}}ϵ^{\mu \nu }\gamma _\mu \gamma _\nu ,`$ (2)
where $`ϵ^{\mu \nu }`$ is the 2-dimensional Levi-Civita tensor, normalized by $`ϵ^{01}=+1`$. For later use, let us quote first a generic trace formula:
$$\mathrm{Tr}(/A/B/C\gamma ^5\gamma ^\mu )=AB\mathrm{Tr}(/C\gamma ^5\gamma ^\mu )AC\mathrm{Tr}(/B\gamma ^5\gamma ^\mu )+BC\mathrm{Tr}(/A\gamma ^5\gamma ^\mu ).$$
(3)
In order to keep the following expressions compact, it is helpful to define the “dual” of a given vector by:
$$\stackrel{~}{A}^\mu ϵ^{\mu \nu }A_\nu ,$$
(4)
as well as the “wedge product” of two vectors:
$$ABϵ^{\mu \nu }A_\mu B_\nu .$$
(5)
According to these definitions, we have the obvious relations
$`\stackrel{~}{\stackrel{~}{A}}=A`$
$`AB=A\stackrel{~}{B}`$
$`(AB)^2=(AB)^2A^2B^2.`$ (6)
When $`A+B+C=0`$, we have also:
$$AB=BC=CA.$$
(7)
Finally, we have
$$\mathrm{Tr}(/A\gamma ^5\gamma ^\mu )=2\stackrel{~}{A}^\mu .$$
(8)
## 3 $`\pi ^0\gamma `$ amplitude
### 3.1 Retarded amplitude
We consider first the 1-loop contribution to the $`\pi ^0\gamma `$ decay amplitude depicted on figure 1.
The Feynman rules for the retarded-advanced formalism<sup>3</sup><sup>3</sup>3One should pay special attention to the chemical potential. Indeed, one should use a chemical potential $`\mu `$ in statistical weights where the Feynman rules give an argument $`k_0`$. In other words, $`\mu `$ appears in the formalism to account for the fact that the fermions carry some conserved charge, and the sign of this charge for a given propagator depends on how one orientates the propagator. give for the retarded amplitude the following expression
$`\mathrm{\Pi }__R^\mu (K)=eg{\displaystyle \frac{d^2L}{(2\pi )^2}\mathrm{Tr}((/L+m)\gamma ^\mu (/L+/K+m)\gamma ^5)}`$
$`\times \{n__F(l_0,\mu )\mathrm{Disc}\mathrm{\Delta }__R(L)\mathrm{\Delta }__R(L+K)`$
$`+n__F(l_0+k_0,\mu )\mathrm{Disc}\mathrm{\Delta }__R(L+K)\mathrm{\Delta }__A(L)\},`$ (9)
where $`\mathrm{\Delta }_{_{R,A}}(L)i/(L^2m^2\pm il_0ϵ)`$ are the retarded and advanced propagators, $`n__F(l_0,\mu )1/(\mathrm{exp}((l_0\mu )/T)+1)`$ is the Fermi-Dirac distribution function, and where the notation “Disc” denotes the discontinuity across the real energy axis:
$$\mathrm{Disc}\mathrm{\Delta }__R(L)\mathrm{\Delta }__R(L)\mathrm{\Delta }__A(L)=2\pi ϵ(l_0)\delta (L^2m^2).$$
(10)
The expression of the trace is very simple
$$\mathrm{Tr}((/L+m)\gamma ^\mu (/L+/K+m)\gamma ^5)=2m\stackrel{~}{K}^\mu ,$$
(11)
and in particular it makes obvious the fact that $`\mathrm{\Pi }__R^\mu `$ is transverse with respect to the photon momentum. Moreover, being independent of the loop momentum $`L`$, it can be immediately factorized out of the integral.
### 3.2 Zero momentum limit
At this point, it is convenient to perform the change of variable $`L+KL`$ on the second term of Eq. (9) in order to make the expression more symmetric. Then, the Dirac distributions hidden in the discontinuities make the integration over $`l_0`$ trivial, which gives<sup>4</sup><sup>4</sup>4We have dropped the R/A prescription for the denominator, since it can easily be recovered at the very end of the calculation by substituting $`k_0k_0+i0^+`$.
$`\mathrm{\Pi }__R^\mu (K)=2imeg\stackrel{~}{K}^\mu {\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \underset{\eta =\pm 1}{}}{\displaystyle \frac{1}{2L_\eta K+K^2}},`$ (12)
where we denote $`\omega _l\sqrt{l^2+m^2}`$ and $`L_\eta (\eta \omega _l,l)`$. It is now trivial to perform an expansion in powers of the external momentum $`K`$. The first term in this expansion, of degree $`0`$ in $`K`$, vanish because the corresponding integrand is an odd function of $`l`$. The first non vanishing term in this expansion comes at the next order, and is of degree $`1`$ in $`K`$:
$$\mathrm{\Pi }__R^\mu (K)=imeg\stackrel{~}{K}^\mu I(K),$$
(13)
where $`I(K)`$ is an homogeneous function of degree $`0`$ in $`K`$, containing the non trivial part of the momentum dependence, and defined by:
$$I(K)\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{dl}{2\pi }\frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}\frac{K^2}{(L_+K)^2}.$$
(14)
This integral is studied in some important limits in the appendix A.
### 3.3 Discussion
We observe for this amplitude the same features as in four dimensions. The most striking effect is related to what happens near the chiral limit $`m0`$, and is visible in formulas Eqs. (31) and (38). In the limit $`m/T,m/\mu 0`$, the function $`I(K)`$ behaves as follows:
$`\mathrm{If}K^2=0,I(K)={\displaystyle \frac{1}{2\pi m^2}},`$
$`\mathrm{If}K^20\mathrm{and}k_00,I(K)={\displaystyle \frac{1}{8mT}}{\displaystyle \frac{k_0}{\sqrt{K^2}}}{\displaystyle \frac{1}{\mathrm{cosh}^2(\mu /2T)}},`$
$`\mathrm{If}k_0=0,I(K)=\{\begin{array}{cc}& {\displaystyle \frac{7\zeta (3)}{8\pi ^3T^2}}\mathrm{if}\mu T\\ & {\displaystyle \frac{1}{4\pi \mu ^2}}\mathrm{if}T\mu \end{array}.`$ (15)
The configuration obtained with $`k_0=0`$ and $`\mu T`$ corresponds to Salcedo’s result, previously obtained in the imaginary time formalism at the static point. This point is particular because the first term in the expansion of $`I(K)`$ at small $`m`$ vanishes if $`k_0=0`$. For any other point, the expansion starts with a term behaving like $`I(K)1/mT`$. However, as one turns the chemical potential on, we see that this leading term is exponentially suppressed when $`\mu T`$. As a consequence, in a dense and cold system, the function $`I(K)`$ starts by a term in $`1/\mu ^2`$, whether $`k_0=0`$ or not.
The reason why the on-shell value is so singular when $`m0`$ is related to collinear singularities: in $`1+1`$ dimensions, all the spatial vectors are aligned, so that we are always at the most singular point<sup>5</sup><sup>5</sup>5This is to be contrasted with what happens in four dimensions: there the collinear singularities are softened by subsequent angular integrations, so that the on-shell amplitude is not exceptionally singular.. Here also, in order to be able to apply Pisarski’s argument, one should first regularize the theory by resumming a thermal mass. Then, all the powers of $`m`$ in the denominators would be replaced by powers of the thermal mass, leaving an uncompensated power of $`m0`$ in the numerator. From Eq. (13), we can write an effective Lagrangian coupling the neutral pion to the photon:
$`_{\pi ^o\gamma }`$ $`=egmϵ_{\mu \nu }{\displaystyle d^2xA^\mu (x)I(i_x)_x^\nu \pi ^o(x)}`$ (16)
$`=eg^2ϵ_{\mu \nu }{\displaystyle d^2xA^\mu (x)I(i_x)\sigma _x^\nu \pi ^o(x)}`$
This result completes the effective coupling found by Salcedo in , by incorporating all the nonlocal terms. The reason why the non-locality of this coupling has been missed in can be traced back in a misuse of the imaginary time techniques to get the zero momenta limit.
## 4 $`\pi ^0\sigma \gamma `$ amplitude
### 4.1 Retarded amplitude
We now consider the one-loop contribution to the $`\pi ^0\sigma \gamma `$ amplitude, represented on figure 2.
In the retarded-advanced formalism, the $`\mathrm{\Gamma }_{_{ARR}}^\mu `$ component of the vertex receives the following contribution from the first diagram:
$`\mathrm{\Gamma }_{_{ARR}}^{\mu ,(1)}(K,P,S)=ieg^2{\displaystyle \frac{d^2L}{(2\pi )^2}\mathrm{Tr}((/L+m)\gamma ^5(/L/P+m)\gamma ^\mu (/L+/S+m))}`$
$`\times \{n__F(l_0,\mu )\mathrm{Disc}\mathrm{\Delta }__R(L)\mathrm{\Delta }__R(L+S)\mathrm{\Delta }__A(LP)`$
$`+n__F(l_0p_0,\mu )\mathrm{Disc}\mathrm{\Delta }__R(LP)\mathrm{\Delta }__R(L)\mathrm{\Delta }__R(L+S)`$
$`+n__F(l_0+s_0,\mu )\mathrm{Disc}\mathrm{\Delta }__R(L+S)\mathrm{\Delta }__A(L)\mathrm{\Delta }__A(LP)\}.`$ (17)
Again, the expression becomes simpler if we perform the changes of variables $`LPL`$ on the second term, and $`L+SL`$ on the third one. This enables one to have common statistical weight and discontinuity for the three terms, the tradeoff being that the trace becomes different for the three terms. We can apply the same manipulations to the contribution of the second diagram. With the additional change $`LL`$ on the second diagram, we can merge the two contributions and find
$`\mathrm{\Gamma }_{_{ARR}}^\mu (K,P,S)=ieg^2{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \underset{\eta =\pm 1}{}}`$
$`\times \{{\displaystyle \frac{\mathrm{Tr}_a^\mu }{(2L_\eta K+K^2)(2L_\eta S+S^2)}}+{\displaystyle \frac{\mathrm{Tr}_b^\mu }{(2L_\eta P+P^2)(2L_\eta K+K^2)}}`$
$`+{\displaystyle \frac{\mathrm{Tr}_c^\mu }{(2L_\eta S+S^2)(2L_\eta P+P^2)}}\},`$ (18)
where the traces are given by
$`\mathrm{Tr}_a^\mu =2m^2\mathrm{Tr}(/P\gamma ^5\gamma ^\mu )2L_\eta S\mathrm{Tr}(/L_\eta \gamma ^5\gamma ^\mu )\mathrm{Tr}(/L_\eta /S/K\gamma ^5\gamma ^\mu )`$
$`\mathrm{Tr}_b^\mu =2m^2\mathrm{Tr}(/P\gamma ^5\gamma ^\mu )+2L_\eta P\mathrm{Tr}(/L_\eta \gamma ^5\gamma ^\mu )\mathrm{Tr}(/K/P/L_\eta \gamma ^5\gamma ^\mu )`$
$`\mathrm{Tr}_a^\mu =2m^2\mathrm{Tr}(/P\gamma ^5\gamma ^\mu )\mathrm{Tr}(/S/L_\eta /P\gamma ^5\gamma ^\mu ).`$ (19)
### 4.2 Zero momenta limit
We can now proceed with the expansion in powers of the external momenta. The first term in this expansion is of degree $`1`$ in the external momenta. An explicit calculation of this term shows that it vanishes thanks to energy-momentum conservation: $`P+K+S=0`$. The next term, of degree $`0`$ in the external momenta, vanishes also because the corresponding integrand is an odd function of $`l`$. Therefore, the first non vanishing term is of degree $`1`$ in the external momenta. An explicit extraction of this term gives
$`\mathrm{\Gamma }_{_{ARR}}^\mu (K,P,S)=i{\displaystyle \frac{eg^2}{8}}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \underset{\eta =\pm 1}{}}\{m^2\mathrm{Tr}(/P\gamma ^5\gamma ^\mu )`$
$`\times [{\displaystyle \frac{P^2}{(L_\eta P)^2}}{\displaystyle \frac{S^2}{(L_\eta S)^2}}+{\displaystyle \frac{S^2}{(L_\eta S)^2}}{\displaystyle \frac{K^2}{(L_\eta K)^2}}+{\displaystyle \frac{K^2}{(L_\eta K)^2}}{\displaystyle \frac{P^2}{(L_\eta P)^2}}`$
$`+{\displaystyle \frac{1}{(L_\eta K)(L_\eta P)(L_\eta S)}}({\displaystyle \frac{K^4}{L_\eta K}}+{\displaystyle \frac{P^4}{L_\eta P}}+{\displaystyle \frac{S^4}{L_\eta S}})]`$
$`+{\displaystyle \frac{\mathrm{Tr}(/L_\eta \gamma ^5\gamma ^\mu )}{L_\eta K}}\left[{\displaystyle \frac{K^2}{L_\eta K}}\left({\displaystyle \frac{P^2}{L_\eta P}}{\displaystyle \frac{S^2}{L_\eta S}}\right)+{\displaystyle \frac{S^4}{(L_\eta S)^2}}{\displaystyle \frac{P^4}{(L_\eta P)^2}}\right]`$
$`+{\displaystyle \frac{\mathrm{Tr}(/L_\eta /S/K\gamma ^5\gamma ^\mu )}{(L_\eta K)(L_\eta S)}}\left[{\displaystyle \frac{S^2}{L_\eta S}}{\displaystyle \frac{K^2}{L_\eta K}}\right]`$
$`+{\displaystyle \frac{\mathrm{Tr}(/K/P/L_\eta \gamma ^5\gamma ^\mu )}{(L_\eta P)(L_\eta K)}}\left[{\displaystyle \frac{K^2}{L_\eta K}}{\displaystyle \frac{P^2}{L_\eta P}}\right]`$
$`+{\displaystyle \frac{\mathrm{Tr}(/S/L_\eta /P\gamma ^5\gamma ^\mu )}{(L_\eta S)(L_\eta P)}}[{\displaystyle \frac{P^2}{L_\eta P}}{\displaystyle \frac{S^2}{L_\eta S}}]\}.`$ (20)
At this point, we can make use of Eqs. (3) and (8). It is now obvious that the result can be expressed in terms of the following integrals<sup>6</sup><sup>6</sup>6Note that we need only $`J_{_{AABC}}`$ when $`A+B+C=0`$. Having calculated $`J_{_{AABC}}`$ under this restrictive assumption, one cannot obtain $`J_{_{AABB}}`$ from it by enforcing $`C=B`$.
$`J_{_{AB}}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \frac{1}{L_+A}}{\displaystyle \frac{1}{L_+B}}`$
$`J_{_{AA}}{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \frac{1}{(L_+A)^2}}={\displaystyle \frac{I(A)}{A^2}}`$
$`J_{_{AAB}}^\mu {\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \frac{L_+^\mu }{(L_+A)^2}}{\displaystyle \frac{1}{L_+B}}`$
$`J_{_{AABC}}m^2{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \frac{1}{(L_+A)^2}}{\displaystyle \frac{1}{L_+B}}{\displaystyle \frac{1}{L_+C}}`$
$`J_{_{AABB}}m^2{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle \frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}}{\displaystyle \frac{1}{(L_+A)^2}}{\displaystyle \frac{1}{(L_+B)^2}},`$ (21)
where $`A,B,C`$ can be any of $`P,S,K`$.
### 4.3 Transversality
The major difference compared to the case of the $`\pi ^0\gamma `$ amplitude is that the above three-point amplitude is not manifestly transverse with respect to the photon momentum. The fact that the Dirac’s traces depend upon the loop momentum $`L`$ indicates that it could be necessary to calculate all the integrals before one can see the transversality. The situation is not so intricate though, since it happens that we only need to establish some relations between the various integrals defined in Eqs. (21). In appendix B, we show how the last three integrals of Eq. (21) can be expressed as functions of the first two. Using the relation given in Eqs. (40), (43) and (44), as well as Eq. (45), it is a simple matter of algebra to check the transversality of the $`\pi ^0\sigma \gamma `$ amplitude with respect to $`K`$, without the need of calculating the various $`J_{_{AA}}`$ and $`J_{_{AB}}`$.
Once we know that the result can indeed be written as $`\mathrm{\Gamma }^\mu =\mathrm{\Gamma }\stackrel{~}{K}^\mu `$, one can contract the amplitude with $`S`$ for instance in order to extract the coefficient $`\mathrm{\Gamma }`$. A straightforward calculation gives
$$\mathrm{\Gamma }_{_{ARR}}^\mu (K,P,S)=ieg^2\stackrel{~}{K}^\mu F(P,S),$$
(22)
with
$$F(P,S)\frac{P^2S^2I(S)+(PS)[P^2I(P)+(PK)I(K)]}{(PS)^2},$$
(23)
where implicitly $`K=PS`$. At first sight, this expression could explode whenever two momenta become parallel. However, one can check that this is not the case, because the numerator behaves like $`(PS)^2`$ when $`PS`$ becomes small.
### 4.4 Discussion
We see that this 3-point amplitude involving the $`\sigma `$ field depends on the same function $`I(.)`$ defined above, and has one power of the mass $`m`$ less when compared to the $`\pi ^o\gamma `$ amplitude, in agreement with the general arguments of .
Again, it is found that this limit depends on the kinematics, i.e. on the way one is approaching the zero momenta limit. In particular, the way this amplitude depends on $`m`$ at small $`m`$ depends strongly on the kinematics. This is to be contrasted with the result of , which only picked one particular limit. Except at the static point ($`k_0=p_0=s_0=0`$), this amplitude becomes singular at the critical point where $`m0`$, indicating the necessity of regularizing the fermion propagator by a thermal mass. After this resummation has been performed, this amplitude is regular but does not vanish in the chirally symmetric phase. Therefore, the conjecture of holds, but only after infrared regularization.
One can also write an effective coupling associated with this amplitude:
$`_{\pi ^o\sigma \gamma }`$ $`=eg^2ϵ_{\mu \nu }{\displaystyle d^2xA^\mu (x)F(i_{x_1},i_{x_2})}`$ (24)
$`\times [\sigma (x_2)_{x_1}^\nu \pi ^o(x_1)+\pi ^o(x_1)_{x_2}^\nu \sigma (x_2)]|_{x_1=x_2=x}.`$
Of course, if one uses Eq. (25), one finds a local limit for this effective coupling in the limit of zero temperature and density.
## 5 Conclusions
In this paper, we have studied the $`\pi ^o\gamma `$ and $`\pi ^o\sigma \gamma `$ amplitudes in the 2-dimensional $`\sigma `$ model at finite temperature and density. For both of these amplitudes, the zero momenta limit is not unique and strongly depends on the kinematical configuration. In particular, the imaginary time formalism should be used with great care when looking at this limit. Indeed, if one first sets the external discrete energies to zero, then all the information regarding the non-locality of the amplitude is lost, and in particular the physical limit cannot be recovered. A proper way to use the imaginary time formalism would be to perform the sum over the loop discrete energies while keeping nonzero external discrete energies. After that, one should perform the analytical continuation to real external energies, and only then consider the zero momenta limit.
The other conclusion of this work is that collinear or infrared singularities spoil the general symmetry arguments given in to justify the nullity of the pion decay into photons in the chirally symmetric phase: the on-shell decay amplitude in the bare theory does not vanish. For these arguments to be valid, one should perform the resummation of a thermal mass that will regularize the fermion propagators.
If such a regularization is used, then the conclusion is that $`\pi ^o\gamma `$ vanishes when $`m0`$, while $`\pi ^o\sigma \gamma `$ does not, in agreement with the conjecture of .
## Acknowledgements
We thank the Erwin Schrödinger Institute for support and hospitality during the worshop “BRST cohomology, quantization and anomalies”, where this work started. The work of F.G. is supported by DOE under grant DE-AC02-98CH10886.
## Appendix A Properties of the function $`I(K)`$
### A.1 Vacuum limit
The purpose of this appendix is to study the integral $`I(K)`$ since all the quantities calculated in this paper can be expressed in terms of this function. A first check is to look at the zero temperature and chemical potential limit of this function, which gives immediately
$$\underset{T,\mu 0^+}{lim}I(K)=\frac{1}{2\pi m^2}.$$
(25)
As one could expect, this $`T=\mu =0`$ limit does not exhibit any non-locality in its momentum dependence, since this is a purely thermal feature.
### A.2 Transformation into a sum
A convenient way to look at the high temperature or density limit is to turn the integral defining $`I(K)`$ in Eq. (14) into a sum, by making use of the following identity<sup>7</sup><sup>7</sup>7To derive this formula, one can start from Mittag-Leffler’s expansion of the $`\mathrm{cot}`$ function :
$$\mathrm{cot}(z)=\frac{1}{z}+2\underset{n=1}{\overset{+\mathrm{}}{}}\frac{z}{z^2n^2\pi ^2}.$$
(26) :
$$\frac{1}{e^x+1}=\frac{1}{2}2x\underset{n=0}{\overset{+\mathrm{}}{}}\frac{1}{x^2+(2n+1)^2\pi ^2}.$$
(27)
This identity enables one to rewrite $`n__F(\omega _l,\mu )n__F(\omega _l,\mu )`$ as a series, from which it is straightforward to first obtain
$$I(K)=\frac{1}{\pi T^2}\frac{\kappa ^2+1}{\kappa ^21}\underset{n=0}{\overset{+\mathrm{}}{}}\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑v\frac{v^2+\xi _+^2}{(v^2+\xi _{}^2)^2}\frac{v^2+A_n^2}{(v^2+A_n^2)^2+B_n^2},$$
(28)
where we denote
$`v{\displaystyle \frac{l}{T}},\kappa k_0/k,\xi _\pm ^2{\displaystyle \frac{m^2}{T^2}}{\displaystyle \frac{\kappa ^2}{\kappa ^2\pm 1}}`$
$`A_n^2\pi ^2(2n+1)^2+{\displaystyle \frac{m^2}{T^2}}{\displaystyle \frac{\mu ^2}{T^2}},B_n^24\pi ^2(2n+1)^2{\displaystyle \frac{\mu ^2}{T^2}}.`$ (29)
At this stage, it remains to perform term by term the integration over $`dv`$, which is elementary and yields
$`{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}dv{\displaystyle \frac{v^2+\xi _+^2}{(v^2+\xi _{}^2)^2}}{\displaystyle \frac{v^2+A_n^2}{(v^2+A_n^2)^2+B_n^2}}={\displaystyle \frac{\pi }{2\xi _{}}}{\displaystyle \frac{}{\xi _{}}}[{\displaystyle \frac{1}{\xi _{}}}{\displaystyle \frac{(\xi _+^2\xi _{}^2)(A_n^2\xi _{}^2)}{(A_n^2\xi _{}^2)^2+B_n^2}}`$
$`+{\displaystyle \frac{A_n^4+B_n^22\xi _+^2A_n^2+(\xi _+^2\xi _{}^2)\sqrt{A_n^4+B_n^2}+\xi _+^2\xi _{}^2}{(A_n^2\xi _{}^2)^2+B_n^2}}`$
$`\times \sqrt{{\displaystyle \frac{\sqrt{A_n^4+B_n^2}+A_n^2}{2(A_n^4+B_n^2)}}}].`$ (30)
### A.3 Chiral limit far from the light cone
Another interesting limit is the chiral limit, where $`m/T`$ goes to zero, while $`\mu /T`$ is kept fixed. It is very easy to extract from the above formula a systematic expansion in powers of $`m/T`$. It is just a matter of expanding at small $`\xi _\pm `$ the above expression, which gives for the first two orders:
$`I(K)={\displaystyle \frac{1}{mT}}\left[{\displaystyle \frac{k_0}{\sqrt{K^2}}}F_0\left({\displaystyle \frac{\mu }{T}}\right)+{\displaystyle \frac{m}{T}}{\displaystyle \frac{K^2+2k^2}{K^2}}F_1\left({\displaystyle \frac{\mu }{T}}\right)+𝒪\left({\displaystyle \frac{m^2}{T^2}}\right)\right],`$ (31)
where the coefficients are given by
$`F_0\left({\displaystyle \frac{\mu }{T}}\right)={\displaystyle \frac{1}{\pi ^2}}\mathrm{Re}{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{1}{(2n+1+i\mu /\pi T)^2}},`$
$`F_1\left({\displaystyle \frac{\mu }{T}}\right)={\displaystyle \frac{1}{\pi ^3}}{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}(2n+1){\displaystyle \frac{(2n+1)^23\mu ^2/\pi ^2T^2}{((2n+1)^2+\mu ^2/\pi ^2T^2)^3}}.`$ (32)
Introducing the “digamma” function
$$\psi (z)\frac{d}{dz}\mathrm{ln}\mathrm{\Gamma }(z)$$
(33)
and a series representation of its first derivative
$$\psi ^{}\left(\frac{1+z}{2}\right)=4\underset{n=0}{\overset{+\mathrm{}}{}}\frac{1}{(2n+1+z)^2},$$
(34)
it is immediate to check<sup>8</sup><sup>8</sup>8The equality on the first line is exact and comes from the formula
$$\mathrm{\Gamma }(z)\mathrm{\Gamma }(1z)=\frac{\pi }{\mathrm{sin}(\pi z)},$$
(35) while the limit $`\mu T`$ in the last line is obtained from Stirling’s expansion for $`\mathrm{\Gamma }(z)`$.
$`F_0\left({\displaystyle \frac{\mu }{T}}\right)={\displaystyle \frac{1}{8\pi ^2}}\left[\psi ^{}\left({\displaystyle \frac{1}{2}}+i{\displaystyle \frac{\mu }{2\pi T}}\right)+\psi ^{}\left({\displaystyle \frac{1}{2}}i{\displaystyle \frac{\mu }{2\pi T}}\right)\right]={\displaystyle \frac{1}{8\mathrm{cosh}^2(\mu /2T)}},`$
$`F_1\left({\displaystyle \frac{\mu }{T}}\right)={\displaystyle \frac{1}{32\pi ^3}}\left[\psi ^{\prime \prime }\left({\displaystyle \frac{1}{2}}+i{\displaystyle \frac{\mu }{2\pi T}}\right)+\psi ^{\prime \prime }\left({\displaystyle \frac{1}{2}}i{\displaystyle \frac{\mu }{2\pi T}}\right)\right]`$
$`\{\begin{array}{cc}& {\displaystyle \frac{7\zeta (3)}{8\pi ^3}}\mathrm{if}\mu T\\ & {\displaystyle \frac{T^2}{4\pi \mu ^2}}\mathrm{if}T\mu \end{array}.`$ (36)
### A.4 Light cone limit
Note however that this expansion is not valid near the light cone. Indeed, its derivation assumed that a small $`m/T`$ would imply a small $`\xi _{}`$, which is not true if $`K^2`$ is small (or equivalently $`\kappa ^21`$). Sufficiently close to the light cone, $`\xi _{}`$ becomes large and a different kind of expansion must be considered. In this region of phase space, one can write:
$`\underset{K^20}{lim}I(K)`$ $`=\underset{\xi _{}+\mathrm{}}{lim}{\displaystyle \frac{1}{2\xi _{}T^2}}{\displaystyle \frac{2}{\kappa ^21}}{\displaystyle \underset{n=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{1}{[\xi _{}+(2n+1)\pi ]^2}}`$ (37)
$`=\underset{\xi _{}+\mathrm{}}{lim}{\displaystyle \frac{\xi _{}}{m^2}}{\displaystyle \frac{1}{4\pi ^2}}\psi ^{}\left({\displaystyle \frac{1}{2}}+{\displaystyle \frac{\xi _{}}{2\pi }}\right).`$
where we have used $`(\kappa ^21)^1\xi _{}^2T^2/m^2`$. Making use of Stirling’s formula , we obtain:
$$\underset{K^20}{lim}I(K)=\frac{1}{2\pi m^2}.$$
(38)
We notice that the on-shell value of $`I(K)`$ is totally immune to corrections due to temperature or density. The origin of this property can be understood from Eq. (14): when $`K^2=0`$, the denominator behaves like $`k^2m^4/l^4`$ (this is reminiscent of a collinear singularity cured by the mass $`m`$), and the integral is completely dominated by its ultraviolet sector. As a consequence, $`I(K)`$ is saturated by the vacuum contribution for this value of $`K^2`$. Away from the light cone, thermal corrections of order $`1/mT`$ appear in $`I(K)`$.
## Appendix B Relations between some integrals
In this appendix, we establish some useful relations between the five integrals defined in Eq. (21). Let us start by the study of $`J_{_{AAB}}^\mu `$: this integral satisfies the $`2\times 2`$ linear system
$$\{\begin{array}{cc}A_\mu J_{_{AAB}}^\mu & =J_{_{AB}}\\ B_\mu J_{_{AAB}}^\mu & =J_{_{AA}},\end{array}$$
(39)
the resolution of which gives the two components of $`J_{_{AAB}}^\mu `$ as functions of $`J_{_{AA}}`$ and $`J_{_{AB}}`$:
$$J_{_{AAB}}^\mu =\frac{\stackrel{~}{B}^\mu J_{_{AB}}\stackrel{~}{A}^\mu J_{_{AA}}}{AB}.$$
(40)
In order to obtain $`J_{_{AABB}}`$ it is convenient to define first the second rank tensor
$$J_{_{AABB}}^{\mu \nu }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{dl}{2\pi }\frac{n__F(\omega _l,\mu )n__F(\omega _l,\mu )}{2\omega _l}\frac{L_+^\mu }{(L_+A)^2}\frac{L_+^\nu }{(L_+B)^2},$$
(41)
from which we can obtain $`J_{_{AABB}}`$ as $`g_{\mu \nu }J_{_{AABB}}^{\mu \nu }`$. The three independent components of this symmetric tensor can be obtained via the resolution of the following $`3\times 3`$ linear system
$$\{\begin{array}{cc}A_\mu B_\nu J_{_{AABB}}^{\mu \nu }& =J_{_{AB}}\\ A_\mu A_\nu J_{_{AABB}}^{\mu \nu }& =J_{_{BB}}\\ B_\mu B_\nu J_{_{AABB}}^{\mu \nu }& =J_{_{AA}},\end{array}$$
(42)
which finally gives
$$J_{_{AABB}}=\frac{2(AB)J_{_{AB}}A^2J_{_{AA}}B^2J_{_{BB}}}{(AB)^2}.$$
(43)
The same method can be applied to $`J_{_{AABC}}`$, which gives:
$$J_{_{AABC}}=\frac{A^2J_{_{AA}}+B^2J_{_{AB}}+C^2J_{_{AC}}}{(AB)^2},$$
(44)
under the assumption that $`A+B+C=0`$. We also need the following relation
$$J_{_{AB}}+J_{_{BC}}+J_{_{CA}}=0,$$
(45)
which is valid when $`A+B+C=0`$. |
warning/0002/hep-th0002256.html | ar5iv | text | # References
INFNCA-TH0004
Symmetry Breaking, Central Charges
and the AdS<sub>2</sub>/CFT<sub>1</sub> Correspondence
Mariano Cadoni<sup>a,c,∗</sup> and Salvatore Mignemi<sup>b,c,∗∗</sup>
<sup>a</sup>Dipartimento di Fisica, Università di Cagliari,
Cittadella Universitaria, 09042, Monserrato, Italy.
<sup>b</sup> Dipartimento di Matematica, Università di Cagliari,
viale Merello 92, 09123, Cagliari, Italy.
<sup>c</sup> INFN, Sezione di Cagliari.
## Abstract
When two-dimensional Anti-de Sitter space ($`\mathrm{AdS}_2`$ ) is endowed with a non-constant dilaton the origin of the central charge in the Virasoro algebra generating the asymptotic symmetries of $`\mathrm{AdS}_2`$ can be traced back to the breaking of the $`SL(2,R)`$ isometry group of $`\mathrm{AdS}_2`$ . We use this fact to clarify some controversial results appeared in the literature about the value of the central charge in these models.
E-Mail: CADONI@CA.INFN.IT
<sup>∗∗</sup>E-Mail: MIGNEMI@CA.INFN.IT
The Anti-de Sitter(AdS)/conformal field theory (CFT) correspondence in two spacetime dimensions seems to contradict the general belief that low-dimensional physics is simpler than the higher-dimensional one. Compared with the higher-dimensional cases, the $`\mathrm{AdS}_2/\mathrm{CFT}_1`$ duality has many puzzling and controversial features. Even though general arguments suggest that gravity on $`\mathrm{AdS}_2`$ should be related to some sort of conformal mechanics living on the one-dimensional boundary of $`\mathrm{AdS}_2`$ , the realization of this correspondence is rather involved. Although it has been shown that, analogously to the three-dimensional (3D) case , the conformal symmetry involved is infinite dimensional , the search for a physical system realizing this symmetry has not been very successful. On the other hand, the fact that $`\mathrm{AdS}_2`$ has a disconnected boundary makes the $`\mathrm{AdS}_2/\mathrm{CFT}_1`$ correspondence problematic even at a more fundamental level.
At a technical level there are features of gravity on $`\mathrm{AdS}_2`$ that represent a further difficulty for sheding light on the subject, e.g. the existence of several two-dimensional (2D) gravity models admitting $`\mathrm{AdS}_2`$ as solution ( models with or without dilatons) or the fact that the boundary of $`\mathrm{AdS}_2`$ is one-dimensional.
In a previous work , we have been able, working in the context of 2D dilaton gravity, to show that the asymptotic symmetry group of $`\mathrm{AdS}_2`$ is generated by a Virasoro algebra. Using a canonical realization of the symmetry we have also computed the central charge of the algebra. Unfortunately, using our value of the central charge for the computation of the statistical entropy of 2D black holes, we found a discrepancy of a factor of $`\sqrt{2}`$ with respect to the thermodynamical entropy.
Later Navarro-Salas and Navarro , using a boundary field realization of the symmetry, found a value of the central charge (half of our result) that produces a statistical entropy in agreement with the thermodynamical one.
In this letter we try to clarify this controversial point. We show that in models where $`\mathrm{AdS}_2`$ is endowed with a non-constant dilaton, the origin of the central charge of the Virasoro algebra can be traced back to the breakdown of the $`SL(2,R)`$ isometry group of $`\mathrm{AdS}_2`$ . The authors of Ref. use instead a non-scalar dilaton, which enables them to keep the $`SL(2,R)`$ symmetry of $`\mathrm{AdS}_2`$ unbroken. However, using a non-scalar dilaton, they have to give up the diffeomorphisms invariance of the 2D dilaton gravity theory. We will show that keeping the dilaton to transform as a scalar, hence considering a truly diffeomorphisms invariance 2D dilaton gravity theory, our previous result of Ref. for the central charge can be recovered also adopting a boundary field realization of the asymptotic symmetry of $`\mathrm{AdS}_2`$ .
Let us consider a generic two-dimensional (2D) dilaton gravity model,
$$S=\frac{1}{2}d^2x\sqrt{g}\left[\eta R+\lambda ^2V(\eta )\right],$$
(1)
where $`\eta `$ is a scalar field related to the usual definition of the dilaton $`\varphi `$ by $`\eta =\mathrm{exp}(2\varphi )`$. In the discussion of the symmetries of the model a crucial role is played by the scalar $`\eta `$. In general, the scalar character of $`\eta `$ implies that the symmetries (isometries) of the 2D spacetime are broken by a non-costant dilaton. In fact under the isometric transformations generated by a Killing vector $`\chi `$,
$$\delta \eta =_\chi \eta =\chi ^\mu _\mu \eta .$$
(2)
A non-constant dilaton in general implies $`_\chi \eta 0`$. On the other hand, 2D dilaton gravity always admits a Killing vector given by :
$$\chi _{(1)}^\mu =ϵ^{\mu \nu }_\nu \eta ,$$
(3)
which leaves $`\eta `$ invariant. Thus, a non-constant field $`\eta `$ in general breaks down the isometry group of the metric to the subgroup generated by the Killing vector $`\chi _{(1)}`$.
The previous features of 2D dilaton gravity theories have strong analogies with spontaneous symmetry breaking in ordinary field theory. In the case under consideration the quantity characterizing the breaking of the symmetry is $`_\mu \eta `$. The analogy is particularly evident if one considers classical solutions of the model (1). These solutions are characterized by a mass $`M`$ and a temperature $`T`$ (when the solutions can be interpreted as black holes). Both $`M`$ and $`T`$ are geometric objects, so that they are invariant under the isometry group of the metric, and can be expressed in terms of $`\eta `$
$$M=F_0\left[𝑑\eta \lambda ^2V(\eta )(\eta )^2\right],TV(\eta _h),$$
(4)
where $`F_0`$ is a constant related to the normalization of the Killing vectors (in the following we use $`F_0=1/2\lambda \eta _0)`$ and $`\eta _h`$ is the scalar evaluated on the horizon of the 2D black hole. Using the field equations one can easily show that the solutions characterized by a constant dilaton, thus preserving the isometry group of the metric, have zero mass and temperature. Only symmetry-breaking excitations, characterized by a non constant dilaton, can have $`M0`$, $`T0`$.
Let us now apply the previous considerations to $`\mathrm{AdS}_2`$ . With a dilaton potential $`V=2\eta `$ the action (1) describes the Jackiw-Teitelboim (JT) model . $`\mathrm{AdS}_2`$ or more generally black holes in $`\mathrm{AdS}_2`$ , are solutions of the model . Owing to Birkhoff’s theorem, the solutions can always be written in a static form, with dilaton $`\eta =\eta _0\lambda x`$, where $`\eta _0`$ is an integration constant (the field equations, but not the action, are invariant for rescaling of $`\eta `$ by a constant) and metric
$$ds^2=(\lambda ^2x^2a^2)dt^2+(\lambda ^2x^2a^2)^1dx^2,$$
(5)
where $`a^2`$ is proportional to the mass of the black hole.
Being a maximally symmetric spacetime $`\mathrm{AdS}_2`$ admits three Killing vectors generating the $`SO(1,2)SL(2,R)`$ group of isometries. It is evident that the static solution for $`\eta `$ is invariant only under the action of the Killing vector (3), which in this case describes time-translations $`T`$. Thus, the non-constant value of the dilaton breaks $`SL(2,R)T`$. The parameter characterizing the symmetry breaking is $`_x\eta =\eta _0\lambda `$.
Similar considerations hold when one considers the asymptotic symmetries of $`\mathrm{AdS}_2`$ . These are defined as the transformations which leave the asymptotic form of the $`\mathrm{AdS}_2`$ metric invariant, and where shown in to generate a Virasoro algebra with non-trivial central charge. Our aim here is to give a realization of this algebra in terms of fields which describe the degrees of freedom of the boundary. For this purpose, it is useful to adopt the formalism introduced in .
We define a two-dimensional metric to be asymptotically $`\mathrm{AdS}_2`$ if, for $`x\mathrm{}`$, it behaves as
$`g_{tt}`$ $`=`$ $`\lambda ^2x^2+\gamma _{tt}(t)+o\left({\displaystyle \frac{1}{x^2}}\right),`$
$`g_{tx}`$ $`=`$ $`{\displaystyle \frac{\gamma _{tx}(t)}{\lambda ^3x^3}}+o\left({\displaystyle \frac{1}{x^5}}\right),`$
$`g_{xx}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda ^2x^2}}+{\displaystyle \frac{\gamma _{xx}(t)}{\lambda ^4x^4}}+o\left({\displaystyle \frac{1}{x^6}}\right),`$ (6)
where the fields $`\gamma _{\mu \nu }`$ parametrize the first sub-leading terms in the expansion and can be interpreted as deformations on the boundary.
The asymptotic form (6) of the metric is preserved by infinitesimal diffeomorphisms $`\chi ^\mu (x,t)`$ of the form
$`\chi ^t`$ $`=`$ $`ϵ(t)+{\displaystyle \frac{\ddot{ϵ}(t)}{2\lambda ^4x^2}}+{\displaystyle \frac{\alpha ^t(t)}{x^4}}+o\left({\displaystyle \frac{1}{x^5}}\right),`$
$`\chi ^x`$ $`=`$ $`x\dot{ϵ}(t)+{\displaystyle \frac{\alpha ^x(t)}{x}}+o\left({\displaystyle \frac{1}{x^2}}\right).`$ (7)
where $`ϵ(t)`$ and $`\alpha ^\nu (t)`$ are arbitrary and a dot denotes time derivative. $`\alpha ^\nu `$ describes ”pure gauge” 2D diffeomorphisms which affect only the fields on the boundary. In Ref. is shown that the symmetries (7) are generated by a Virasoro algebra.
In view of (2), the asymptotic behaviour of the scalar field $`\eta `$, compatible with the transformations (7), must take the form
$$\eta =\eta _0\left(\lambda \rho (t)x+\frac{\gamma _{\varphi \varphi }(t)}{2\lambda x}\right)+o\left(\frac{1}{x^3}\right),$$
(8)
where $`\rho `$ and $`\gamma _{\varphi \varphi }`$ play a role analogous to that of the $`\gamma _{\mu \nu }`$. On shell hold the constraints
$$\lambda ^2\ddot{\rho }=\rho (\gamma _{tt}\gamma _{xx})\gamma _{\varphi \varphi },$$
(9)
$$\dot{\rho }\gamma _{tt}+\frac{\rho }{2}\dot{\gamma }_{xx}+\dot{\gamma }_{\varphi \varphi }=0.$$
(10)
Eq. (8) implies that the full, infinite dimensional, asymptotic symmetry of $`\mathrm{AdS}_2`$ is broken by the boundary condition for $`\eta `$. In fact, only the asymptotic transformations generated by the Killing vector (3), which correspond to $`ϵ=\rho `$ in Eqs. (7), leave $`\eta `$ asymptotically invariant. The quantity characterizing the symmetry breaking is now the boundary field $`\rho =_x\eta +o(x^1)`$.
In it was assumed that $`\rho (t)1`$. However, this is at variance with the transformation (2) of $`\eta `$ under diffeomorphisms, which implies that the background dilaton $`\eta =\eta _0\lambda x`$, is transformed by (7) into the form (8) with generic $`\rho (t)`$. In order to avoid this problem, the authors of had instead to assume that under diffeomorphisms the dilaton transform as $`\delta \eta =\chi ^\mu _\mu \eta +_t\chi ^t\eta `$ , but this seems rather ad hoc and moreover spoils the diffeomorphism invariance of the action (1), presumably leading to inconsistencies. It is important to notice that, if $`\eta `$ transforms according to the previous transformation law, the asymptotic symmetry group of $`\mathrm{AdS}_2`$ is no longer broken. In fact, in this case, asymptotically $`\delta \eta =0`$ under the full group generated by the Killing vectors (7).
Although the asymptotic symmetry of $`\mathrm{AdS}_2`$ is broken by $`\eta `$, the boundary fields $`\gamma _{tt}`$, $`\gamma _{xx}`$, $`\gamma _{\varphi \varphi }`$, $`\rho `$ still span a representation of the full infinite dimensional group generated by the Killing vectors (7). In fact, under the asymptotic symmetries (7), the boundary fields transform as
$`\delta \gamma _{tt}`$ $`=`$ $`ϵ\dot{\gamma }_{tt}+2\dot{ϵ}\gamma _{tt}{\displaystyle \frac{\stackrel{\mathrm{}}{ϵ}}{\lambda ^2}}2\lambda ^2\alpha ^x,`$
$`\delta \gamma _{xx}`$ $`=`$ $`ϵ\dot{\gamma }_{xx}+2\dot{ϵ}\gamma _{xx}4\lambda ^2\alpha ^x,`$
$`\delta \gamma _{\varphi \varphi }`$ $`=`$ $`ϵ\dot{\gamma }_{\varphi \varphi }+\dot{ϵ}\gamma _{\varphi \varphi }+{\displaystyle \frac{\ddot{ϵ}\dot{\rho }}{\lambda ^2}}+2\lambda ^2\rho \alpha ^x,`$
$`\delta \rho `$ $`=`$ $`ϵ\dot{\rho }\dot{ϵ}\rho .`$ (11)
These transformations are easily recognized as (anomalous) transformation laws for conformal fields of weight respectively $`2,2,1,1`$. The anomalous parts of the transformation (the terms proportional to $`\ddot{ϵ}`$ and $`\stackrel{\mathrm{}}{ϵ}`$) are related with the $`SL(2,R)`$ symmetry breaking. In fact, the anomalous terms are connected one with the other by using the equation of motion (9), whereas the anomalous term appearing in the transformation of $`\gamma _{\varphi \varphi }`$ is proportional to $`\dot{\rho }`$. It follows that, when $`\mathrm{AdS}_2`$ is endowed with a non-constant dilaton, the origin of the anomalous terms in the transformation laws for conformal boundary fields (hence of the central charge) can be traced back to the breakdown of the full asymptotic isometry group of the spacetime. This implies that in general one has a $`\rho `$-dependent central charge. However, we are only interested in classical solutions of the 2D gravity model. Because of Birkhoff’s theorem, we can, without loss of generality, limit ourselves to the field configurations with $`\rho =`$ const. Moreover, the equations of motions of the JT model are invariant under the rescaling of $`\eta `$ by a constant. It will therefore be sufficient to consider only the $`\rho =1`$ configuration.
The next step in our analysis is the construction of the generator of the conformal symmetry in terms of the boundary fields $`\gamma _{tt},\gamma _{xx},\gamma _{\varphi \varphi },\rho `$. This generator has a natural interpretation as the stress-energy tensor associated with the one-dimensional conformal field theory living on the boundary of $`\mathrm{AdS}_2`$ . A natural candidate for such a generator is the charge $`J(ϵ)`$, which in the canonical formalism can be associated with the asymptotic symmetries (7). The charge $`J(ϵ)`$ is defined in terms of the boundary contribution one must add to the Hamiltonian in order to have well defined variational derivatives ,
$$\delta J=\underset{x\mathrm{}}{lim}[N(\sigma ^1\delta \eta ^{}\sigma ^2\eta ^{}\delta \sigma )N^{}(\sigma ^1\delta \eta )+N^x(\mathrm{\Pi }_\eta \delta \eta \sigma \delta \mathrm{\Pi }_\sigma )].$$
(12)
Using the boundary conditions (6), (8) one finds,
$$\delta J[ϵ]=\eta _0\left[\lambda ϵ(\gamma _{tt}\delta \rho +\frac{\rho }{2}\delta \gamma _{xx}+\delta \gamma _{\varphi \varphi })+\frac{\dot{ϵ}\delta \dot{\rho }}{\lambda }\frac{\ddot{ϵ}\delta \rho }{\lambda }\right].$$
(13)
This expression is locally integrable in a neighborhood of the classical solutions , but unfortunately is not integrable globally in the full phase space. The non-integrability of $`\delta J`$ is, again, a consequence of the $`SL(2,R)`$ symmetry breaking. For instance, if one uses the symmetry-preserving boundary conditions for $`\eta `$ used in Ref. , $`\delta J`$ is globally integrable and Eq. (13) yields the form of $`J`$ obtained in that paper. Fortunately, we do not need to know $`J`$ globally, but is sufficient to know the form of $`J`$ near the $`\rho =1`$ configuration. The expression of $`J`$ can be further simplified considering only on-shell field configurations. Using the equations of motion (9) and (10), and expanding around the classical solutions, $`\rho =1+\overline{\rho }`$ in Eq. (13), one obtains at leading order in $`\overline{\rho }`$,
$$J(ϵ)=\frac{\eta _0}{\lambda }\left(\dot{ϵ}\dot{\overline{\rho }}\ddot{ϵ}\overline{\rho }\right)+ϵM,$$
(14)
Where $`M`$ is the mass (4), which on shell becomes constant. The charge (14) is defined up to an additive constant, we use a normalization such that $`J(ϵ=1)=M`$. In the following we will consider, for sake of semplicity, only the case $`M=0`$, i.e variations near the ground state. We are mainly interested in the value of the central charge of the Virasoro algebra, which is independent of $`M`$.
For generic $`ϵ`$ the charges $`J(ϵ)`$ are not conserved. The only conserved charge is obtained for $`ϵ=1`$. This fact is a consequence of the $`SL(2,R)T`$ symmetry breaking: only the charge associated with the residual symmetry $`T`$ is conserved. This fact makes it impossible to use the charges $`J(ϵ)`$ to give a realization of the Virasoro algebra, as it has been already shown in the canonical framework . To solve the problem, we proposed to introduce the time-integrated charges,
$$\widehat{J}(ϵ)=\frac{\lambda }{2\pi }_0^{\frac{2\pi }{\lambda }}𝑑tJ(ϵ).$$
(15)
The charges $`\widehat{J}`$ are now trivially conserved and generate the asymptotic symmetries of $`\mathrm{AdS}_2`$ trough the relation
$$\widehat{\delta _\omega J(ϵ)}=[\widehat{J}(ϵ),\widehat{J}(\omega )],$$
(16)
where the hat means overall time-integration as defined in Eq. (15).
The definitions (15) imply that $`J`$ is defined up to a total time-derivative. Using this freedom, we can always write,
$$J(ϵ)=2\frac{\eta _0}{\lambda }ϵ\ddot{\overline{\rho }}=ϵ\mathrm{\Theta }_{tt}.$$
(17)
Using the transformation laws (11), one gets
$$ϵ\delta _\omega \mathrm{\Theta }_{tt}=ϵ\left(\omega \dot{\mathrm{\Theta }}_{tt}+2\dot{\omega }\mathrm{\Theta }_{tt}\right)+c(ϵ,\omega ),$$
(18)
where $`c(ϵ,\omega )=\frac{\eta _0}{\lambda }(\ddot{ϵ}\dot{\omega }\ddot{\omega }\dot{ϵ})`$.
The previous equation tells us that $`\mathrm{\Theta }_{tt}`$ can be considered as the one-dimensional stress-energy tensor associated with the conformal symmetry, whereas $`\widehat{J}(ϵ)`$ are the charges generating a central extension of the Virasoro algebra. Notice that the central charge $`c`$ in general is $`\rho `$-dependent, the expression quoted above being the central charge evaluated near $`\rho =1`$. Expanding in Fourier modes,
$$\widehat{J}(ϵ)=\underset{m}{}a_mL_m,ϵ=\underset{m}{}a_me^{i\lambda mt},\mathrm{\Theta }_{tt}=\underset{m}{}L_me^{i\lambda mt},$$
(19)
and using Eqs. (16) and (18) one finds that the $`L_m`$ span a Virasoro algebra,
$$[L_m,L_n]=(mn)L_{m+n}+\frac{C}{12}m^3\delta _{m+n},$$
(20)
with central charge $`C=24\eta _0`$. The value of the central charge of the algebra is exactly the same as that found by us using a canonical realization of the symmetry and differs by a factor $`1/2`$ from that found by Navarro-Salas and Navarro . The origin of the mismatch is easily understood. As already pointed out, the authors of Ref. use different boundary conditions and a different transformation law for the field $`\eta `$. They use a non-scalar dilaton and $`\rho =1`$, identically, so that the $`SL(2,R)`$ isometry of $`\mathrm{AdS}_2`$ is not broken, and the origin of the central charge is completely different from our case. In Ref. the central charge arises as consequence of the anomalous term in the transformation law of $`\gamma _{tt}`$ (see Eq. (11)). This is very similar to the three-dimensional case . Their results hold only for a 2D dilaton gravity model that is not invariant under space-time diffeomorphisms. Therefore, the mismatch between statistical and thermodynamical entropy of 2D black hole discussed in our previous work still remains an open question. |
warning/0002/math0002239.html | ar5iv | text | # Polynomials on Schreier’s space
## 1 The weak Dunford-Pettis property
We say that a Banach space $`E`$ has the weak Dunford-Pettis property (wDPP, for short) if, given a uniformly weakly null sequence $`(x_n)E`$ and a weakly null sequence $`(\varphi _n)E^{}`$, we have $`lim\varphi _n(x_n)=0`$.
The space $`\mathrm{}_2`$ fails the wDPP since its unit vector basis is uniformly weakly null. Clearly, if $`E`$ has the DPP, then $`E`$ has the wDPP.
Denote by $`T`$ the dual of the original Tsirelson space $`T^{}`$ . Then the uniformly weakly convergent sequences in $`T`$ are norm convergent. Indeed, suppose $`(x_n)`$ is uniformly weakly convergent to $`xT`$ and $`x_nx\delta >0`$. Passing to a subsequence, we may assume that the sequence $`(x_nx)`$ is basic and equivalent to a subsequence of the unit vector basis $`(t_n)`$ of $`T`$ \[4, Proposition II.7\]. If $`A`$ is admissible, by the definition of the norm of $`T`$, we have:
$$\underset{iA}{}t_i\frac{1}{2}cardA$$
and so, $`(t_n)`$ has no uniformly weakly null subsequence, which yields a contradiction.
Therefore, $`T`$ enjoys the wDPP, but $`T^{}`$ does not since the unit vector basis of $`T^{}`$ is a Banach-Saks set. We conclude that the wDPP of a Banach space neither implies nor is implied by the wDPP of its dual.
The following simple remark will be useful:
###### Proposition 1
A Banach space $`E`$ has the wDPP if and only if whenever $`(x_n)E`$ is uniformly weakly null and $`(\varphi _n)E^{}`$ is weak Cauchy, we have $`lim\varphi _n(x_n)=0`$.
Proof. For the nontrivial part, if $`\varphi _n(x_n)\delta >0`$, we can find $`k_1<\mathrm{}<k_n<\mathrm{}`$ such that $`\left|\varphi _n\left(x_{k_n}\right)\right|<\delta /2`$. Then,
$$\delta \varphi _{k_n}\left(x_{k_n}\right)\left|\left(\varphi _{k_n}\varphi _n\right)\left(x_{k_n}\right)\right|+\left|\varphi _n\left(x_{k_n}\right)\right|$$
and the right hand side is less than $`\delta `$ for $`n`$ large enough, since the sequence $`\left(\varphi _{k_n}\varphi _n\right)`$ is weakly null. $`\mathrm{}`$
Denoting by $`𝒲𝒞o(E,F)`$ the space of all weakly compact (linear) operators from $`E`$ into the Banach space $`F`$, and by $`𝒞_w(E,F)`$ the space of all operators taking uniformly weakly null sequences in $`E`$ into norm null sequences in $`F`$, we have:
###### Proposition 2
The Banach space $`E`$ satisfies the wDPP if and only if, for all Banach spaces $`F`$, we have $`𝒲𝒞o(E,F)𝒞_w(E,F)`$.
Proof. Suppose $`E`$ has the wDPP and $`(x_n)E`$ is uniformly weakly null. Take $`L𝒲𝒞o(E,F)`$ with adjoint $`L^{}`$. Choose $`(\varphi _n)`$ in the unit ball of $`F^{}`$ such that $`\varphi _n(Lx_n)=Lx_n`$. There is a subsequence $`\left(\varphi _{n_k}\right)`$ such that $`\left(L^{}\varphi _{n_k}\right)`$ is weakly convergent. Hence, $`\varphi _n(Lx_n)=(L^{}\varphi _n)x_n0`$. Conversely, if $`E`$ fails the wDPP, we can find $`(x_n)`$ uniformly weakly null in $`E`$ and $`(\varphi _n)`$ weakly null in $`E^{}`$ such that $`\varphi _n(x_n)\delta >0`$. We define an operator $`L:Ec_0`$ by $`Lx:=(\varphi _n(x))`$. Then, $`L`$ is weakly compact but $`Lx_n|\varphi _n(x_n)|\delta >0`$ for all $`n`$. $`\mathrm{}`$
The following easy fact characterizes the reflexive Banach spaces with the wDPP:
###### Proposition 3
Let $`E`$ be a reflexive Banach space. Then $`E`$ has the wDPP if and only if every uniformly weakly null sequence in $`E`$ is norm null.
Proof. Suppose there is a uniformly weakly null sequence $`(x_n)E`$ with $`x_n=1`$. We can assume that $`(x_n)`$ is basic and the sequence of coefficient functionals $`(\varphi _n)`$ is weakly null in $`E^{}`$. Since $`\varphi _n(x_n)=1`$, we conclude that $`E`$ does not have the wDPP. The converse is clear. $`\mathrm{}`$
Recall that a Banach space $`E`$ has the Banach-Saks property if every bounded subset in $`E`$ is a Banach-Saks set. We then have:
###### Corollary 4
If $`E`$ has the Banach-Saks property and the wDPP, then $`E`$ is finite dimensional.
A space $`E`$ has the weak Banach-Saks property if every weakly null sequence in $`E`$ contains a subsequence whose arithmetic means converge. Equivalently \[18, Theorem 2.10\], every weakly null sequence has a subsequence which converges to zero uniformly weakly in $`E`$. The space $`L^1[0,1]`$ has the weak Banach-Saks property. The following result is clear:
###### Proposition 5
Assume $`E`$ has the weak Banach-Saks property. Then $`E`$ has the DPP if and only if $`E`$ has the wDPP.
We say that $`E`$ has the hereditary weak Dunford-Pettis property if every closed subspace of $`E`$ has the wDPP.
###### Proposition 6
A Banach space $`E`$ has the hereditary wDPP if and only if every normalized uniformly weakly null sequence in $`E`$ contains a subsequence equivalent to the $`c_0`$-basis.
Proof. Suppose that the uniformly weakly null sequence $`(x_n)E`$, $`x_n=1`$, has no subsequence equivalent to the $`c_0`$-basis. We can assume that $`(x_n)`$ is basic. Let $`(\varphi _n)[x_n]^{}`$ be the sequence of coefficient functionals, where $`[x_n]`$ denotes the closed linear span of the set $`\{x_n\}`$ in $`E`$. After taking a subsequence, we can assume that either $`(\varphi _n)`$ is equivalent to the $`\mathrm{}_1`$-basis or $`(\varphi _n)`$ is weak Cauchy \[10, Ch. XI\]. In the first case, we define an operator $`L:[x_n]c_0`$ by $`L(x):=(\varphi _n(x))`$. Clearly, $`L`$ is injective and has dense range. The adjoint $`L^{}:\mathrm{}_1[x_n]^{}`$ takes the unit vector basis of $`\mathrm{}_1`$ into the sequence $`(\varphi _n)`$ and has therefore closed range. Hence, $`L`$ is a surjective isomorphism, which contradicts our assumption. So, $`(\varphi _n)`$ must be weak Cauchy. Since $`\varphi _n(x_n)=1`$, the subspace $`[x_n]`$ fails to have the wDPP.
For the converse, it is enough to show that $`E`$ has the wDPP. Suppose it does not. Then we can find a uniformly weakly null sequence $`(x_n)E`$ and a weakly null sequence $`(\varphi _n)E^{}`$ such that $`\varphi _n(x_n)1`$ for all $`n`$. Passing to a subsequence, we can assume that $`(x_n)`$ is equivalent to the $`c_0`$-basis. Since the dual of $`c_0`$ has the Schur property, the restriction of $`(\varphi _n)`$ to the subspace $`[x_k]`$ is norm null, and we get a contradiction. $`\mathrm{}`$
###### Remark 7
This simple proof also shows that a Banach space $`E`$ has the hereditary DPP if and only if every normalized weakly null sequence in $`E`$ has a subsequence equivalent to the $`c_0`$-basis \[8, Proposition 2\]. From this we get that every infinite dimensional Banach space without a copy of either $`c_0`$ or $`\mathrm{}_1`$ contains a subspace without the DPP \[10, p. 254\]. The original proofs of these two results were based on a characterization of $`c_0`$’s unit vector basis that Elton obtained by using Ramsey’s theorem.
Our aim now is to show that Schreier’s space enjoys the hereditary wDPP.
###### Proposition 8
If $`(x_n)`$ is a uniformly weakly null sequence in $`S`$, then $`x_n_{\mathrm{}}0`$.
Proof. Let $`x_n=\left(x_n^i\right)_{i=1}^{\mathrm{}}`$. Since a set of $`\pm 1`$’s on an admissible set is a norm-one functional on $`S`$, given $`ϵ>0`$, there is $`N(ϵ)`$ such that
$$card\{n:\underset{iA}{}\left|x_n^i\right|ϵ\}N(ϵ)$$
for each admissible $`A`$. Suppose our statement fails; then we can find $`\delta >0`$ and two increasing sequences of indices $`(n_k),(l_k)`$ such that
$$\left|x_{n_k}^{l_k}\right|\delta \text{ for all }k.$$
The set $`A_m:=\{l_{m+1},\mathrm{},l_{2m}\}`$ is admissible for each $`m`$, and
$$card\{n:\underset{iA_m}{}\left|x_n^i\right|\delta \}m,$$
a contradiction which finishes the proof. $`\mathrm{}`$
The converse is not true. Indeed, take $`x_n:=(e_1+\mathrm{}+e_n)/n`$. The set $`A_k:=\{2^{k1},\mathrm{},2^k1\}`$ is admissible for each $`k`$. Denoting by $`\left(e_i^{}\right)`$ the unit vector basis of $`S^{}`$, the functional
$$\varphi _k:=\underset{i=2^{k1}}{\overset{2^k1}{}}e_i^{}S^{}$$
has norm one. Choosing $`n`$ so that $`2^{k2}+2^{k1}n2^k1`$, we have
$$\varphi _k(x_n)\frac{2^{k2}}{n}>\frac{2^{k2}}{2^k}=\frac{1}{4}.$$
Therefore, $`x_n_{\mathrm{}}0`$, but $`(x_n)`$ does not converge to zero uniformly weakly. The proof of the following result is essentially contained in . We give it for completeness.
###### Proposition 9
Let $`(x_n)`$ be a normalized sequence in $`S`$ such that $`x_n_{\mathrm{}}0`$. Then $`(x_n)`$ contains a subsequence equivalent to the $`c_0`$-basis.
Proof. Let us denote by $`supp(x)`$ the support of $`x`$. Passing to a subsequence and perturbing it with a null sequence, we can assume that $`\mathrm{max}supp(x_n)<\mathrm{min}supp(x_{n+1})`$, and
$`x_n_{\mathrm{}}{\displaystyle \frac{1}{2^n\mathrm{max}supp(x_{n1})}}.`$ (1)
Given $`x_{n_1},\mathrm{},x_{n_m}`$ and an admissible set $`A`$, we take $`k_0`$ to be the minimum value of $`k`$ such that $`Asupp\left(x_{n_k}\right)\mathrm{}`$. In particular, this implies that $`cardA\mathrm{max}supp\left(x_{n_{k_0}}\right)`$. Denoting $`x_n(i):=x_n^i`$, we have:
$`{\displaystyle \underset{iA}{}}\left|\left({\displaystyle \underset{k=1}{\overset{m}{}}}x_{n_k}\right)(i)\right|`$ $`=`$ $`{\displaystyle \underset{iA}{}}\left|\left({\displaystyle \underset{k=k_0}{\overset{m}{}}}x_{n_k}\right)(i)\right|`$
$`=`$ $`{\displaystyle \underset{k=k_0}{\overset{m}{}}}{\displaystyle \underset{iAsupp\left(x_{n_k}\right)}{}}\left|x_{n_k}(i)\right|`$
$``$ $`x_{n_{k_0}}+{\displaystyle \underset{k=k_0+1}{\overset{m}{}}}x_{n_k}_{\mathrm{}}cardA`$
$``$ $`x_{n_{k_0}}+{\displaystyle \underset{k=k_0+1}{\overset{m}{}}}2^{n_k}`$
$``$ $`2,`$
where we have used (1). Thus we have proved that
$$\underset{k=1}{\overset{m}{}}x_{n_k}2$$
and hence the series $`x_n`$ is weakly unconditionally Cauchy. Therefore, $`(x_n)`$ has a subsequence equivalent to the $`c_0`$-basis \[10, Corollary V.7\]. $`\mathrm{}`$
Combining the last two results with Proposition 6 yields:
###### Theorem 10
Schreier’s space $`S`$ has the hereditary wDPP.
We now show that the dual $`S^{}`$ of Schreier’s space fails the wDPP. The next result follows the lines of \[17, Example 13,(H)\].
###### Proposition 11
Let $`(\varphi _n)`$ be a normalized block basis of the unit basis of $`S^{}`$ such that $`\varphi _n_{\mathrm{}}0`$. Then $`(\varphi _n)`$ contains a subsequence equivalent to the $`\mathrm{}_1`$-basis.
Proof. Let $`(x_n)`$ be a sequence in $`S`$ such that $`x_n<2`$, $`supp(x_n)=supp(\varphi _n)`$ and $`\varphi _n(x_n)=1`$ for every $`n`$.
First we select $`n_1`$ such that $`\mathrm{min}supp\left(\varphi _{n_1}\right)>2^2`$ and $`\varphi _{n_1}_{\mathrm{}}<2^4`$. Since $`x_{n_1}<2`$, the set
$$A_1=\{i:\left|x_{n_1}(i)\right|2^1\}$$
has fewer than $`2^2`$ elements. We define $`y_{n_1}(i)=0`$ if $`iA_1`$ and $`y_{n_1}(i)=x_{n_1}(i)`$ otherwise, and obtain $`y_{n_1}S`$ such that $`y_{n_1}<2`$, $`y_{n_1}_{\mathrm{}}<2^1`$ and
$$\left|\varphi _{n_1}\left(y_{n_1}\right)\right|\varphi _{n_1}(x_{n_1})\left|\varphi _{n_1}(y_{n_1}x_{n_1})\right|>12(2^2)2^4=2^1.$$
Next we select $`n_2>n_1`$ such that $`\mathrm{min}supp\left(\varphi _{n_2}\right)>2^3`$ and $`\varphi _{n_2}_{\mathrm{}}<2^5`$. Since $`x_{n_2}<2`$, the set
$$A_2=\{i:\left|x_{n_2}(i)\right|2^2\}$$
has fewer than $`2^3`$ elements. We define $`y_{n_2}(i)=0`$ if $`iA_2`$ and $`y_{n_2}(i)=x_{n_1}(i)`$ otherwise, and obtain $`y_{n_2}S`$ such that $`y_{n_2}<2`$, $`y_{n_2}_{\mathrm{}}<2^2`$ and
$$\left|\varphi _{n_2}\left(y_{n_2}\right)\right|\varphi _{n_2}(x_{n_2})\left|\varphi _{n_2}(y_{n_2}x_{n_2})\right|>12(2^3)2^5=2^1.$$
In this way we get a subsequence $`\left(\varphi _{n_j}\right)`$ and a sequence $`\left(y_{n_j}\right)S`$ such that $`\left|\varphi _{n_j}\left(y_{n_j}\right)\right|>2^1`$, $`y_{n_j}<2`$ and $`y_{n_j}_{\mathrm{}}<2^j`$. Passing to a subsequence we can assume by Proposition 9 that $`\left(y_{n_j}\right)`$ is equivalent to the $`c_0`$-basis, from which it easily follows that $`\left(\varphi _{n_j}\right)`$ is equivalent to the $`\mathrm{}_1`$-basis. $`\mathrm{}`$
###### Proposition 12
The dual $`S^{}`$ of Schreier’s space $`S`$ has the weak Banach-Saks property.
Proof. Let $`(\varphi _n)`$ be a normalized weakly null sequence in $`S^{}`$. Passing to a subsequence we can assume that $`(\varphi _n)`$ is equivalent to a block basis of the unit basis. We have that $`((\varphi _1+\mathrm{}+\varphi _n)/n))`$ is a weakly null sequence and $`(\varphi _1+\mathrm{}+\varphi _n)/n_{\mathrm{}}0`$. If $`(\varphi _1+\mathrm{}+\varphi _n)/n`$ does not converge to $`0`$, passing to a subsequence, it follows from Proposition 11 that $`((\varphi _1+\mathrm{}+\varphi _n)/n))`$ contains a subsequence equivalent to the $`\mathrm{}_1`$-basis, a contradiction. $`\mathrm{}`$
###### Corollary 13
The dual $`S^{}`$ of Schreier’s space does not have the wDPP.
## 2 Applications to polynomials
In this Section, we describe the $`𝒫`$-convergence of sequences in $`S`$, obtaining thereby some polynomial properties of this space, and characterize the Banach-Saks sets in it.
We shall use the fact that $`S`$ may be algebraically embedded in $`\mathrm{}_2`$, and that the natural inclusion $`j:S\mathrm{}_2`$ is continuous. To see this, take $`x:=(x_i)S`$, $`x_S=1`$, and call $`y:=(y_i)`$ the sequence $`(|x_i|)`$, reordered in a nonincreasing way. Then $`y_2=x_2`$, and $`y_S1`$. This implies $`y_{2k1}k^1`$ for each $`k`$. Therefore,
$$y_2^2=\underset{i=1}{\overset{\mathrm{}}{}}y_i^21+1+\frac{1}{2^2}+\frac{1}{2^2}+\frac{1}{3^2}+\frac{1}{3^2}+\mathrm{}=\frac{\pi ^2}{3},$$
from which $`j\pi /\sqrt{3}`$.
As a consequence, $`P(x):=x_2^2`$ defines a $`2`$-homogeneous polynomial on $`S`$.
###### Proposition 14
Let $`(x_n)`$ be a sequence in $`S`$. The following assertions are equivalent:
(a) $`(x_n)`$ is $`𝒫`$-null;
(b) $`(x_n)`$ is bounded in $`S`$ and $`x_n_20`$;
(c) $`(x_n)`$ is bounded in $`S`$ and $`x_n_{\mathrm{}}0`$.
Proof. (a) $``$ (b) since $`P(x):=x_2^2`$ is a polynomial on $`S`$.
(b) $``$ (c) is clear.
(c) $``$ (a). It is enough to show that $`(x_n)`$ has a $`𝒫`$-null subsequence. If $`infx_n>0`$, then there is a subsequence of $`(x_n)`$ equivalent to the $`c_0`$-basis (Proposition 9) and so $`𝒫`$-null, since the $`c_0`$-basis is $`𝒫`$-null. If $`infx_n=0`$, then there is a norm null subsequence, which is $`𝒫`$-null a fortiori. $`\mathrm{}`$
A Banach space has the hereditary polynomial DPP if every closed subspace has the polynomial DPP.
###### Theorem 15
The space $`S`$ has the hereditary polynomial DPP.
Proof. By Propositions 14 and 9, every normalized $`𝒫`$-null sequence in $`S`$ contains a subsequence equivalent to the $`c_0`$-basis. Obvious modifications in the “if” part of the proof of Proposition 6 yield the result. $`\mathrm{}`$
It is shown in that, given two $`𝒫`$-null sequences $`(x_n)`$, $`(y_n)`$ in a space with the polynomial DPP, the sequence $`(x_n+y_n)`$ is $`𝒫`$-null. A Banach space where this is not true was recently found by Castillo et al. .
###### Proposition 16
Let $`A`$ be a subset of $`S`$. The following assertions are equivalent:
(a) $`A`$ is a Banach-Saks set;
(b) $`A`$ is relatively $`𝒫`$-compact;
(c) $`A`$ is relatively weakly compact in $`S`$, and relatively compact as a subset of $`\mathrm{}_{\mathrm{}}`$.
Proof. (a) $``$ (b). Let $`A`$ be a Banach-Saks set. Given a sequence $`(x_n)A`$, passing to a subsequence, we may assume that $`(x_n)`$ converges to some $`x`$ uniformly weakly in $`S`$. Then, $`(x_nx)`$ has a subsequence which is either norm null or equivalent to the $`c_0`$-basis. In both cases, $`(x_n)`$ is $`𝒫`$-convergent to $`x`$.
(b) $``$ (c). If $`A`$ is relatively $`𝒫`$-compact, it is relatively weakly compact. Moreover, given a sequence $`(x_n)A`$, we can assume that $`(x_nx)`$ is $`𝒫`$-null for some $`x`$. By Proposition 14, $`x_nx_{\mathrm{}}0`$ and so $`A`$ is relatively compact as a subset of $`\mathrm{}_{\mathrm{}}`$.
(c) $``$ (a). Choose a sequence $`(x_n)A`$. We may assume that $`(x_n)`$ is weakly convergent to some $`x`$ and $`x_nx_{\mathrm{}}0`$. Passing to a subsequence, we have either $`x_nx0`$ or, by Proposition 9, $`(x_nx)`$ is equivalent to the $`c_0`$-basis and is therefore uniformly weakly null. $`\mathrm{}`$
###### Corollary 17
If $`A`$ is a Banach-Saks set in $`S`$, then the absolutely convex closed hull of $`A`$ is a Banach-Saks set.
The following two properties were introduced in and studied by various authors (see, e.g., ).
(a) A Banach space $`E`$ has property (P) if, given two bounded sequences $`(u_n),(v_n)`$ in $`E`$ such that $`P(u_n)P(v_n)0`$ for every $`P𝒫(^kE)`$ and all $`k`$, it follows that the sequence $`(u_nv_n)`$ is $`𝒫`$-null. Every superreflexive space and every space with the DPP have property (P). A Banach space failing to have property (P) has been found by Castillo et al. .
(b) A Banach space $`E`$ has property (RP) if, given two bounded sequences $`(u_n),(v_n)`$ in $`E`$ such that the sequence $`(u_nv_n)`$ is $`𝒫`$-null, it follows that $`P(u_n)P(v_n)0`$ for every $`P𝒫(^kE)`$ and all $`k`$. Every $`\mathrm{\Lambda }`$-space and every predual of a Banach space with the Schur property have property (RP). The spaces $`L_1[0,1]`$, $`C[0,1]`$ and $`L_{\mathrm{}}[0,1]`$ fail to have property (RP) .
We now show that $`S`$ has property (P) and fails property (RP).
###### Proposition 18
The space $`S`$ fails property (RP).
Proof. Consider the vectors
$$v_n:=e_n;u_n:=e_n+2^{1n}\left(e_{2^{n1}}+\mathrm{}+e_{2^n1}\right).$$
Then $`u_nv_n_{\mathrm{}}0`$ and so $`(u_nv_n)`$ is $`𝒫`$-null in $`S`$. Define
$$P(x):=\underset{n=1}{\overset{\mathrm{}}{}}x_n^2\left(\underset{k=2^{n1}}{\overset{2^n1}{}}x_k\right),\text{ for }x=(x_n)S.$$
Since
$$|P(x)|x_Sx_2^2\frac{\pi ^2}{3}x_S^3,$$
we get that $`P𝒫(^3S)`$. We have $`P(v_n)=0`$ and $`P(u_n)=1`$ for all $`n>1`$. $`\mathrm{}`$
In the above proof, we need a polynomial of degree greater than or equal to three. Indeed, if $`P𝒫(^2S)`$ and $`(u_n),(v_n)S`$ are bounded with $`(u_nv_n)`$ $`𝒫`$-null, denoting $`w_n:=u_nv_n`$, we have
$$P(u_n)P(v_n)=P(w_n+v_n)P(v_n)=2\widehat{P}(w_n,v_n)+P(w_n),$$
where $`\widehat{P}`$ is the symmetric bilinear form associated to $`P`$. Let $`\overline{P}:SS^{}`$ be the operator defined by $`\overline{P}(x)(y):=\widehat{P}(x,y)`$. Since $`S`$ has an unconditional basis and contains no copy of $`\mathrm{}_1`$, the space $`S^{}`$ has an unconditional basis and is weakly sequentially complete. Therefore, every operator from $`S`$ into $`S^{}`$ is weakly compact. Passing to a subsequence, we can assume that $`(w_n)`$ is uniformly weakly null. Since $`S`$ has the wDPP, $`\overline{P}(w_n)0`$. Hence, $`\widehat{P}(w_n,v_n)=\overline{P}(w_n)(v_n)0`$. Clearly, $`P(w_n)0`$ and so $`P(u_n)P(v_n)0`$.
###### Proposition 19
The space $`S`$ enjoys property (P).
Proof. Let $`(u_n),(v_n)S`$ be bounded sequences such that $`(u_nv_n)`$ is not $`𝒫`$-null. We wish to find $`Q𝒫(^kS)`$ for some $`k`$ so that $`(Q(u_n)Q(v_n))`$ does not tend to zero. By $`u_n^i`$ and $`v_n^i`$ we shall denote the $`i`$th coordinate of $`u_n`$ and $`v_n`$, respectively.
If $`(u_nv_n)`$ is not weakly null, then $`\varphi (u_n)\varphi (v_n)\to ̸0`$ for some $`\varphi S^{}`$. It is enough to take $`Q:=\varphi `$.
If $`(u_nv_n)`$ is weakly null, passing to a subsequence and perturbing it by a norm null sequence, we can assume that $`(u_nv_n)`$ is a block basis:
$$u_nv_n=\underset{i=k_n}{\overset{l_n}{}}a_ie_i.$$
Take $`p_n`$ with $`k_np_nl_n`$ and $`\left|a_{p_n}\right|=u_nv_n_{\mathrm{}}`$. We know that $`u_nv_n_{\mathrm{}}`$ does not go to zero. Passing to a subsequence, we may assume:
$$v_n^{p_n}v;u_n^{p_n}u;uv.$$
Let $`P𝒫(^2S)`$ be given by
$$P(x):=\underset{i=1}{\overset{\mathrm{}}{}}\left(x_{p_i}\right)^2.$$
If $`P(u_n)P(v_n)\to ̸0`$, we are done. If
$$0=lim\left[P(u_n)P(v_n)\right]=lim\left[\left(u_n^{p_n}\right)^2\left(v_n^{p_n}\right)^2\right]=u^2v^2=(uv)(u+v),$$
we have $`u=v=\alpha `$, for some $`\alpha 0`$. Defining $`Q𝒫(^3S)`$ by
$$Q(x):=\underset{i=1}{\overset{\mathrm{}}{}}\left(x_{p_i}\right)^3,$$
we have $`Q(u_n)Q(v_n)2\alpha ^30`$, and the proof is finished. $`\mathrm{}`$
###### Proposition 20
Let $`(x_n),(y_n)S`$ be $`𝒫`$-null sequences. Then:
(a) the set $`\{x_ny_n\}`$ is a Banach-Saks set in $`S_\pi S`$;
(b) the sequence $`(x_ny_n)`$ is $`𝒫`$-null in $`S_\pi S`$.
Proof. (a) Since $`(x_n)`$ and $`(y_n)`$ have subsequences equivalent to the $`c_0`$-basis, it is enough to show that $`(e_ne_n)`$ is uniformly weakly null in $`c_0_\pi c_0`$. Take $`L(c_0_\pi c_0)^{}`$, which may be viewed as an operator from $`c_0`$ into $`\mathrm{}_1`$. Since the series $`e_n`$ is weakly unconditionally Cauchy, using \[14, Theorem 2\], we can find $`C>0`$ such that $`|Le_n,e_n|C`$ whenever $`L1`$. Therefore, given $`ϵ>0`$, choosing $`N`$ with $`NC/ϵ`$, we have
$$card\{n:|Le_n,e_n|ϵ\}N$$
if $`L1`$, and the result is proved.
(b) Since $`S`$ has the polynomial DPP, part (b) follows from \[3, Theorem 2.1\]. $`\mathrm{}`$
As a consequence, if $`A,BS`$ are Banach-Saks sets, then $`AB`$ is a Banach-Saks set in $`S_\pi S`$.
| Manuel González | Joaquín M. Gutiérrez |
| --- | --- |
| Departamento de Matemáticas | Departamento de Matemáticas |
| Facultad de Ciencias | ETS de Ingenieros Industriales |
| Universidad de Cantabria | Universidad Politécnica de Madrid |
| 39071 Santander (Spain) | C. José Gutiérrez Abascal 2 |
| e-mail gonzalem@ccaix3.unican.esunican | 28006 Madrid (Spain) |
| | e-mail gutierrezj@member.ams.org |
file poss.tex |
warning/0002/hep-lat0002018.html | ar5iv | text | # 1. Introduction
## 1. Introduction
In a series of papers it has been shown that the technique of stochastic quantisation introduced by Parisi and Wu can be implemented as a practical algorithm which enables to reach unprecedented high orders in lattice perturbation theory (e.g. $`\alpha ^8`$ in the plaquette expectation value in 4-D $`SU(3)`$ lattice gauge theory). Evaluating perturbative expansion coefficients by a Monte Carlo technique opens the back–door to a number of errors, both statistical and systematic, which should be understood and, hopefully, kept under control. It has been shown in Ref. how systematic errors due to the finite volume can be estimated, putting a limit on the perturbative order reachable on a given lattice size. The aim of the present paper is to present a rather detailed analysis of statistical errors, which present some novel features with respect to the ordinary practice in lattice gauge Monte Carlo. This work was triggered by the observation made by M. Pepe of unexpected large discrepancies with respect to known values in the perturbative coefficients of the non–linear $`O(3)\sigma `$model. Starting from this observation (“Pepe effect”) we performed a systematic study of simple models with a small number of degrees of freedom trying to trace the origin of the discrepancies and to resolve them in a reliable way. The crucial fact that our analysis has uncovered is the following: the statistical nature of the processes which enter into the calculation of perturbative coefficients is rapidly deviating from normality as we increase the perturbative order, i.e. the distribution function of a typical coefficient estimator is strongly non–Gaussian, exhibiting a large skewness and a long tail; as a result very rare events give a substantial contribution to the average. A simple minded statistical analysis based on the assumption of normality may grossly fail to identify the confidence intervals; some non–parametric statistical analysis, like the bootstrap method, is necessary to assess the statistical error and provide reliable confidence intervals. This idea will be shown at work in the analysis of a lattice toy model (Weingarten’s “pathological” model ) which nevertheless presents many features of interest.
In view of this analysis one should of course worry about the results obtained in Lattice Gauge Theory (LGT). Our main conclusion with this respect is that one is not going to jeopardize the picture we drew in our previous works. As amazing as it can appear at first sight, the application of the method to a by far less trivial model stands on a by far firmer ground. We shall in fact show how the distribution function of coefficient estimators for a typical LGT observable does not exhibit the strongly non–Gaussian nature we find in simpler models.
The content of the paper is organized as follows: in sec.2 we recall the basis facts about Parisi-Wu stochastic technique applied to the numerical calculation of perturbative coefficients in quantum field theory on the lattice (hereafter the “Parisi-WU process” or PW-process for short). In sec.3 we discuss the probability distributions of the PW-process; we shall argue that the customary asymptotic analysis of “sum of identically distributed independent random variables” does not really help at high orders; in this case we shall present numerical evidence showing that the distribution functions still present large deviations from normality, in particular a whole window where the density presents a power–law rather than Gaussian behaviour. We present some detail on the algorithmic implementation of the PW-process in Sec.4 and some numerical results in Sec.5. A discussion of an alternate route is given in Sec.6, based on Girsanov’s formula. We then show (sec.7) how a bootstrap analysis can be very effective in estimating confidence intervals. Finally, in sec. 8, some evidence is produced that both convergence time of the processes and statistical errors are under control for LGT. We present our conclusions in sec. 9. Some details on the bootstrap method and a formal analysis of convergence of perturbative correlation functions are given in appendix.
## 2. Stochastic perturbation theory
Starting from Parisi and Wu’s pioneering paper, stochastic equations have been used in various forms to investigate quantum field–theoretical models, both perturbatively and non–perturbatively. In particular a Langevin approach can be used as a proposal step subject to a Metropolis check to implement a non–perturbative MonteCarlo for Lattice Gauge Theories. We are concerned here with another approach which makes use of the Langevin equation to calculate weak coupling expansions. The idea is very simple: we start from the Langevin equation (let us focus our attention on a simple scalar field $`\phi `$)
(1)
$$\frac{\phi (𝒙,t)}{t}=\frac{S}{\phi (𝒙,t)}+\eta (𝒙,t)$$
where $`\eta (𝒙,t)`$ is the standard white–noise generalized process; assuming that the action $`S`$ is splitted into a free $`S_0`$ and an interaction part $`gS_1`$, we expand the process $`\phi `$ into powers in the coupling constant
(2)
$$\phi (𝒙,t)=\underset{n=0}{\overset{\mathrm{}}{}}g^n\phi _n(𝒙,t).$$
The Langevin equation is thus translated into a hierarchical system of partial differential equations
(3) $`{\displaystyle \frac{\phi _0(𝒙,t)}{t}}`$ $`=`$ $`{\displaystyle \frac{S_0}{\phi _0(𝒙,t)}}+\eta (𝒙,t)`$
(4) $`{\displaystyle \frac{\phi _n(𝒙,t)}{t}}`$ $`=`$ $`G_0^1\phi _n(𝒙,t)+D_n(\phi _0,\mathrm{},\phi _{n1}),\mathrm{for}n1`$
$`G_0`$ being the free propagator and $`D_n`$ representing source terms which can be expressed in terms of higher functional derivatives of the interaction term $`S_1`$. Notice that the source of randomness is confined to the first (free–field) equation; the system can be truncated at any order $`n`$ due to its peculiar structure ($`D_n`$ depends only on $`\phi _m`$ with $`m<n`$).
The system can be used to generate a diagrammatic expansion, as Parisi and Wu did for gauge theories, in terms of the free propagator $`G_0`$; or, it can be studied numerically by simulating the white–noise process, as we have discussed in Ref.. Any given observable $`𝒪(\phi )`$ can be expanded in powers of the coupling constant
(5)
$$𝒪(\phi )𝒪\left(\underset{n0}{}g^n\phi _n\right)=\underset{n0}{}g^n𝒪_n(\phi _0,\mathrm{},\phi _n)$$
and its expectation value is given by
(6)
$$𝒪=\underset{n0}{}g^n𝒪_n=\underset{n0}{}c_n(𝒪)g^n.$$
The operator $`𝒪_n`$ is therefore an unbiased estimator of the $`n`$th expansion coefficient of $`𝒪`$.
It is rather straightforward to implement this idea in a practical algorithm, once the theory has been formulated on a space–time lattice. The application to gauge theory was presented in Ref. using the Langevin algorithm of Batrouni et al. Finite size errors where studied in a subsequent paper .
The key problem we want to discuss in the present paper consists in finding a reliable way to estimate the statistical errors in the measure of $`𝒪_n`$. Admittedly, it would be desirable to have some analytic information on the nature of the multidimensional coupled stochastic processes (3). As a substitute we choose to study some toy models where we can perform high statistics calculations and compare the results with the exact coefficients.
## 3. Toy models.
We have extensively studied the application of numerical stochastic perturbation theory to the following simple models:
1. quartic random variable: $`S(\phi ,g)=\frac{1}{2}\phi ^2+\frac{1}{4}g\phi ^4,\phi `$.
2. dipole random variable: $`S(\phi ,g)=[1\mathrm{cos}(g\phi )]/g^2,\phi (\pi ,\pi )`$.
3. Weingarten’s “pathological model” :
$$S=\frac{1}{2}\underset{\mathrm{}}{}\phi _{\mathrm{}}^2+\frac{1}{4}g\underset{𝔭}{}\phi _𝔭$$
$$\phi _𝔭=\underset{\mathrm{}𝔭}{}\phi _{\mathrm{}},\phi _{\mathrm{}},$$
where $`\mathrm{}`$ runs over links and $`𝔭`$ over plaquettes in a simple $`n`$-dimensional cell of a cubic lattice. We consider $`n=2,3,4`$.
The integration measure is $`d\phi \mathrm{exp}(S)`$ with the ordinary Lebesgue measure $`d\phi `$ over all degrees of freedom; the expectation value of any field observable is then given by
$$𝒪=Z^1d\phi 𝒪[\phi ,g]\mathrm{exp}\{S(\phi ,g)\}.$$
The calculation of the weak coupling expansion for the typical observables can be easily performed, except for for the model iii) in four dimensions, where we have been unable to go beyond the $`5^{th}`$ order. We have for instance for the “quartic” integral
$`\phi ^4`$ $`=\frac{1}{4}{\displaystyle \frac{d\mathrm{log}(Z(g))}{dg}}=13g24g^2297g^34896g^4100278g^4`$
$`2450304g^569533397g^62247492096g^781528066378g^8+O(g^9),`$
for the “dipole” integral
$`\mathrm{cos}(gx)=1{\displaystyle \frac{g^2}{2}}{\displaystyle \frac{g^4}{8}}+{\displaystyle \frac{17g^6}{96}}{\displaystyle \frac{29g^8}{512}}{\displaystyle \frac{251g^{10}}{61440}}+O(g^{12}),`$
and finally
$`\phi _𝔭`$ $`{\displaystyle \frac{d\mathrm{log}(Z_3(g))}{dg}}`$
$`={\displaystyle \frac{g}{4}}{\displaystyle \frac{13g^3}{64}}{\displaystyle \frac{103g^5}{256}}{\displaystyle \frac{23797g^7}{16384}}{\displaystyle \frac{2180461g^9}{262144}}{\displaystyle \frac{72763141g^{11}}{1048576}}{\displaystyle \frac{13342715521g^{13}}{16777216}}+O(g^{14})(n=2)`$
$`={\displaystyle \frac{g}{4}}{\displaystyle \frac{17g^3}{64}}{\displaystyle \frac{595g^5}{1024}}{\displaystyle \frac{34613g^7}{16384}}{\displaystyle \frac{3059191g^9}{262144}}{\displaystyle \frac{388561373g^{11}}{4194304}}{\displaystyle \frac{67903544647g^{13}}{67108864}}+O(g^{14})(n=3)`$
$`={\displaystyle \frac{g}{4}}{\displaystyle \frac{21g^3}{64}}{\displaystyle \frac{411g^5}{512}}+O(g^7)(n=4)`$
for Weingarten’s model. We have limited ourselves to models involving real random variables which make it very simple to implement the stochastic differential equation Eq.(3). To reach high orders we used a symbolic language to build the right-hand side (i.e. the terms $`D_n(\phi _0,\mathrm{},\phi _{n1})`$). In the numerical experiments $`n`$ is in the range $`10`$ to $`16`$.
For the sake of brevity we shall discuss our numerical experiments for the third model, which is the nearest to lattice gauge theory, but the qualitative aspects of the results are indistinguishable in the other models.
Of course one should be aware that Weingarten’s model is rather peculiar; its action is not bounded from below and it may make sense only in Minkowski space. From the perturbative viewpoint, however, it is perfectly admissible and for imaginary $`g`$ the integrals converge absolutely. So this is an example where perturbation theory for a model living in Minkowski space only can be computed by the stochastic method even if the model does not allow for a Euclidean formulation.
## 4. Algorithms
Let us describe in some detail the numerical implementation of the stochastic differential equations (3). The free field $`\phi _0`$ can be computed exactly (as a discrete Markov chain) since it is a Ornstein-Uhlenbeck process . For a time increment $`\tau `$ we can write
(7)
$$\phi _0(t+\tau )=e^\tau \phi _0(t)+\sqrt{1\mathrm{exp}(2\tau )}N(0,1)$$
where $`N(m,\sigma )`$ is a normal deviate with mean $`m`$ and variance $`\sigma `$. By integrating over the time interval $`[t,t+\tau ]`$ we get
(8) $`\phi _n(t+\mathrm{\Delta }t)=e^\tau \phi _n(t)+{\displaystyle _t^{t+\tau }}\mathrm{exp}(t^{}t\tau )D_n(\phi _0(t^{}),\mathrm{},\phi _{n1}(t^{}))𝑑t^{}`$
which can be approximated, for instance, by the trapezoidal rule. Due to the peculiar nesting of the differential equations it is not necessary to perform the usual “predictor” step to implement the trapezoidal rule, which results in a faster algorithm.
The bias introduced by the approximation is of order $`O(\tau ^\nu )`$, with $`\nu 2`$, which should be taken into account as it is usual with the numerical implementation of Langevin equation. In this case however we cannot apply a Metropolis step, as it can be done in the non–perturbative case. The trapezoidal rule and the absence of bias on the free field, however, conspire to keep the finite-$`\tau `$ error to a small value.
The numerical experiments have been performed by running several independent trajectories in parallel (typically $`10^2`$ to $`10^4`$); a crude estimate of autocorrelation gives a value near to $`\tau _{int}=\frac{1}{2}`$ which is the relaxation time of the Ornstein-Uhlenbeck process which drives the whole system. Hence we measure the observable $`𝒪_n`$ every $`n_{skip}`$ steps with $`n_{skip}\tau 1`$.
The pseudo-random numbers needed to generate the normal deviates are produced through two different algorithms, the lagged–Fibonacci algorithm described by Knuth and Luscher’s algorithm . They do not give appreciably different results in our experiments.
## 5. Results
We report the results of experiments carried out on the model iii) since it permits a graduation in dimensionality. Other models have been studied in detail and they show the same general pattern. The first three figures are organized as follows: the first plot is a log-log histogram of the raw data for a given observable (i.e. $`𝒪_{11}`$); the third plot represents $`\frac{1}{T}_0^T𝒪_n𝑑t`$ averaged over 1000 independent histories. The middle plot represents a blow up of the same average together with some bootstrap sample which allow to estimate confidence intervals, as we describe in some detail later on.
The main observation consists in the fact that the histograms corresponding to high order coefficients deviate substantially from a Gaussian distribution. We do not have a complete characterization of the densities; however the study of these histograms sheds some light on the problem. We observe a rather large window where the densities are dominated by a power–law fall–off. Above a certain order, typically 10-14, the power is approximately $`x^{2\nu }`$ with small $`\nu `$, which means that there are large deviations from the average which occur as very rare events but contribute to the average. An example is found in Fig. 1 where the time average shows a sharp kink at $`t1.12e4`$ and a relaxation thereof. In the first picture the log-log–histogram is contrasted with a $`x^2`$ and $`x^3`$ behaviour. Several runs on the same model $`iii_{3D})`$ are necessary to get a satisfactory estimate of the perturbative coefficient (here $`c_{11}=388561373/419430492.64`$). Another experiment gives more regular histories (see Fig. 2).
At order 15 we encounter a similar pattern. We observe that the large jumps at $`t=1e4`$ and $`t=1.7e4`$ both contribute to reach a value rather close to the exact one. Were it not for the histograms, which, by exhibiting a large $`x^{2\nu }`$ window, warn about very slow convergence, one would be tempted to conclude that a plateau had been reached well before $`t10^4`$.
On the contrary, low order coefficients are rather close to Gaussian behaviour and the convergence is very fast (see Fig. 4):
We attempted an analytical characterization of the distribution functions $`P_n(𝒪_n<x)`$ for the processes $`𝒪_n`$. The problem is rather tricky and no explicit formula was reached. One can only show that if $`𝒪`$ is a generic correlation of fields, and $`P_n(𝒪_1,\mathrm{},𝒪_n,t)`$ the distribution function of its first $`n`$ perturbative coefficients at a Langevin time $`t`$, then a limit distribution (for $`t\mathrm{}`$) always exists, and that all its moments are finite. In fact this reduces to show that all correlation functions
(9)
$$\underset{j}{}\phi _{p_j}(t)$$
are finite in the limit $`t\mathrm{}`$. Details can be found in App. A.
Since all moments turn out to be finite, one is tempted to use general theorems regarding the sum of independent identically distributed random variables . These could in principle make it unnecessary to have a detailed knowledge of $`P_n`$, since we are averaging over a large number of independent histories and the outcome should be a Gaussian distribution with corrections which can be parameterized and fitted to the data. According to Chebyshev and Petrov
(10)
$$_{\mathrm{}}^x[P_n^{(k)}(x)G(x)]𝑑x=\frac{e^{x^2/2}}{\sqrt{2\pi }}(\frac{Q_1^{(k)}(x)}{\sqrt{n}}+\frac{Q_2^{(k)}(x)}{n}+\mathrm{})$$
where the $`Q_i(x)`$ are polynomials, $`G(u)`$ is the normal distribution and $`k`$ is the number of independent histories and the convergence is even uniform on $`x`$. One has for instance
$`Q_1^{(k)}(x)=\lambda _3^{(k)}{\displaystyle \frac{(1x^2)}{6}};Q_2^{(k)}(x)=\lambda _4^{(k)}x(3x^2);`$
$`\lambda _q^{(k)}={\displaystyle \frac{(\varphi ^{(k)})^q_c}{\sigma _k^q}};\sigma _k^2=(\varphi ^{(k)})^2.`$
Unfortunately the expansion is, at best, an asymptotic one and no useful estimate on the error was obtained on the basis of these formulas. This fact has triggered our interest on non–parametric methods of analysis which will be presented in sec. 7.
## 6. Girsanov’s formula
The huge jumps in processes $`𝒪_n`$ at high $`n`$ could in principle be caused by numerical inaccuracy, and for some time we suspected that this could in fact happen. It was soon realized that the large deviations are in fact needed to reach the right average in the long runs, but it was nevertheless reassuring that an alternative method with zero autocorrelation and a totally different algorithm gave the same pattern of configurations. As a matter of fact the method turns out to be less efficient than the direct solution of the system given by Eq.(3), but it is important in order to show that the peculiar properties of the $`𝒪_n`$ histories are not an algorithmic artifact.
The method applies Cameron-Martin-Girsanov formula (see e.g., , §3.12)
(11)
$$E\left(\mathrm{\Phi }\left(𝒙(T)\right)|_w\right)=E\left(\mathrm{\Phi }\left(𝒚(T)\right)\mathrm{exp}(\xi _T)|_w\right)$$
where
$`d𝒙(t)`$ $`=`$ $`A𝒙dt+𝒃(𝒙,t)dt+d𝒘(t)`$
$`d𝒚(t)`$ $`=`$ $`A𝒚dt+d𝒘(t),(y(0)=x(0))`$
$`\xi (T)`$ $`=`$ $`{\displaystyle _0^T}𝒃(𝒚(\tau ),\tau )𝑑𝒘(\tau )\frac{1}{2}{\displaystyle _0^T}𝒃(𝒚(\tau ),\tau )^2𝑑\tau `$
and $`E(.|_w)`$ denotes the average with respect to the standard Brownian process $`𝒘`$. We apply the formula like in Ref., hence $`A`$ is given by the free inverse propagator and $`𝒃(𝒙,t)`$ is proportional to the coupling constant $`g`$; therefore we can explicitly expand the “exponential martingale” $`\xi (t)`$ as a power series in $`g`$ and get an explicit characterization of the expansion coefficients:
(12)
$$𝒪_n=\underset{T\mathrm{}}{lim}E\left(H_n(I_1(T))I_2^{\frac{n}{2}}(T)/n!\right)$$
where
$`I_1(T)`$ $`=`$ $`\frac{1}{2}{\displaystyle _0^T}𝒃(𝒚(\tau ),\tau )𝑑𝒘(\tau )`$
$`I_2(T)`$ $`=`$ $`\frac{1}{2}{\displaystyle _0^T}𝒃(𝒚(\tau ),\tau )𝑑\tau `$
and $`H_n`$ are Hermite polynomials. This representation of the perturbative coefficients is completely independent on our expansion Eq.(3). It has been implemented as a numerical algorithm and used to estimate the perturbative expansion for model $`i)`$. Each iteration consists in selecting a normally distributed starting point $`x(0)`$ and following the evolution up to a time $`T`$ where transient effects are sufficiently damped; since $`A=1`$ in our model, $`T1`$ is adequate. The samples are statistically independent, modulo the quality of the random number generator. By monitoring the average of $`\xi (t)`$, which should be exactly one, we have a handle on the accuracy of the algorithm (finite time step and statistics).
The application of Girsanov’s formula gives a cross–check on the existence of large deviations; the algorithm, however, is much more cumbersome to implement on models other than the simple scalar field and even in this latter case it turns out to be less efficient.
## 7. Bootstrap analysis
The consideration of a totally trivial random variable, a high power of a simple normal deviate, is sufficient to convince ourselves that the occurrence of large deviation is actually very natural. Consider a standard Gaussian random variable $`x`$ and let $`𝒪=x^n`$. The probability density for $`𝒪`$ contains a power–law prefactor which dominates, for high $`n`$, over the exponential term. It is clear, by a saddle point argument, that for $`n5`$ the average $`x^n`$ is dominated by large deviations in the process; this appears to be a common crux in all models we have considered. Having established that the high order coefficients will suffer from large fluctuations, it is important to use reliable tools to estimate the statistical fluctuations. Since our estimators $`𝒪_n`$ will in general have strongly non-Gaussian distribution functions, which moreover are only known empirically through the numerical experiments, it appears a natural choice to apply some non–parametric analysis. We present here the results obtained by applying the bootstrap method . Let us first consider experiments on the 3-D Weingarten’s model, restricted to an elementary cube (a total of 12 link variables, 6 plaquettes). We subdivide the data coming from several runs at the same value $`\tau =.01`$ (measuring every 100 steps) in bins of $`N`$ samples (say $`N=1000`$) and on each bin we perform $`B`$ bootstrap replicas (in this example $`B=100`$). Fig.7 reports the result for $`𝒪_3`$. In this case the bootstrap analysis is indistinguishable from a standard gaussian analysis, which means that the distribution is not too far from normal. The standard gaussian analysis is performed by taking a weighted mean of of averages in each sample, i.e.
$`X_{gauss}`$ $`=`$ $`{\displaystyle \frac{_jx_j/\sigma _j^2}{_j1/\sigma _j^2}}`$
$`\sigma ^2`$ $`=`$ $`{\displaystyle \underset{j}{}}\sigma _j^2`$
A discrepancy between these values and the bootstrap’s signals a significant deviation from normality, like in Fig.8.
The next plot refers to the estimator $`𝒪_{11}`$. Here the deviation from normal is relevant: it is clearly reflected in a strong discrepancy between gaussian–like weighted mean and bootstrap estimate. The bootstrap apparently gives a reliable estimate for the confidence interval, as can be seen in the close up picture (Fig. 9).
Finally let us give a look at the overall pattern obtained in 3-D. Fig. 10 is remarkably similar to the one obtained for the plaquette expansion coefficients in Wilson’s 4-D $`SU(3)`$ Lattice Gauge Theory.
## 8. Lattice Gauge Theories
A good message from this paper is that the picture we just drew for simple models does not spoil the confidence we have in Numerical Stochastic Perturbation Theory for a field of real physical interest like Lattice Gauge Theory. As amazing as it can appear at first sight, this does not totally come as a surprise. In such a theory both the huge number of degrees of freedom and their coupling make it less likely applicable the simple argument of the previous section referring to a high power of a normal variable. One should also keep in mind that the gauge degrees of freedom make it possible to keep the norms of the perturbative components of the field under control by providing a restoring force which does not spoil the convergence of the process.
As hard as formal arguments can be and since an exact characterization of the stochastic processes is missing also for the simple models we considered above, we take a pragmatic attitude. The lesson we can learn from simple models is that monitoring frequencies histograms is a good tool in order to assess huge deviations from normality. By inspecting histograms from Lattice Gauge Theory simulations one can convince oneself that in this case convergence is much more reliable. This of course turns out to be consistent with the bootstrap and standard error analysis giving the same results. For instance, the plot in fig. 11 refers to a $`SU(3)`$ plaquette high order perturbative coefficient on a small lattice: the process is manifestly safe from “Pepe-effect”.
An independent hint for the reliability of Numerical Stochastic Perturbation Theory for Lattice Gauge Theory comes from . A couple of new perturbative orders in the expansion of the plaquette in 4-D $`SU(3)`$ (which means reaching $`\alpha ^{10}`$) are shown to be fully consistent with both the expected renormalon behaviour and with finite size effects on top of that. A more organic report on the status of the art concerning the application of the stochastic method to LGT will be given in .
## 9. Conclusions
We have performed a series of simulations to test the numerical stochastic perturbation algorithm previously introduced in the context of Lattice Gauge Theories. The emergence of large non–Gaussian fluctuations in the statistical estimators for high order coefficients was observed in all the simple models we considered with a similar pattern. The study of histograms, giving an estimate of the distribution functions for these estimators, exhibits a large window where exponential fall–off is masked by an approximately power–like behaviour; this fact entails that in any finite run there exist large deviations which are very rare but contribute to the average on the same foot with more frequent events. The lesson we draw from this experiment is twofold: it is necessary to monitor the histograms of the estimators in order to assess the deviation from normality; in case of large deviations, a reliable estimate of statistical error should be obtained by non–parametric methods, such as the bootstrap. These problems do not appear to plague the application of the method to Lattice Gauge Theories. The same viewpoint presented for simple models results in histograms hinting at good convergence properties, fully consistent with all our previous experience in this field.
## Acknowledgments
Many people contributed ideas which finally merged in this paper. We would like to thank A. Sokal, who before anybody else raised some doubts about the kind of convergence of the PW-process as we implement it and stimulated a deeper analysis, P. Butera and G. Marchesini for their continuous interest and support, M.Pepe, for the original observation which triggered our analysis, G.Burgio, for many interesting discussions. E.O. would like to warmly thank CERN’s Theory Division for hospitality while part of this work was performed.
## Appendix A Correlation functions at high Langevin time.
To simplify the argument, let us consider the model $`(i)`$ with a single degree of freedom (it may be generalized also to realistic systems). First of all we observe that any correlation function of free fields $`\phi _0(t)`$ may be exactly computed, in fact
(13)
$$\phi _0(t)^{2m}=(2m1)!!(1e^{2t})^m$$
converges exponentially to a limit. Then one can proceed by induction: it is sufficient to show that the correlation function Eq.(9), which has a total perturbative order $`p_T=_jp_j`$, may be written as a finite sum of correlation functions which have a total perturbative order strictly lower than $`p_T`$, plus correlation functions which have a total perturbative order equal to $`p_T`$ but with less free fields, plus exponentially damped terms.
To this end it is useful to re-write the formal solution given by Eq.(8) in discretized time ($`t=N\tau `$). The idea is to separate the solution into the memory terms, representing the last $`N1`$ steps, and the most recent contribution from lower orders and noise.
(14) $`\phi _0(t)`$ $`=`$ $`e^\tau \phi _0(t\tau )+\sqrt{\tau }\eta (t)`$
$`\phi _j(t)`$ $`=`$ $`e^\tau \phi _j(t\tau )+\tau f^{(j)}(t)`$
where
$$f^{(j)}(t)=\underset{i=0}{\overset{n1}{}}\underset{j=0}{\overset{i}{}}\phi _j(t)\phi _{ij}(t)\phi _{n1i}(t)$$
and $`\eta (t)`$ are independent gaussian variables with mean zero and variance $`2`$. We now should put Eq.(14) inside Eq.(9). Let’s do it for a correlation function of two fields (the extension to a general correlation function will be considered later). For any perturbative order $`i,j>0`$ we have:
$`\phi _i(N\tau )\phi _j(N\tau )`$ $`=`$ $`e^{2\tau }\phi _i(N\tau \tau )\phi _j(N\tau \tau )+`$
$`+\tau e^\tau \phi _i(N\tau \tau )f^{(j)}(N\tau )`$
$`+\tau e^\tau \phi _j(N\tau \tau )f^{(i)}(N\tau )`$
$`+\tau ^2f^{(i)}(N\tau )f^{(j)}(N\tau )`$
Let us use the general solution of the recursive formula:
$$y_l=\alpha _ly_{l1}+\beta _l$$
which is given by
$$y_N=\left(\underset{m=1}{\overset{N}{}}\alpha _m\right)y_0+\underset{n=1}{\overset{N}{}}\left(\underset{m=n+1}{\overset{N}{}}\alpha _m\right)\beta _n$$
In our case we simply have $`\alpha _m=e^{2\tau }`$, independent on $`m`$, while $`\beta _m`$ is composed of a linear ($`B_m`$) and a quadratic ($`C_m`$) part in $`\tau `$:
$`\beta _m`$ $`=`$ $`\tau e^\tau B_m(\tau )+\tau ^2C_m`$
$`B_m(\tau )`$ $`=`$ $`\phi _i(N\tau \tau )f^{(j)}(N\tau )+\phi _j(N\tau \tau )f^{(i)}(N\tau )`$
$`C_m`$ $`=`$ $`f^{(i)}(N\tau )f^{(j)}(N\tau ).`$
As a result:
$`\phi _i(N\tau )\phi _j(N\tau )=e^{2N\tau }\phi _i(0)\phi _j(0)+`$
$`+`$ $`e^{2N\tau }\left[{\displaystyle \underset{n=1}{\overset{N}{}}}e^{2n\tau }(\tau e^\tau B_n(\tau )+\tau ^2C_n)\right]`$
We observe that $`B_n(\tau )`$ is a correlation function of total perturbative order strictly lower than $`i+j`$. By inductive hypothesis we know that $`B_n(\tau )`$ has a finite limit for $`\tau 0`$ (at $`t=N\tau `$ fixed), let’s say
(17)
$$B_N(\tau )=B_t+O(\tau ),$$
and that $`B_t`$ has a finite limit for $`t\mathrm{}`$. Let us parameterize the remaining dependence on $`t`$ as:
(18)
$$B_t=B_{\mathrm{}}+a_1e^{2t}t^p+\underset{q_j>2}{}a_je^{q_jt}$$
From Eq.(13) it follows that the dependence on $`t^p`$ is absent in correlations of free fields, but such correction factors can appear at higher orders. By inserting Eqs.(17,18) in Eq.(A) the geometric series can be re-summed:
$`\phi _i(N\tau )\phi _j(N\tau )=e^{2N\tau }\phi _i(0)\phi _j(0)+`$
$`+`$ $`e^{2N\tau }\left[\left(\tau B_t+O(\tau ^2)\right)\left({\displaystyle \frac{1e^{2(N+1)\tau }}{1e^{2\tau }}}\right)+O(\tau )\right]`$
If we take first the limit $`\tau 0`$ at $`t=N\tau `$ fixed and then $`t\mathrm{}`$, we find
$$\underset{t\mathrm{}}{lim}\phi _i(t)\phi _j(t)=\frac{1}{2}B_{\mathrm{}}.$$
Up to now, we have not considered the case of correlation functions with some free fields together with higher order fields. In this case the argument runs almost in the same way, but $`\tau f^{(j)}`$ is substituted by $`\sqrt{\tau }\eta `$:
$`\phi _0(N\tau )\phi _0(N\tau )(\mathrm{})`$ $`=`$ $`e^{2\tau }\phi _0(N\tau \tau )\phi _0(N\tau \tau )(\mathrm{})+`$
$`+2\sqrt{\tau }e^\tau \phi _0(N\tau \tau )\eta (N\tau )(\mathrm{})`$
$`+\tau \eta (N\tau )\eta (N\tau )(\mathrm{})`$
Only terms with an even number of $`\eta `$ at time $`N\tau `$ contributes. In this case the inductive hypothesis must be applied to correlation functions with the same total perturbative order but with a lower number of free fields.
The same argument can be applied to a general correlation function, obtaining in the limit $`t\mathrm{}`$
$`{\displaystyle \underset{l=0}{\overset{p}{}}}\phi _l^{n_l}`$ $`=`$ $`{\displaystyle \frac{1}{_{l=0}^pn_l}}[n_0(n_01){\displaystyle \underset{l=0}{\overset{p}{}}}\phi _l^{n_l}\widehat{(\phi _0)^2}`$
$`{\displaystyle \underset{m=1}{\overset{p}{}}}n_m{\displaystyle \underset{i=0}{\overset{m1}{}}}{\displaystyle \underset{j=0}{\overset{i}{}}}{\displaystyle \underset{l=0}{\overset{p}{}}}\phi _l^{n_l}\widehat{\phi _m}\phi _j\phi _{ij}\phi _{m1i}]`$
where $`p`$ is the maximum perturbative order present, and $`\widehat{\phi }`$ means that a factor $`\phi `$ should be dropped in the expression.
Induction allows us to conclude that all correlation functions reach a finite limit.
If one goes through the previous argument by keeping the first correction in Eq.(18) it is possible to show that convergence to equilibrium is dominated by an exponential factor which is at least $`e^{2t}`$. However, even if in correlations of free fields the dependence on $`t`$ has exactly an exponential form (or sum of exponentials), in general (at higher orders) one must expect power corrections like $`t^\alpha e^{2t}`$. The precise determination of them is tedious and beyond our scope.
The formula given in Eq.(A) is useful also because it allows to compute not only the mean values of observables but also some high moments of their distributions, at least for the model ($`i`$). In this way we have been able to show that the large fluctuations were neither an artifact of the numeric simulation nor the effect of a slow convergence to equilibrium.
## Appendix B The bootstrap algorithm
The idea of bootstrap is very simple. Let $`X`$ be a random variable and $`x_1,x_2,\mathrm{},x_N`$ an $`N`$tuple of values generated by a physical process, by the stock market, by your Monte Carlo, whichever the case you are currently studying. In the absence of any a priori information on the distribution function for $`X`$, one makes the most conservative assumption, namely one adopts as distribution function the empirical distribution with density
$$\rho _{boot}^N(x)=\frac{1}{N}\underset{j=1}{\overset{N}{}}\delta (xx_j).$$
One then uses $`\rho _{boot}`$ to generate other $`N`$-tuples; all values $`x_j`$, $`(j=1,\mathrm{},N)`$ being equally probable. Having produced $`B`$ such $`N`$-tuples one can estimate various statistics of interest, like mean, median, quartiles etc. For instance the definition of standard deviation is simply
(20)
$$\mathrm{\Delta }x=\sqrt{\frac{1}{B1}\underset{k=1}{\overset{B}{}}\left(x_kx\right)^2}$$
where $`x_k`$ is the average computed on the $`k`$-th $`N`$-tuple and $`x`$ is their mean value. Recommended values for $`B`$ are in the range $`25B200`$. It is obvious that no improvement can be obtained on mean values; the virtue of the method should be found in the easy estimation of standard errors which do not rely on the assumption of normality. This does not mean of course that one should not take care of autocorrelation problems or other sources of statistical inaccuracies.
The method is implemented very easily in any language. A one–liner to produce a new $`N`$-tuple in matlab is the following (if the vector $`X`$ contains the original data)
> XB = X(ceil(N * rand(N,1)));
We refer to Refs. for a thorough discussion of the method. |
warning/0002/hep-ph0002253.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The search for supersymmetric (SUSY) particles and the determination of their properties will be one of the main goals of a future $`e^+e^{}`$ linear collider. Particularly interesting will be the experimental study of the neutralinos, which are the quantum mechanical mixtures of the neutral gauginos and higgsinos, the SUSY partners of the neutral gauge and Higgs bosons. In the Minimal Supersymmetric Standard Model (MSSM) there are four neutralinos $`\stackrel{~}{\chi }_i^0`$, $`i=1,\mathrm{},4`$. In extensions of the MSSM there may be more than four neutralinos. Usually the lightest neutralino $`\stackrel{~}{\chi }_1^0`$ is the lightest SUSY particle LSP. Therefore, the production of the lightest and the second lightest neutralino can presumably be studied at an $`e^+e^{}`$ linear collider with a CMS energy $`\sqrt{s}=500`$ GeV. The aim will be to precisely determine the SUSY parameters of the neutralino system. By a detailed study of the neutralino system one can also examine the question whether the MSSM or another SUSY model is realized in nature. Models with an extended neutralino sector have been discussed in . Recently a method for determining the SUSY parameter $`M_2`$, $`\mu `$ and $`\mathrm{tan}\beta `$ by measuring suitable observables in chargino production $`e^+e^{}\stackrel{~}{\chi }_i^+\stackrel{~}{\chi }_j^{}`$ ($`i,j=1,2`$) has been proposed in . The process $`e^+e^{}\stackrel{~}{\chi }_i^+\stackrel{~}{\chi }_j^{}`$, $`\stackrel{~}{\chi }_i^+\stackrel{~}{\chi }_k^0\mathrm{}^+\nu `$ including the full spin correlation has been studied in .
In previous papers (see and references therein) the cross sections of $`e^+e^{}\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$, $`i,j=1\mathrm{}4`$, and the branching ratios and energy distributions of neutralino decays were studied, which do not depend on spin correlations. In the calculation of the decay angular distributions, however, one has to take into account the spin correlations between production and decay of the neutralinos. In we have studied the process $`e^+e^{}\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$, $`i,j=1\mathrm{}4`$, with unpolarized beams and the subsequent leptonic decays $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_k^0\mathrm{}^+\mathrm{}^{}`$ including the complete spin correlations. In we have given the complete analytical formulae for longitudinally polarized beams, fully including the spin correlations between production and decay. The formulae have been given in the laboratory system in terms of the basic kinematic variables.
In the present paper we extend our analysis in and study neutralino production in the case that both beams are polarized. We show that by a suitable choice of the $`e^+`$ beam polarization not only higher cross sections but also additional information can be gained. We give numerical predictions for the process $`e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$ with $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0e^+e^{}`$ including the full spin correlations between production and decay. The framework of our studies is the MSSM. We first assume the GUT relation $`M_1/M_2=\frac{5}{3}\mathrm{tan}^2\mathrm{\Theta }_W`$ for the gaugino mass parameters (note that in Refs. we used the notation $`M^{}`$ and $`M`$ for $`M_1`$ and $`M_2`$). We consider two gaugino–like and one higgsino–like scenario and study the dependence of the cross section, the forward–backward asymmetry of the decay electron, and the opening angle distribution of the decay lepton pair $`e^+`$ and $`e^{}`$ on the beam polarizations and on the masses of the exchanged $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$. Then we relax the GUT relation between $`M_1`$ and $`M_2`$ and study the $`M_1`$ dependence of $`\sigma (e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)\times BR(\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0e^+e^{})`$ and the forward-backward asymmetry of the decay electron for various beam polarizations and slepton masses. We also discuss which observables are suitable for the determination of the neutralino parameters and selectron masses.
## 2 Spin correlations between production and decay
Both the helicity amplitudes $`T_P^{\lambda _i\lambda _j}`$ for the production process $`e^{}(p_1)e^+(p_2)\stackrel{~}{\chi }_i^0(p_3,\lambda _i)\stackrel{~}{\chi }_j^0(p_4,\lambda _j)`$ and the helicity amplitudes $`T_{D,\lambda _i}`$ and $`T_{D,\lambda _j}`$ for the decay processes $`\stackrel{~}{\chi }_i^0(p_3,\lambda _i)\stackrel{~}{\chi }_k^0(p_5)\mathrm{}^+(p_6)\mathrm{}^{}(p_7)`$ and $`\stackrel{~}{\chi }_j^0(p_4,\lambda _j)\stackrel{~}{\chi }_l^0(p_8)\mathrm{}^+(p_9)\mathrm{}^{}(p_{10})`$, respectively, receive contributions from $`Z^0`$ exchange in the direct channel and from $`\stackrel{~}{e}_{L,R}`$ exchange in the crossed channels.
The amplitude squared of the combined process of production and decay is:
$$|T|^2=|\mathrm{\Delta }(\stackrel{~}{\chi }_i^0)|^2|\mathrm{\Delta }(\stackrel{~}{\chi }_j^0)|^2\rho _P^{\lambda _i\lambda _j\lambda _i^{^{}}\lambda _j^{^{}}}\rho _{D,\lambda _i^{^{}}\lambda _i}\rho _{D,\lambda _j^{^{}}\lambda _j}\text{(summed over helicities)}.$$
(1)
It is composed of the (unnormalized) spin density production matrix
$$\rho _P^{\lambda _i\lambda _j\lambda _i^{^{}}\lambda _j^{^{}}}=T_P^{\lambda _i\lambda _j}T_P^{\lambda _i^{^{}}\lambda _j^{^{}}},$$
(2)
the decay matrices
$$\rho _{D,\lambda _i^{^{}}\lambda _i}=T_{D,\lambda _i}T_{D,\lambda _i^{^{}}}^{}\text{and}\rho _{D,\lambda _j^{^{}}\lambda _j}=T_{D,\lambda _j}T_{D,\lambda _j^{^{}}}^{},$$
(3)
and the propagators
$$\mathrm{\Delta }(\stackrel{~}{\chi }_{i,j}^0)=1/[p_{3,4}^2m_{i,j}^2+im_{i,j}\mathrm{\Gamma }_{i,j,}].$$
(4)
Here $`p_{3,4}^2`$, $`\lambda _{i,j}`$, $`m_{i,j}`$ and $`\mathrm{\Gamma }_{i,j}`$ denote the four–momentum squared, helicity, mass and total width of $`\stackrel{~}{\chi }_{i,j}^0`$. For these propagators we use the narrow–width approximation.
Introducing a suitable set of polarization vectors for each of the neutralinos the density matrices can be expanded in terms of Pauli matrices $`\sigma ^a`$:
$`\rho _P^{\lambda _i\lambda _j\lambda _i^{^{}}\lambda _j^{^{}}}`$ $`=`$ $`(\delta _{\lambda _i\lambda _i^{^{}}}\delta _{\lambda _j\lambda _j^{^{}}}P(\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0)+\delta _{\lambda _j\lambda _j^{^{}}}{\displaystyle \underset{a}{}}\sigma _{\lambda _i\lambda _i^{^{}}}^a\mathrm{\Sigma }_P^a(\stackrel{~}{\chi }_i^0)`$ (5)
$`+\delta _{\lambda _i\lambda _i^{^{}}}{\displaystyle \underset{b}{}}\sigma _{\lambda _j\lambda _j^{^{}}}^b\mathrm{\Sigma }_P^b(\stackrel{~}{\chi }_j^0)+{\displaystyle \underset{ab}{}}\sigma _{\lambda _i\lambda _i^{^{}}}^a\sigma _{\lambda _j\lambda _j^{^{}}}^b\mathrm{\Sigma }_P^{ab}(\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0)),`$
$`\rho _{D,\lambda _i^{^{}}\lambda _i}`$ $`=`$ $`(\delta _{\lambda _i^{^{}}\lambda _i}D(\stackrel{~}{\chi }_i^0)+{\displaystyle \underset{a}{}}\sigma _{\lambda _i^{^{}}\lambda _i}^a\mathrm{\Sigma }_D^a(\stackrel{~}{\chi }_i^0)),`$ (6)
$`\rho _{D,\lambda _j^{^{}}\lambda _j}`$ $`=`$ $`(\delta _{\lambda _j^{^{}}\lambda _j}D(\stackrel{~}{\chi }_j^0)+{\displaystyle \underset{b}{}}\sigma _{\lambda _j^{^{}}\lambda _j}^b\mathrm{\Sigma }_D^b(\stackrel{~}{\chi }_j^0))\text{, a, b=1,2,3}.`$ (7)
We choose the polarization vectors such that $`\mathrm{\Sigma }_P^1(\stackrel{~}{\chi }_{i,j}^0)`$ describes the transverse polarization in the production plane, $`\mathrm{\Sigma }_P^2(\stackrel{~}{\chi }_{i,j}^0)`$ denotes the polarization perpendicular to the production plane and $`\mathrm{\Sigma }_P^3(\stackrel{~}{\chi }_{i,j}^0)`$ describes the longitudinal polarization of the respective neutralino. $`\mathrm{\Sigma }_P^{ab}(\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0)`$ is due to correlations between the polarizations of both neutralinos. The complete analytical expressions for the production density matrix and for the decay matrices are given in .
If CP is conserved the neutralino couplings are real. It can be seen in that in this case, due to the Majorana character of the neutralinos, the quantities $`P`$, $`\mathrm{\Sigma }_P^1`$, $`\mathrm{\Sigma }_P^{11}`$, $`\mathrm{\Sigma }_P^{22}`$, $`\mathrm{\Sigma }_P^{33}`$, $`\mathrm{\Sigma }_P^{23}`$ are forward–backward symmetric, whereas the quantities $`\mathrm{\Sigma }_P^2`$, $`\mathrm{\Sigma }_P^3`$, $`\mathrm{\Sigma }_P^{12}`$, $`\mathrm{\Sigma }_P^{13}`$ are forward–backward antisymmetric.
The amplitude squared of the combined process of production and decay can be written as:
$`|T|^2=4|\mathrm{\Delta }(\stackrel{~}{\chi }_i^0)|^2|\mathrm{\Delta }(\stackrel{~}{\chi }_j^0)|^2(P(\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0)D(\stackrel{~}{\chi }_i^0)D(\stackrel{~}{\chi }_j^0)+{\displaystyle \underset{a=1}{\overset{3}{}}}\mathrm{\Sigma }_P^a(\stackrel{~}{\chi }_i^0)\mathrm{\Sigma }_D^a(\stackrel{~}{\chi }_i^0)D(\stackrel{~}{\chi }_j^0)`$
$`+{\displaystyle \underset{b=1}{\overset{3}{}}}\mathrm{\Sigma }_P^b(\stackrel{~}{\chi }_j^0)\mathrm{\Sigma }_D^b(\stackrel{~}{\chi }_j^0)D(\stackrel{~}{\chi }_i^0)+{\displaystyle \underset{a,b=1}{\overset{3}{}}}\mathrm{\Sigma }_P^{ab}(\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0)\mathrm{\Sigma }_D^a(\stackrel{~}{\chi }_i^0)\mathrm{\Sigma }_D^b(\stackrel{~}{\chi }_j^0)).`$ (8)
The differential cross section in the laboratory system is then given by
$$d\sigma _e=\frac{1}{2s}|T|^2(2\pi )^4\delta ^4(p_1+p_2\underset{i}{}p_i)d\mathrm{lips}(p_3\mathrm{}p_{10}),$$
(9)
$`d\mathrm{lips}(p_3,\mathrm{},p_{10})`$ is the Lorentz invariant phase space element.
If one neglects all spin correlations between production and decay only the first term in (8) contributes. The second and third term in (8) describe the spin correlations between the production and the decay process and the last term is due to spin–spin correlations between both decaying neutralinos.
## 3 Numerical analysis and discussion
In the MSSM the masses and couplings of neutralinos are determined by the parameters $`M_1`$, $`M_2`$, $`\mu `$, $`\mathrm{tan}\beta `$, which can be chosen real if CP violation is neglected. Moreover, one usually makes use of the GUT relation
$$M_1=\frac{5}{3}M_2\mathrm{tan}^2\mathrm{\Theta }_W.$$
(10)
The neutralino mass mixing matrix in the convention used can be found in .
The total cross section of $`e^+e^{}\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$ and the decay rate $`\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_k^0e^+e^{}`$ further depend on the masses of $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$. In the following numerical analysis we study neutralino production and decay in three scenarios, which we denote by A1, A2, and B. The corresponding parameters are given in Table 1.
Scenario A1 corresponds to that in . In this scenario $`\stackrel{~}{\chi }_1^0`$ is $`\stackrel{~}{B}`$–like and $`\stackrel{~}{\chi }_2^0`$ is $`\stackrel{~}{W}^3`$–like. The selectron masses are $`m_{\stackrel{~}{e}_L}=176`$ GeV, $`m_{\stackrel{~}{e}_R}=132`$ GeV. Scenario A2 differs from A1 only by the mass of $`\stackrel{~}{e}_L`$. In scenario B the same masses as in scenario A1 are taken for $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, $`\stackrel{~}{e}_L`$, $`\stackrel{~}{e}_R`$. $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ are, however, higgsino–like due to the choice $`\mu <M_2`$.
### 3.1 Effects of beam polarizations on the total cross section
The cross section $`\sigma (e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)`$ is shown in Figs. 1a, 1b, and 1c as a function of the longitudinal beam polarizations $`P_{}^3`$ for electrons and $`P_+^3`$ for positrons, for scenario A1, A2 and B, respectively, (with $`P_\pm ^3=\{1,0,1\}`$ for $`\{`$left–, un–, right–$`\}`$polarized ). The cross section $`\sigma (e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)`$ is shown at $`\sqrt{s}=(m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0})+30`$ GeV.
In Figs. 1a, 1b, and 1c the white area is covered by an electron polarization $`|P_{}^3|85\%`$ and a positron polarization $`|P_+^3|60\%`$. The cross section can be enhanced by a factor 2–3 by polarizing both beams. Theoretically, for pure gaugino–like neutralinos and $`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}`$ ($`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}`$) and $`P_{}^3=+1`$, $`P_+^3=1`$ ($`P_{}^3=1`$, $`P_+^3=+1`$), the cross section could be enlarged by a factor 4. For pure higgsino–like neutralinos and $`P_{}^3=+1`$, $`P_+^3=1`$ ($`P_{}^3=1`$, $`P_+^3=+1`$) this factor would be 1.7 (2.3) .
One clearly recognizes, Figs. 1a, 1b, and 1c, the sensitive dependence of the cross section on the selectron masses $`m_{\stackrel{~}{e}_L}`$ and $`m_{\stackrel{~}{e}_R}`$ as well as on the mixing character of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$. In scenario A2 with $`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}`$ and gaugino–like $`\stackrel{~}{\chi }_{1,2}^0`$, one expects the largest cross section for $`P_{}^3=+1`$ and $`P_+^3=1`$, see Fig. 1b. In scenario B the cross section is governed by $`Z^0`$ exchange and is therefore rather symmetric for $`P_{}^3=\pm P_+^3=`$. The beam polarizations are a useful tool for getting more information about $`m_{\stackrel{~}{e}_L}`$ and $`m_{\stackrel{~}{e}_R}`$.
If the polarizations of both beams are varied, the relative size of the cross sections strongly depends on the mixing character of both neutralinos $`\stackrel{~}{\chi }_{1,2}^0`$ and on the selectron masses . If $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ are pure higgsinos, one obtains for $`|P_{}^3|=85\%`$ and $`|P_+^3|=60\%`$ the sequence
$$\sigma _e^+>\sigma _e^+>\sigma _e^0>\sigma _e^{00}>\sigma _e^{+0}>\sigma _e^{}>\sigma _e^{++}.$$
(11)
Here $`\sigma _e=\sigma (e^{}e^+\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)\times BR(\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0\mathrm{}^+\mathrm{}^{})`$ and $`(+)`$ etc. denotes the sign of the electron polarization $`P_{}^3`$ and of the positron polarization $`P_+^3`$, respectively.
If $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$ are pure gauginos the order of the cross sections depends on the relative magnitude of the selectron masses $`m_{\stackrel{~}{e}_L}`$ and $`m_{\stackrel{~}{e}_R}`$. For $`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}`$ only right selectron exchange contributes, and one obtains
$$\sigma _e^+>\sigma _e^{+0}>\sigma _e^{00}>\sigma _e^{++}>\sigma _e^{}>\sigma _e^0>\sigma _e^+,$$
(12)
whereas for $`m_{\stackrel{~}{e}_R}m_{\stackrel{~}{e}_L}`$, one gets:
$$\sigma _e^+>\sigma _e^0>\sigma _e^{00}>\sigma _e^{}>\sigma _e^{++}>\sigma _e^{+0}>\sigma _e^+.$$
(13)
The case of a heavy right slepton may be realized in extended SUSY models ( and references therein).
Comparing (11) and (13) shows that polarizing both beams allows one to distinguish between a higgsino–like scenario and a gaugino–like scenario with dominating $`\stackrel{~}{e}_L`$ exchange. This is not possible if only the electron beam is polarized.
In Table 2 we show the cross sections for various polarization configurations for our scenarios A1, A2, and B at $`\sqrt{s}=(m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0})+30`$ GeV for $`P_{}^3=0,\pm 85\%`$ and $`P_+^3=0,\pm 60\%`$. For the higgsino–like scenario B the sequence of the cross sections coincides with (11). One notices that one obtains the same ordering of the polarized cross sections for the gaugino–like scenario A1. This shows that the relative size of the cross sections sensitively depends on the mass difference between $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$, which is rather small in scenario A1, $`m_{\stackrel{~}{e}_L}m_{\stackrel{~}{e}_R}=44`$ GeV. Comparing the sequence of cross sections for the gaugino–like scenario A2, see Table 2, with that for pure $`\stackrel{~}{e}_R`$ exchange, see (12), one sees a small influence of $`\stackrel{~}{e}_L`$ exchange despite the rather high $`\stackrel{~}{e}_L`$ mass, $`m_{\stackrel{~}{e}_L}=500`$ GeV.
### 3.2 Lepton forward–backward asymmetry
Owing to the Majorana character of the neutralinos the angular distribution of the production process is forward–backward symmetric . The angular distribution of the decay lepton, however, depends sensitively on the polarization of $`\stackrel{~}{\chi }_i^0`$. Since the longitudinal polarization $`\mathrm{\Sigma }_P^3`$ and the transverse polarization $`\mathrm{\Sigma }_P^1`$ of $`\stackrel{~}{\chi }_i^0`$ are forward–backward antisymmetric, the lepton forward–backward asymmetry $`A_{FB}`$ of the decay lepton may become quite large. The lepton forward–backward asymmetry $`A_{FB}`$ is defined as
$$A_{FB}=\frac{\sigma _e(\mathrm{cos}\mathrm{\Theta }_e>0)\sigma _e(\mathrm{cos}\mathrm{\Theta }_e<0)}{\sigma _e(\mathrm{cos}\mathrm{\Theta }_e>0)+\sigma _e(\mathrm{cos}\mathrm{\Theta }_e<0)},$$
(14)
where $`\sigma _e`$ is a short–hand notation for $`\sigma _e=\sigma (e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)\times BR(\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0e^+e^{})`$. We will show $`A_{FB}`$ not too far from threshold because it decreases with $`\sqrt{s}`$ for fixed neutralino masses.
In Figs. 2a, 2b, and 2c we show $`A_{FB}`$ of the decay electron as a function of the electron and positron polarizations for the scenarios A1, A2, and B, respectively, at $`\sqrt{s}=(m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0})+30`$ GeV. First one notices that polarizing suitably both beams gives a larger asymmetry. Actually, in the scenarios A1 and B $`A_{FB}`$ turns out to be practically zero for both beams unpolarized. The figures also exhibit a very different pattern. When comparing Fig. 2a with Fig. 2b, the different behaviour is due to the different masses of $`\stackrel{~}{e}_L`$. In the higgsino scenario B the asymmetries are much smaller. Notice again the symmetry of $`P_{}^3=\pm P_+^3=`$ in scenario B as already observed in the total cross section. Measuring the lepton forward–backward asymmetry $`A_{FB}`$ in addition to the total cross section strongly constrains the selectron masses $`m_{\stackrel{~}{e}_L}`$ and $`m_{\stackrel{~}{e}_R}`$ as well as the mixing properties of $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_2^0`$.
We also studied the dependence of the lepton forward–backward asymmetry $`A_{FB}`$ on $`\sqrt{s}`$. For $`\sqrt{s}(m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0})`$ the angular distribution of the decay lepton is essentially the same as that of the decaying neutralino $`\stackrel{~}{\chi }_2^0`$ . Therefore the lepton forward–backward asymmetry practically vanishes.
### 3.3 Opening angle distribution
The opening angle distribution between the two leptons from the decay of one of the neutralinos, $`\stackrel{~}{\chi }_2^0\mathrm{}^+\mathrm{}^{}\stackrel{~}{\chi }_1^0`$, is independent of the spin correlations due to the Majorana nature of the neutralinos . Therefore, it factorizes into the contributions from production and decay. For the same reason this is also valid for the energy distribution of the neutralino decay products. For both distributions it is suitable to parametrize the phase space by the scattering angle $`\mathrm{\Theta }`$ between the incoming $`e^{}(p_1)`$ beam and the outgoing neutralino $`\stackrel{~}{\chi }_2^0(p_4)`$, the azimuthal angle $`\mathrm{\Phi }_{\stackrel{~}{\chi }_2^0\mathrm{}^{}}`$ between the scattering plane and the ($`\stackrel{~}{\chi }_2^0\mathrm{}^{}`$)–plane and the opening angle $`\mathrm{\Theta }_+`$ between the leptons $`\mathrm{}^+`$ and $`\mathrm{}^{}`$ from the decay of the neutralino $`\stackrel{~}{\chi }_2^0`$. Since the phase space is independent of the azimuthal angle, the contributions of the transverse polarizations of the neutralino vanish after integration over $`\mathrm{\Phi }_{\stackrel{~}{\chi }_2^0\mathrm{}^{}}`$. As for the longitudinal polarization $`\mathrm{\Sigma }_P^3`$, the Majorana character of the neutralino is crucial. If CP is conserved $`\mathrm{\Sigma }_P^3`$ is forward–backward antisymmetric, $`\mathrm{\Sigma }_P^3(\mathrm{cos}\mathrm{\Theta })=\mathrm{\Sigma }_P^3(\mathrm{cos}\mathrm{\Theta })`$, so that the contribution of the longitudinal polarization vanishes after integration over the scattering angle $`\mathrm{\Theta }`$ .
In Figs. 3a and 3b we show the distribution of the angle $`\mathrm{\Theta }_+`$ between the decay leptons from $`e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$, $`\stackrel{~}{\chi }_2^0e^+e^{}\stackrel{~}{\chi }_1^0`$ in the scenarios A1, A2 and B at $`\sqrt{s}=(m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0})+30`$ GeV and for various beam polarizations. One notices that the shape of the $`\mathrm{\Theta }_+`$ distribution is mainly determined by the mixing character of the neutralinos. The selectron masses mainly influence the size of the cross section. Due to the factorization of production and decay the beam polarizations have no influence on the shape of this distribution.
### 3.4 Dependence on $`M_1`$
So far we have used the GUT relation (10) for the gaugino masses. In the following we will be more general and not use this relation . We will discuss the $`M_1`$ dependence of the cross section and the forward–backward asymmetry of the decay electron . All other parameters are chosen as in scenario A1 except the mass of $`\stackrel{~}{e}_R`$, $`m_{\stackrel{~}{e}_R}=161`$ GeV. The neutralino masses as well as the $`Z^0\stackrel{~}{\chi }_i^0\stackrel{~}{\chi }_j^0`$ couplings $`O_{ij}^{{}_{}{}^{\prime \prime }L}`$ and the $`\stackrel{~}{\chi }_i^0\stackrel{~}{\mathrm{}}\mathrm{}`$ couplings $`f_\mathrm{}i^{L,R}`$ depend on $`M_1`$.
In Fig. 4 we show the neutralino masses as function of $`M_1`$. The grey areas are excluded by the constraints $`m_{\stackrel{~}{\chi }_1^0}<m_{\stackrel{~}{\chi }_1^\pm }`$, $`m_{\stackrel{~}{\chi }_1^0}>35`$ GeV. We see that in the interval between $``$130 GeV$`<M_1<M_2`$ $`m_{\stackrel{~}{\chi }_1^0}`$ depends very strongly on $`M_1`$, whereas all other neutralino masses are nearly independent of $`M_1`$. On the other hand, in the region $`M_2M_1|\mu |`$ only $`m_{\stackrel{~}{\chi }_2^0}`$ depends strongly on $`M_1`$.
In the formulae for the cross section $`\sigma (e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)`$ and for the decay $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0\mathrm{}^+\mathrm{}^{}`$ the products of couplings $`f_\mathrm{}1^Lf_\mathrm{}2^L`$ and $`f_\mathrm{}1^Rf_\mathrm{}2^R`$ enter. We therefore show for these products the dependence on $`M_1`$ in Fig. 5. We do not consider values $`|M_1|>160`$ GeV, where $`m_{\stackrel{~}{\chi }_2^0}>m_{\stackrel{~}{e}_R}`$, because then the two–body decay $`\stackrel{~}{\chi }_2^0\stackrel{~}{e}_R+e`$ would be possible. One observes a strong variation of $`f_\mathrm{}1^Rf_\mathrm{}2^R`$ and $`f_\mathrm{}1^Lf_\mathrm{}2^L`$ for $`M_1>80`$ GeV. In particular, $`f_\mathrm{}1^Rf_\mathrm{}2^R`$ has a positive maximum at 140 GeV, whereas $`f_\mathrm{}1^Lf_\mathrm{}2^L`$ is zero at $`M_1=120`$ GeV and reaches large negative values for $`M_1160`$ GeV. We therefore have the following regions: $`f_\mathrm{}1^Lf_\mathrm{}2^L>f_\mathrm{}1^Rf_\mathrm{}2^R>0`$ for $`200`$ GeV$`<M_1<80`$ GeV, $`|f_\mathrm{}1^Rf_\mathrm{}2^R|>f_\mathrm{}1^Lf_\mathrm{}2^L`$ for 110 GeV$`<M_1<140`$ GeV, and $`|f_\mathrm{}1^Lf_\mathrm{}2^L|>f_\mathrm{}1^Rf_\mathrm{}2^R`$ for $`M_1>150`$ GeV. The coupling $`O_{12}^{{}_{}{}^{\prime \prime }L}`$ is small in this gaugino–like scenario.
Fig. 6a exhibits the $`M_1`$ dependence of $`\sigma (e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0)\times BR(\stackrel{~}{\chi }_2^0e^+e^{}\stackrel{~}{\chi }_1^0)`$ at $`\sqrt{s}=(m_{\stackrel{~}{\chi }_1^0}+m_{\stackrel{~}{\chi }_2^0})+30`$ GeV in the region 40 GeV$`<M_1<`$160 GeV for various beam polarizations. Since the masses of $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$ are in this case comparable, the curves reflect the behaviour of $`f_\mathrm{}1^Lf_\mathrm{}2^L`$ and $`f_\mathrm{}1^Rf_\mathrm{}2^R`$ of Fig. 5. A left (right) beam polarization of the electron (positron) selects the $`\stackrel{~}{e}_L`$ exchange, while the maximum in the curve with a right (left) electron (positron) polarization is due the maximum of $`f_\mathrm{}1^Rf_\mathrm{}2^R`$ in Fig. 5.
Fig. 6b shows the analogous curves for a heavy $`\stackrel{~}{e}_L`$ ($`m_{\stackrel{~}{e}_L}=500`$ GeV) and all other parameters as in Fig. 6a. One clearly sees that the $`\stackrel{~}{e}_L`$ exchange is strongly suppressed, and one obtains higher cross sections for right polarized $`e^{}`$ beams.
We have also studied the $`M_1`$ dependence of the forward–backward asymmetry $`A_{FB}`$ of the decay electron, eq. (14). It is shown in Fig. 7a for $`m_{\stackrel{~}{e}_L}=176`$ GeV, and in Fig. 7b for $`m_{\stackrel{~}{e}_L}=500`$ GeV. One notices a strong variation with $`M_1`$ and a strong dependence on the beam polarizations. Comparing Fig. 7a and Fig. 7b, one observes a very pronounced difference of the forward–backward asymmetry of the decay electron in the region 40 GeV$`<M_1<`$100 GeV. This is due to the suppression of $`\stackrel{~}{e}_L`$ exchange in Fig. 7b. The beam polarizations enhance the effect considerably. The peak at $`M_1120`$ GeV is again due to the maximum of $`f_\mathrm{}1^Rf_\mathrm{}2^R`$.
## 4 Conclusions
The objective of this paper has been twofold. Firstly, we have studied the advantage of having both the $`e^{}`$ and the $`e^+`$ beam polarized. If the polarizations of $`e^{}`$ and $`e^+`$ are varied, the relative size of the cross sections depends significantly on the mixing character of the neutralinos and on the masses of $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$. By an appropriate choice of the polarizations one can obtain up to three times larger cross sections than in the unpolarized case. Secondly, by taking into account the full spin correlations between production and decay, we have studied the angular distribution, as well as the forward–backward asymmetry of the decay electron $`e^+e^{}\stackrel{~}{\chi }_1^0\stackrel{~}{\chi }_2^0`$, $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0e^+e^{}`$. Measuring this asymmetry for various beam polarizations strongly constrains the masses of $`\stackrel{~}{e}_L`$ and $`\stackrel{~}{e}_R`$ and the mixing properties of the neutralinos. We have also studied the dependence on the gaugino mass parameter $`M_1`$. For a determination of $`M_1`$ the use of polarized $`e^+`$ and $`e^{}`$ beams would be very helpful. Due to the Majorana character of the neutralinos the opening angle distribution between the decay leptons is independent of spin correlations. It is very sensitive to the mixing character of the neutralinos, whereas its shape is only weakly dependent on the selectron masses.
## Acknowledgments
We are grateful to W. Porod and S. Hesselbach for providing the computer programs for neutralino widths. G.M.-P. was supported by Friedrich-Ebert-Stiftung. This work was also supported by the German Federal Ministry for Research and Technology (BMBF) under contract number 05 7WZ91P (0), by the Deutsche Forschungsgemeinschaft under contract Fr 1064/4-1, and the ‘Fonds zur Förderung der wissenschaftlichen Forschung’ of Austria, Project No. P13139-PHY. |
warning/0002/math0002252.html | ar5iv | text | # Birational properties of pencils of del Pezzo surfaces of degree 1 and 2.
## 1. Notions and results.
Due to progress in the Minimal Model Program, the classification problem in the modern birational geometry for varieties of negative Kodaira dimension can be formulated as follows: given a class of birational equivalency, to describe all Mori fibrations and birational maps between them.
Recall that a normal variety $`V`$ with at most $``$-factorial terminal singularities is called a Mori fibration if there exists an extremal contraction $`\phi :VS`$ of fibering type, i.e.
* $`\phi `$ is a morphism with connected fibers onto a normal variety $`S`$ and $`dimS<dimV`$;
* $`K_V`$ is $`\phi `$-ample and the relative Picard number is equal to 1: $`\rho (V/S)=\rho (V)\rho (S)=1`$.
Note that often comparing birational classes of varieties, we suffice to know Mori structures rather than Mori fibrations themselves (roughly speaking, Mori fibrations modulo birational maps over the base). From this viewpoint, the following class of varieties is rather important (and simple for describing):
###### Definition 1.1.
A Mori fibration $`V/S`$ is said to be birationally rigid, if any birational map $`\chi :VV^{}`$ onto another Mori fibration $`V^{}/S^{}`$ is birational over the base (”square”), i.e., there exists a birational map $`\psi :SS^{}`$ making the following diagram to be commutative:
$$\begin{array}{ccc}V& \stackrel{\chi }{}& V^{}\\ & & \\ S& \stackrel{\psi }{}& S^{}\end{array}$$
Now let $`\rho :VS`$ be a Mori fibration. Then any non-empty linear system $`𝒟`$ on $`V`$ is a subsystem of $`|a(K_V)+\rho ^{}(A)|`$, where $`a=\mu (V,𝒟)0`$ is the so-called quasi-effective threshold, and $`A`$ is a divisor on $`S`$. The behaviour of linear systems under birational maps between Mori fibrations is given by the following theorem (, theorem 4.2):
###### Proposition 1.2.
Let $`\rho :VS`$ and $`\rho ^{}:V^{}S^{}`$ be Mori fibrations, $`\chi :VV^{}`$ a birational map, $`𝒟^{}=|n^{}(K_V^{})+\rho ^{}{}_{}{}^{}(A^{})|`$ a very ample linear system on $`V^{}`$, where $`A^{}`$ is an ample divisor on $`S^{}`$. Denote $`𝒟=\chi _{}^1𝒟^{}|n(K_V)+\rho ^{}(A)|`$. Then:
* $`nn^{}`$ always, and if $`n=n^{}`$ then $`\chi `$ is square;
* if a log pair $`K_V+\frac{1}{n}𝒟`$ is canonical and numerically effective, then $`\chi `$ is an isomorphism inducing also an isomorphism between the bases.
The statement $`(i)`$ of the proposition allows us to describe various Mori structures on $`V`$, while $`(\mathrm{𝑖𝑖})`$ can be used to get Mori models (i.e., various Mori fibrations birational to $`V`$).
There is another important number arising from the maximal singularities method (now playing the key role in the birational classification). This is so-called the adjunction threshold.
###### Definition 1.3.
A non-negative number $`\alpha (V,𝒟)`$ is called the adjunction threshold of a pair $`(V,𝒟)`$, where $`𝒟`$ is a non-empty linear system without fixed components, $`V`$ is projective and non-singular in codimension 1, if $`\alpha (V,𝒟)`$ is the smallest number such that for any positive integers $`m`$ and $`n`$ with $`\frac{m}{n}>\alpha (V,𝒟)`$ the linear system $`|nD+mK_V|`$ is empty, where $`D`$ is an element of $`𝒟`$.
The following statement clarifies the role of the adjunction threshold (, 2.1.):
###### Proposition 1.4.
Let $`\chi :VV^{}`$ be a birational map between terminal $``$-factorial varieties, $`𝒟^{}`$ a non-empty linear system on $`V^{}`$ without fixed components, and $`𝒟=\chi _{}^1𝒟^{}`$ its strict transform. If $`\alpha (V,𝒟)>\alpha (V^{},𝒟^{})`$, then a pair $`K_V+\frac{1}{\alpha (V,𝒟)}𝒟`$ is not canonical (in other words, $`𝒟`$ has a maximal singularity)
This condition ”to be not canonical” is strong enough, and we can often either prove its impossibility, or decrease the adjunction threshold applying an ”untwisting” birational automorphism or jumping to another Mori model (as in the Sarkisov program). A relation between both the thresholds is given by the following obvious lemma:
###### Lemma 1.5.
Let $`𝒟|n(K_V)+\rho ^{}(A)|`$ be a non-empty linear system without fixed components on a Mori fibration $`\rho :VS`$. Always $`n=\mu (V,𝒟)\alpha (V,𝒟)`$, and if $`A`$ is effective, then the equality holds; if $`C\rho ^{}(A)<0`$ for a general curve $`C`$ covering the base, then $`\mu (V,𝒟)>\alpha (V,𝒟)`$.
The key idea is that the relation between these thresholds determines whether a variety is rigid. For Mori fibrations that are pencils of del Pezzo surfaces, it looks as follows:
###### Conjecture 1.6.
Let $`V/^1`$ be a pencil of del Pezzo surfaces of degree 1,2 or 3 (we assume it to be a Mori fibration). $`V/^1`$ is birationally rigid if and only if $`\alpha (V,𝒟)=\mu (V,𝒟)`$ for any linear system without fixed components (i.e., a linear system $`|n(K_V)F|`$ either is empty, or has a base divisor; here $`F`$ is the class of a fiber).
Let us note that ”nearly all” (smooth) pencils of del Pezzos of degree 1, 2 or 3 are sufficiently ”twisted” along the base, and in their rigidity is proved. In this paper, we deal with all cases of degree 1 (theorem 2.6) and advance in degree 2 (theorem 3.1). In particular, conjecture 1.6 turns out to be true for smooth cases of degree 1.
The paper is based on using the maximal singularities method (, ) in its most perfect kind .
Everywhere in this paper the characteristic of the ground field is assumed to be 0.
The author would like to thank the Directors and the stuff of the Max-Planck Institute for Mathematics in Bonn for hospitality and excellent conditions during the work.
## 2. Smooth varieties with a pencil of del Pezzo surfaces of degree 1.
### 2.1. The essential construction.
We follow to the idea given in .
Let $`\rho :V^1`$ be a smooth Mori fibration on del Pezzo surfaces of degree 1, $`PicV=[K_V][F]`$, where $`F`$ is the class of a fiber. By the Grauert theorem, $`\rho _{}𝒪(2K_V+mF)`$ is a vector bundle of rank 4 on $`^1`$. We may choose $`m`$ such that
$$=\rho _{}𝒪(2K_V+mF)𝒪𝒪(n_1)𝒪(n_2)𝒪(n_3),$$
where $`0n_1n_2n_3`$. Suppose $`b=n_1+n_2+n_3`$.
Note that there exists a special section $`s_B`$:
$$s_B=Bas|K_V+kF|$$
for $`k0`$. Obviously, any fiber $`S|F|`$ is smooth at the point $`s_BS`$, which is also a unique base point of $`|K_S|`$.
Let $`\pi :X^1`$ be a natural projection for $`X\stackrel{\mathrm{def}}{=}𝐏𝐫𝐨𝐣`$, $`\phi :VX`$ a morphism defined by $`|2K_V+mF|`$. Then $`\rho =\pi \phi `$. It is easy to see that if $`QX`$ is the image of $`\phi `$, then $`\phi :VQ`$ is a degree 2 finite morphism branched along a smooth divisor $`R_Q=RQ`$ that does not intersect $`t_B\stackrel{\mathrm{def}}{=}\phi (s_B)`$, where $`RX`$ is a divisor.
$`Q`$ is restricted to any fiber $`T`$ of $`X/^1`$ as a non-degenerated quadratic cone with $`Tt_B`$ as its vertex, and $`R|_T`$ is a cubic surface. So if $`M`$ is the class of the tautological bundle on $`X`$ and $`L`$ is the class of a fiber, we have
$$\begin{array}{cc}Q=2M+aL;& R=3M+cL.\end{array}$$
Let $`t_0`$ be a section corresponding to the surjection $`𝒪0`$, $`l`$ the class of a line in $`L^3`$; set $`t_Bt_0+\epsilon l`$, where $`\epsilon `$ is obviously nonnegative.
Choose a very ample divisor $`HM+\beta l`$, $`\beta >0`$. So $`D=HQ`$ is a conic bundle with $`N=t_BH`$ degenerations. We may suppose that all fibers of D are reduced and all its singular points are du Val $`𝐀_\mathrm{𝟏}`$. It is easy to compute that $`K_D^2=82b2\beta 3a`$. Let us blow up all singular points of $`D`$ and then contract all (-1)-curves (there will be $`2N`$ of them) onto a ruled surface $`D^{}`$. We have
$$K_D^2+2N=K_D^{}^2=8,$$
hence
$$N=b+\beta +\frac{3}{2}a.$$
Note that $`a`$ must be even. Set $`a=2a^{}`$. Then,
$$N=tH=(t_0+\epsilon l)(M+\beta l)=\beta +\epsilon ,$$
and we obtain $`\epsilon =b+3a^{}`$. Since $`R|_Qt_B=\mathrm{}`$, we have $`0=(t_0+\epsilon l)(3M+cL)=3\epsilon +c`$. So, we have proved the following lemma:
###### Lemma 2.1.
The following equalities hold:
$$\begin{array}{c}Q=2M\frac{2}{3}(b\epsilon )L,\hfill \\ t_B=t_0+\epsilon l,\hfill \\ R=3M3\epsilon L.\hfill \end{array}$$
In the sequel we will always assume $`b>0`$. Really, if $`b=0`$, then $`Q=2M`$ and $`\epsilon =0`$, whence $`V^1\times S`$ for some smooth del Pezzo surface $`S`$. So $`V`$ is rational and even not a Mori fibration because of the Picard group rank.
###### Lemma 2.2.
The only following two cases may occur:
* $`\epsilon =0`$, $`2n_2=n_1+n_3`$, $`n_1`$ and $`n_3`$ are even; so we have
$$\begin{array}{ccc}Q=2M2n_2L,& t_B=t_0,& R=3M;\end{array}$$
* $`\epsilon =n_1>0`$, $`n_3=2n_2`$, $`n_1`$ is even and $`n_23n_1`$; in this case
$$\begin{array}{ccc}Q=2M2n_2L,& t_B=t_0+n_1l,& R=3M3n_1L.\end{array}$$
Proof. Suppose first $`\epsilon >0`$. Then $`\epsilon n_1`$ because of irreducibility of $`t_B`$. Since $`R`$ is irreducible and curves of the class $`t_0+n_2l`$ sweep any effective divisor of the class $`Mn_3L`$, we have $`(t_0+n_2l)R=3n_23\epsilon 0`$, i.e., $`\epsilon n_2`$.
Let $`n_1=n_2=n_3`$. Since $`t_B(Mn_1L)=0`$, the linear system $`|Mn_1L|`$ is two-dimensional, and any of its elements is irreducible, there exists at least one-dimensional subsystem in it that contains $`t_B`$. Let $`C|Mn_1Lt_B|`$ be a general element, then
$$QC=S+S^{}2M^22\left(\frac{1}{3}(b\epsilon )+n_1\right)ML,$$
where $`S`$ and $`S^{}`$ are ruled surfaces covering the base. Note that there exists only a unique curve of the class $`t_0`$, and $`t_Bt_0=\mathrm{}`$. At least one of these surfaces, say, $`S`$, containes both $`t_0`$ and $`t_B`$. It is easy to see that $`S`$ is of the class $`M^22n_1ML`$, whence $`S^{}M^2\frac{4}{3}n_1ML`$. Since $`SS^{}`$, then $`t_0S^{}`$, and denoting by $`P`$ the intersection of two general elements of $`|Mn_1L|`$, we can see that $`dim(S^{}P)=0`$ (note that $`Bas|Mn_1L|=t_0`$). We get a contradiction:
$$0S^{}P=(M^2\frac{4}{3}n_1ML)(M^22n_1ML)=\frac{n_1}{3}<0.$$
So if $`\epsilon n_1`$ then $`n_1<n_3`$; moreover, $`n_1>0`$ because of irreducibility of $`R_Q`$ (otherwise, if $`n_1=0`$, there should be at least 1-dimensional family of lines of the class $`t_0`$ sweeping a surface, and $`R_Q`$ should contain it).
Let $`G`$ be an irreducible ruled surface that is the complete intersection of general elements of $`|Mn_2L|`$ and $`|Mn_3L|`$. Let us show that $`t_BG`$. Indeed, $`G`$ has the type $`𝐅_{𝐧_\mathrm{𝟏}}`$ with $`t_0`$ as its exceptional section. Assume the converse, i.e., $`t_BG`$. Then, for a general fiber $`S|L|`$, the line $`G|_S`$ does not lie on $`Q|_S`$. So $`GQ`$, and we may suppose that $`GQ=t_0+C`$, where $`C`$ is a curve with no $`t_0`$ as a component. It is easy to compute that
$$Ct_0+2\left(bn_2n_3\frac{b\epsilon }{3}\right)l,$$
and
$$2\left(bn_2n_3\frac{b\epsilon }{3}\right)n_1,$$
which, taking into account $`\epsilon n_2`$, gets a contradiction: $`n_1>\frac{4}{3}n_1`$.
So $`t_BG`$, hence $`n_2<n_3`$. In fact, the equality should led $`dim|Mn_3L|`$, thus $`Bas|Mn_3L|=t_0`$, and we could always find $`G`$ such that $`t_BG`$.
By $`T`$ we denote a unique effective element of $`|Mn_3L|`$. We show that $`T|_Q=2G`$. Indeed, assuming the converse, i.e., $`t|_Q=G+G^{}`$, $`GG^{}`$, and taking into account that $`t_BGG^{}`$ and $`t_0G`$, we see $`t_0G^{}`$. Then
$$G^{}|_G=t_0+\text{\#}(t_0t_B=\epsilon n_1)\text{ fibers}.$$
Choose a general $`S|Mn_1L|`$. It is clear that $`S|_T`$ is a ruled surface of the type $`𝐅_{𝐧_\mathrm{𝟐}}`$ with the exceptional section $`t_0`$, and $`GS|_T=t_0`$. Let $`C=G^{}S|_T`$ be an irreducible curve. We have $`Ct_0=\epsilon n_1`$, $`QTS=t_0+C`$, whence
$$Ct_0+2\left(n_2\frac{b\epsilon }{3}\right)l.$$
This yields $`Ct_0=n_2\frac{2}{3}(b\epsilon )=\epsilon n_1`$, and we get a contradiction again:
$$\epsilon =n_1+n_22n_3<0.$$
So we have
$$TQ=2M^22\left(n_3+\frac{b\epsilon }{3}\right)ML=2G=2(M^2(n_2+n_3)ML),$$
thus $`2n_2=n_1+n_3\epsilon `$. If $`\epsilon `$ were greater than $`n_1`$, then $`R_Q`$ should contain $`G`$, because $`G`$ is covered by curves of the class $`t_0+n_1l`$. This contradicts to irreducibility of $`R_Q`$. So $`\epsilon =n_1`$ and then $`2n_2=n_3`$. In order to show $`n_23n_1`$, we note that $`S_Q=SQ`$ is a conic bundle without degenerations, i.e., a ruled surface of the type $`𝐅_{𝐧_\mathrm{𝟐}}`$. Then $`R`$ cuts off an effective curve on $`S_Q`$ that does not contain $`t_0`$, and $`t_0`$ itself. We get $`3n_33n_15n_2`$, which yields $`n_23n_1`$.
It remains to show that $`n_1`$ is even. Let $`Ct_0+n_1l`$ be a general irreducible curve on $`G`$, and $`C_V`$ its inverse image on $`V`$. Then $`C_V`$ covers $`C`$ with the ramification divisor of degree $`t_BC=n_1`$. Clearly, this degree must be even. This completes the case $`\epsilon >0`$.
Now let $`\epsilon =0`$. We may assume $`n_1<n_3`$. Take $`T|Mn_3L|`$. It is easy to observe that $`T=𝐏𝐫𝐨𝐣_T`$, where
$$_T=𝒪𝒪(n_1)𝒪(n_2).$$
Since $`T`$ can be covered by curves of the class $`t_0+n_2l`$ and $`Q2M\frac{2}{3}bL`$ is irreducible, then $`n_2\frac{b}{3}`$, i.e., $`2n_2n_1+n_3`$, and $`n_2>n_1`$ since $`n_3>n_1`$.
Further, let $`Ct_0+n_1l`$ be irreducible. Since $`b=n_1+n_2+n_3>3n_1`$, then $`CQ=2(n_1\frac{b}{3})<0`$, and therefore $`CQ`$. A ruled surface $`G(Mn_2L)(Mn_3L)=M^2(n_2+n_3)ML`$ has the type $`𝐅_{𝐧_\mathrm{𝟐}}`$ and can be covered by curves of the same class as $`C`$, so $`GQ`$. Therefore, $`T|_Q=G+G^{}`$ for some $`G^{}`$.
Let $`M_T=M|_T`$ be the tautological divisor on $`T`$, $`L_T=L|_T`$. Then $`G`$ has the class $`M_Tn_2L_T`$ in $`T`$, $`Q|_T2M_T\frac{2}{3}bL_T`$, whence $`G^{}M_T(\frac{2}{3}bn_2)L_T`$. Note that $`\frac{2}{3}bn_2>n_1`$ since $`n_1<n_2`$.
But $`T`$ has only one divisor of the form $`M_TxL_T`$ for $`x>n_1`$: this is $`M_Tn_2L_T`$. So $`\frac{2}{3}bn_2=n_2`$, hence $`n_1+n_3=2n_2`$.
Now, let $`CQ`$ be a general irreducible curve of the class $`t_0+n_1l`$, $`C_V`$ its inverse image on $`V`$. Then $`C_V`$ is the double cover of $`C`$ branched along a divisor of degree $`RC=3n_1`$, whence $`n_1`$ and $`n_3`$ are even. Lemma 2.2 is proved.
Such a construction of $`V`$ allows us to deal completely with intersections of cycles on $`V`$. Denote $`G_V`$, $`F`$, and $`H`$ the inverse images of $`G(Mn_2L)(Mn_3L)`$, $`L_Q=L|_Q`$, and $`M_Q=M|_Q`$ on $`V`$ respectively. Then, let $`2s_0`$ be the pull back of $`t_0`$, and $`2f`$ the pull back of a line in a fiber of $`L_Q`$.
###### Lemma 2.3.
$`\overline{\mathrm{𝐍𝐄}}(V)=\mathrm{𝐍𝐄}(V)=_+[s_0]_+[f]`$, and:
* if $`\epsilon =0`$, then $`K_V=G_V+(\frac{1}{2}n_12)F`$, $`H=2G_V+n_3F`$, $`K_V^2=s_0+(4n_2)f`$, $`s_0F=fG_V=1`$, $`s_0G_V=\frac{1}{2}n_3`$, $`fF=0`$.
* if $`\epsilon =n_1>0`$, then $`K_V=G_V(\frac{1}{2}n_1+2)F`$, $`H=2(G_V+n_2F)`$, $`s=s_0+\frac{1}{2}n_1f`$, $`K_V^2=s_0+(4+\frac{3}{2}n_1n_2)f`$, $`s_0F=fG_V=1`$, $`s_0G_V=\frac{1}{2}n_3=n_2`$, $`fF=0`$.
Proof. It can be easily checked by the following way. First blow up $`t_B`$ on $`Q`$ with an exceptional divisor $`S`$; then take the double cover branched along a smooth divisor composed from the pre-image of $`R_Q`$ and $`S`$; after that, contract the pre-image of $`S`$ onto $`s_B`$.
###### Lemma 2.4.
Let $`\epsilon =0`$ (i.e., $`s_B=s_0`$). If $`|n(K_V)mF|`$, $`m>0`$, has no fixed components, then $`n_3=2`$.
Proof. Let $`D|n(K_V)mF|`$. Since
$$D|_{G_V}ns_0\left(m+n\left(\frac{n_1}{2}+n_3n_22\right)\right)f$$
is an effective curve, we obtain
$$0<\frac{m}{n}2+n_2n_3\frac{n_1}{2}=2\frac{n_3}{2}.$$
It only remains to take into account that $`n_3`$ is positive and even. The lemma is proved.
###### Lemma 2.5.
Let $`\epsilon =n_1>0`$ (so $`s_Bs_0`$). If $`|n(K_V)+mF|`$ has no fixed components, then $`m>0`$.
Proof. Let $`D|n(K_V)+mF|`$. Since
$$D|_{G_V}ns_0+[m+n(n_1n_2+2)]f$$
is an effective curve, we have $`2+n_1n_2+\frac{m}{n}0`$. Using $`n_23n_1`$, we see that $`2+n_1n_222n_1<2`$, whence $`m>0`$. The lemma is proved.
### 2.2. Results about rigidity.
###### Theorem 2.6.
Conjecture 1.6 is true for smooth Mori fibrations on del Pezzo surfaces of degree 1 (over $`^1`$).
In order to prove this theorem, we shall use the following proposition:
###### Proposition 2.7.
Let $`V/^1`$ be a smooth Mori fibration on del Pezzo surfaces of degree 1, $`V^{}/S^{}`$ another Mori fibration, $`\chi :VV^{}`$ a birational map. Suppose that
$$𝒟^{}=|n^{}(K_V^{})+\text{pull back of an ample divisor from the base }|$$
is a very ample linear system. Suppose aslo
$$𝒟=\chi _{}^1𝒟^{}|n(K_V)+mF|$$
for some $`m0`$.
Then $`n=n^{}`$ (so $`\chi `$ is birational over the base), possibly except the case $`\epsilon =n_1=0`$, $`n_2=1`$, $`n_3=2`$.
Proof of the proposition. By lemma 1.5, $`\mu (V^{},𝒟^{})=n^{}=\alpha (V^{},𝒟^{})`$ and $`\mu (V,𝒟)=n=\alpha (V,𝒟)`$. If $`n=n^{}`$, the assertion follows from proposition 1.2.
Suppose $`n>n^{}`$. Then a log pair $`K_V+\frac{1}{n}𝒟`$ is not canonical, i.e., the linear system $`𝒟`$ has a maximal singularity over some point. But using , we may say a little more.
Namely, let $`D_1`$, $`D_2`$ be general elements of $`𝒟`$. We observe that
$$D_1D_2=\{\begin{array}{cc}n^2s_0+((4n_2)n^2+2mn)f,\hfill & \text{if }\epsilon =0\hfill \\ n^2s_0+((4+\frac{3}{2}n_1n_2)n^2+2mn)f,\hfill & \text{if }\epsilon =n_1>0\hfill \end{array}$$
Put down $`D_1D_2=Z^h+Z^v`$, where $`Z^h`$ and $`Z^v`$ are horizontal and vertical cycles respectively. Then there exist a geometric discrete valuation $`\nu `$ centered at a point $`B_0`$ over some point $`t^1`$, and a positive number $`e=\nu (𝒟)n\delta `$, where $`\delta `$ is the canonical multiplicity with respect to $`\nu `$, such that for the component $`Z_t^vF_t`$ of the cycle $`Z^v`$ in the fiber $`F_t`$ over the point $`t`$ we have
$$\mathrm{deg}Z_t^v<\{\begin{array}{cc}(4n_2)n^2+2n\frac{e}{\nu (F_t)},\hfill & \text{if }\epsilon =0\hfill \\ (4+\frac{3}{2}n_1n_2)n^2+2n\frac{e}{\nu (F_t)},\hfill & \text{if }\epsilon =n_1>0,\hfill \end{array}$$
where $`\mathrm{deg}Z_t^v\stackrel{\mathrm{def}}{=}Z_t^v(K_V)`$ (this is so-called the supermaximal singularity condition, in the notions of ). Besides, there exist positive numbers $`\mathrm{\Sigma }_0`$, $`\mathrm{\Sigma }_0^{}\mathrm{\Sigma }_0`$, and $`\mathrm{\Sigma }_1`$ such that multiplicities $`m^h`$ and $`m^v`$ of the cycles $`Z^h`$ and $`Z_t^v`$ at $`B_0`$ satisfy
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)(m^h\mathrm{\Sigma }_0+m^v\mathrm{\Sigma }_0^{})>(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+e)^2.$$
Note that $`m^h\mathrm{deg}Z^h\stackrel{\mathrm{def}}{=}Z^hF=n^2`$ and $`m^v2\mathrm{deg}Z_t^v`$. So, in the case $`\epsilon =n_1>0`$ we have
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)((9+3n_12n_2)n^2\mathrm{\Sigma }_0+4ne)>(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+e)^2,$$
hence
$$(2n_23n_15)\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)+(n\mathrm{\Sigma }_1e)^2<0,$$
i.e., $`2n_2<5+3n_1`$. Taking into account the condition $`n_23n_1`$, we obtain $`n_11`$, which is impossible because $`n_1`$ must be even.
Now let $`\epsilon =0`$. Then
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)((82n_2)n^2\mathrm{\Sigma }_0+4ne)>(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+e)^2,$$
whence
$$(2n_25)\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)+(n\mathrm{\Sigma }_1e)^2<0.$$
So $`n_2<\frac{5}{2}`$, so the only three cases are possible:
* $`n_1=0`$, $`n_2=2`$, $`n_3=4`$;
* $`n_1=2`$, $`n_2=2`$, $`n_3=2`$;
* $`n_1=0`$, $`n_2=1`$, $`n_3=2`$.
We show that the first two cases can not occur. Suppose first that there are no sections of the class $`s_0`$ through the point $`B_0`$. Let $`Z_1^h`$ be the sum of all horizontal cycles through $`B_0`$. So $`Z_1^hAs_0+Bf`$ for some $`BA`$. Since $`n_2=2`$ in the considering cases, we have
$$\underset{p^1}{}\mathrm{deg}Z_p^v2n^2+2mnB2n^2+2mnA,$$
which implies the following inquality for the supermaximal singularity:
$$m^v2\mathrm{deg}Z_t^v<4n^2+4n\frac{e}{\nu (F_t)}2A.$$
But we know that $`m^hA`$, and then
(2.1)
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)(A\mathrm{\Sigma }_0+(4n^22A)\mathrm{\Sigma }_0+4ne)>(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+e)^2,$$
so we get a contradiction:
$$A\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)+(n\mathrm{\Sigma }_1e)^2<0.$$
It remains to prove that there are no maximal singularities over points that lie on sections of the class $`s_0`$.
Consider the case $`n_1=n_2=n_3=2`$. Note that $`s_Bs_0`$ is a unique section of such a class in this case. If $`𝒟`$ has a maximal singularity over the point $`B_0s_B`$, then for a general $`D𝒟`$ it holds $`\nu _1=mult_{B_0}D>n`$, and for a general curve $`Cf`$ in the fiber through $`B_0`$ we get $`n=DC\nu _1>n`$, which is impossible.
Now let $`n_1=0`$, $`n_2=2`$, $`n_3=4`$. Proving of this case consists of two lemmas below and an observation that any curve of the class $`s_0`$ lies on the divisor $`G_V`$. Recall that the valuation $`\nu `$ can be realized as a chain of blow-ups centered in nonsingular centers $`B_0,B_1,\mathrm{}`$; a general $`D𝒟`$ has multiplicities $`\nu _1,\nu _2,\mathrm{}`$ in these centers respectively.
###### Lemma 2.8.
Let $`𝒟|n(K_V)+mF|`$, $`m0`$, has a maximal singularity over a point $`B_0S`$, where $`S`$ is a fiber of $`V/^1`$. By $`l|K_S|`$ we denote the curve through $`B_0`$. Then
* $`S`$ is smooth at $`B_0`$;
* $`B_0`$ is a double point of $`l`$;
* $`\nu `$ defines an infinitely near singularity, and if $`B_1S^1\mathrm{}`$ (upper indices denote the strict transform on the correspondig floor of the chain of blow-ups realizing $`\nu `$), then $`B_1l^1\mathrm{}`$ and $`\nu _1+\nu _23n`$.
Proof. If $`S`$ is singular at $`B_0`$, then a general curve $`C|2(K_S)B_0|`$ is also singular at this point, so we get a contradiction:
$$2n=DC2\nu _1>\nu _1.$$
Further, suppose $`l`$ is nonsingular at $`B_0`$. Set
$$D|_Skl+C,$$
where $`lSuppC`$. We have
$$n<\nu _1mult_{B_0}D|_S=k+mult_{B_0}C.$$
But $`mult_{B_0}CCl=(n(K_S)kl)l=nk`$, and then
$$n<\nu _1k+nk=n.$$
So $`l`$ has a double point at $`B_0`$. Now we show that $`\nu `$ defines an infinitely near singularity. Assuming the converse, we get $`\nu _1>2n`$. Supposing $`D|_Skl+C`$, we see that
$$2mult_{B_0}CCl=nk,$$
whence
$$2n<\nu _12k+mult_{B_0}C\frac{3}{2}k+\frac{1}{2}n,$$
i.e., $`k>n`$, a contradiction.
Now let $`B_1S^1\mathrm{}`$, but $`B_1l^1=\mathrm{}`$. Denote $`E_1`$ the exceptional divisor of the blow-up of $`B_0`$. Then
$$D^1|_{S^1}=(DS)^1+mE_1|_{S^1},$$
and for the decomposition $`D|_Skl+C`$ we get
$$\stackrel{~}{\nu }_1+\stackrel{~}{\nu }_22k+mult_{B_0}C+mult_{\stackrel{~}{B}_1}C^1+m,$$
where $`\stackrel{~}{B}_1=B_1S^1`$, $`\stackrel{~}{\nu }_1=mult_{B_0}D|_S`$, $`\stackrel{~}{\nu }_2=mult_{\stackrel{~}{B}_1}(D|_S)^1`$. Then,
$$mult_{\stackrel{~}{B}_1}C^1mult_{B_0}C\frac{1}{2}(n(K_S)kl)l=\frac{1}{2}(nk),$$
so
$$\stackrel{~}{\nu }_1+\stackrel{~}{\nu }_22k+2\frac{nk}{2}+m=n+k+m.$$
But $`\stackrel{~}{\nu }_1\nu _1+m`$ and $`\stackrel{~}{\nu }_2\nu _2`$, thus
$$2n<\nu _1+\nu _2n+k,$$
i.e., $`k>n`$, which is impossible.
So $`B_1l^1\mathrm{}`$. In order to show that $`\nu _1+\nu _23n`$, we argue as before, only taking into account that
$$\stackrel{~}{\nu }_1+\stackrel{~}{\nu }_23k+mult_{B_0}C+mult_{\stackrel{~}{B}_1}C^1+m.$$
We have $`\nu _1+\nu _22k+n3n`$. The lemma is proved.
###### Lemma 2.9.
Let $`\epsilon =0`$, $`n_1=0`$. Then $`G_VC\times ^1`$, where $`C`$ is an elliptic curve.
Proof. We use the notation of section 2.1. It is obvious that $`G_VC\times ^1`$. We shall prove the smoothness of $`C`$.
Notice that $`R|_G3t_0`$. Let $`B`$ be an irreducible component of $`SuppR|_G`$. Clearly, $`Bt_0`$. Let $`\psi :\stackrel{~}{V}V`$ be the blow-up of $`B`$ with the exceptional divisor $`E`$. $`E`$ is a ruled surface of the type $`𝐅_{𝐧_\mathrm{𝟐}}`$. Denote $`t_E`$ and $`l_E`$ the classes of the exceptional section and a fiber of $`E`$ respectively. Then for strict transforms of the divisors $`G`$ and $`R`$ on $`\stackrel{~}{V}`$ we have (using their smoothness along $`B`$)
$$\begin{array}{cc}\stackrel{~}{G}|_Et_E,& \stackrel{~}{R}|_Et_E+n_2l_E,\end{array}$$
hence $`\stackrel{~}{G}|_E\stackrel{~}{R}|_E=\mathrm{}`$. So any fiber of $`G^1\times ^1`$ meets $`R`$ transversally, i.e., at three different points. Thus, $`C`$ is nonsingular.
The lemma is proved. This completes also the proof of proposition 2.7.
###### Remark 2.10.
In fact, proposition 2.7 works also in the case $`n_1=0`$, $`n_2=1`$, $`n_3=2`$, but the proof is more complicated. Since this case is non-rigid, satisfies theorem 2.6, and can not be advanced to a more complete description yet (see below), we omit proving of it.
###### Corollary 2.11.
For $`\epsilon =n_1>0`$, all smooth Mori fibrations on del Pezzo surfaces of degree 1 over $`^1`$ are rigid over the base; for $`\epsilon =0`$, all cases are also rigid, except the following:
* $`n_1=n_2=n_3=2`$;
* $`n_1=0`$, $`n_2=1`$, $`n_3=2`$.
Proof. Rigidity follows from proposition 2.7, using lemma 2.4 in the case $`\epsilon =0`$ or lemma 2.5 when $`\epsilon =n_1>0`$.
We will deal with non-rigid cases after the following remark.
###### Remark 2.12.
In it was proved that (smooth) pencils of del Pezzo surfaces (of degree 1, 2 or 3) are all rigid, if the so-called $`𝐊^\mathrm{𝟐}`$-condition holds:
cycles $`aK^2bf`$ are not effective for any $`a,b>0`$.
For degree 1, this condition is satisfied exactly when $`n_24`$ (for $`\epsilon =0`$) or $`n_24+\frac{3}{2}n_1`$ (for $`\epsilon =n_1>0`$). The second is always true except $`n_1=2`$, $`n_2=6`$, $`n_3=12`$.
Another sufficient condition of rigidity was proposed in (conjecture 1.2):
$`(K_V)^3+m_0+12`$,
where
$$m_0=\mathrm{min}\{r:(K_V)+rF\text{ is nef}\}.$$
It holds exactly when $`n_23`$ ($`\epsilon =0`$) or $`n_22+\frac{3}{2}n_1`$ ($`\epsilon =n_1>0`$). Notice that the second is always true.
Thus, the second sufficient condition is more exact. Both these conditions do not include the rigid case $`n_1=0`$, $`n_2=2`$, $`n_3=4`$ for $`\epsilon =0`$.
Now we consider the non-rigid cases. We start with the case $`n_1=0`$, $`n_2=1`$, $`n_3=2`$. It was already considered in some works (see , section 2, §2, and ). It is easy to see that a unique effective divisor of the class $`G_V`$ is the direct product of $`^1`$ and an elliptic curve, the linear system $`=|2G_V+2F|=|2(K_V)2F|`$ is base points free and defines a morphism contracting $`G_V`$ along $`^1`$ onto a curve $`l`$ on a Fano variety $`U`$ of index 2 and degree 5 (this is so-called the double cone over the Veronese surface). Indeed, the linear system $`|M_Q|`$ maps $`Q`$ onto a cone $`\stackrel{~}{Q}^6`$ over the Veronese surface in $`^5`$, and $`U`$ is obtained by taking the double cover branched along a cubic section that does not pass through the cone vertex. $`l`$ covers one of the generators of $`\stackrel{~}{Q}`$. Let
$$\mu _l:V=V_lU$$
be such a contraction. Notice that $`K_VG_V+2F`$, and $`V`$ is a Fano itself. Thus, $`V/^1`$ is not rigid, and $`1=\alpha (V,)<\mu (V,)=2`$, i.e., the conditions of conjecture 1.6 hold. Notice that $`U`$ has a lot of such structures. Unfortunately, the technique of the maximal singularities method is not enough yet to deal completely with this case (for example, it fails when $`l`$ has a double point; see also remark 2.10). Nevertheless, we can formulate the following conjecture, which seems to be true:
###### Conjecture 2.13.
Any smooth Mori fibration in the class of birational equivalency of $`U`$ is biregular to either $`U`$ or $`V_l`$ for some $`l`$.
The remaining case is $`\epsilon =0`$, $`n_1=n_2=n_3=2`$. Let $`V/^1`$ be the corresponding pencil of del Pezzo surfaces of degree 1. Notice that
$$𝒩_{s_0|V}𝒪(1)𝒪(1)$$
and $`|mH|`$ for $`m0`$ gives a contraction of $`s_0`$. So there exists a flop
$$\mu :VU$$
centered at $`s_0`$. Then, $`dim|G_V|=1`$:
$$2dim|G_V|=2dim|G|dim|2G|=dim\left|(M2L)|_Q\right|=2.$$
Let $`G_U=\overline{\mu (F)}`$, and $`F_U`$ the class of a fiber of $`U/^1`$. Obviously, $`\mu ^1|G_U|=|F|`$ and $`\mu ^1|F_U|=|G_V|`$. Moreover, $`U/^1`$ is a pencil of del Pezzo surfaces of degree 1 with the same structure parameters $`n_1=n_2=n_3=2`$.
###### Proposition 2.14.
Let $`\chi :VW`$ be a birational map onto a Mori fibration $`\gamma :WS`$. Then either $`\chi `$ or $`\chi \mu ^1:UW`$ are birational over the base. Moreover, $`U`$ and $`V`$ are unique smooth Mori models in their class of birational equivalency; $`Bir(V)=Aut(V)`$, and in general case, $`Bir(V)=<\sigma >_2`$, where $`\sigma `$ is the double cover involution.
Proof. Let $`𝒟_W=|n^{}(K_W)+\gamma ^{}(A)|`$ be a very ample linear system, where $`A`$ is an ample divisor on $`S`$, and $`𝒟_V=\chi _{}^1𝒟_W|n(K_W)+mF|`$.
Suppose $`m0`$. Then from proposition 1.2 it follows that $`n=n^{}`$ and $`\chi `$ is birational over the base. So we may assume $`m=l<0`$. Denote $`𝒟_U=\mu ^1𝒟_V`$. Since $`K_U=G_U+F_U`$, we have
$$𝒟_U|(nl)(K_U)+lF_U|,$$
and $`\chi \mu ^1`$ is birational over the base by proposition 2.7.
Now let $`\chi :VW`$ be a birational map onto a Mori fibration $`W/^1`$ with $`K_W`$ to be nef. We may assume $`\chi `$ to be birational over the base. Suppose $`𝒴=|n(K_W)+m^{}F_W|`$ is a very ample linear system, $`m^{}>0`$, and $`𝒟=\chi _{}^1𝒴|n(K_V)+mF|`$ for some $`m0`$. If $`K_V+\frac{1}{n}𝒟`$ were canonical, i.e., $`D`$ had no maximal singularities, then $`\chi `$ would be an isomorphism by proposition 1.2.
So let $`𝒟`$ have maximal singularities. We may assume that they are all infinitely near (see ). By $`Y`$ and $`D`$ we denote general elements of $`𝒴`$ and $`𝒟`$ respectively. Choose a resolution of singularities of $`\chi `$:
$$\begin{array}{ccccccccc}& & & & Z& & & & \\ & & \stackrel{\phi }{}& & & & \stackrel{\psi }{}\\ V& & & & \stackrel{\chi }{}& & W\end{array}$$
Notice that
$$nK_Z+\psi ^1Y=m^{}\psi ^{}(F_W)+a_iE_i^{}$$
and
$$nK_Z+\phi ^1D=m\phi ^{}(F)+b_iE_i,$$
where $`E_i^{},E_i`$ are exceptional divisors, $`a_i`$ are all positive. Since $`𝒟`$ has maximal singularities, then
$$=\{i:b_i<0\}$$
is not empty. For a point $`t^1`$ we denote $`_t=\{i:\phi (E_i)F_t\}`$.
Obviously, $`dim|nK_Z+\psi ^1Y|=m^{}`$, so
$$dim|m\phi ^{}(F)+b_iE_i|=dim|(m^{}+mm^{})\phi ^{}(F)+b_iE_i|=m^{}.$$
Denote $`I=\{t^1:_t\mathrm{}\}`$. Then there exists a split
$$mm^{}=\underset{tI}{}k_t,$$
where $`k_t`$ are all positive, such that
$$dim|\underset{tI}{}k_t\phi ^1(F_t)+\underset{tI}{}\underset{i_t}{}(k_tc_ib_i)E_i|=0,$$
where for $`tI`$ positive numbers $`c_i`$ are defined by
$$\phi ^{}(F_t)=\phi _{}^1(F_t)+\underset{i_t}{}c_iE_i.$$
So for any $`tI`$ and $`i_t`$ we have $`k_tc_ib_i0`$, and, for every i, $`k_i`$ is the smallest positive number with such a property. Thus
$$k_t=\underset{i_t}{\mathrm{max}}\mathrm{}b_ic_i^1\mathrm{},$$
where $`\mathrm{}\mathrm{}`$ is the round up.
Let $`D_1,D_2𝒟`$ be general. We can put down
$$D_1D_2=Z^h+Z^vn^2s+(2n^2+2mn)f,$$
where $`Z^v`$ and $`Z^h`$ are vertical and horizontal cycles. Since $`s_0`$ is unique in its class and centers of maximal singularities can not lie on $`s_0`$, then horizontal cycles being caught by at least one of these centers can put down as $`As_0+Bf`$ for some $`BA`$. Taking into account that
$$mm^{}+\underset{tI}{}k_t,$$
we get an estimation
$$\mathrm{deg}Z^v=\underset{t^1}{}\mathrm{deg}Z_t^v2n^2+2mnA2n^2+2m^{}n+2n\underset{tI}{}k_tA$$
(the degrees of vertical cycles are estimated by intersection with $`K_V`$). Thus there exist $`pI`$ and $`j_p`$ such that
$$\mathrm{deg}Z_p^v2n^2+2m^{}n+2n\mathrm{}b_jc_j^1\mathrm{}A.$$
As in , §4, there exist positive numbers $`\mathrm{\Sigma }_0`$, $`\mathrm{\Sigma }_1`$, and $`\mathrm{\Sigma }_0^{}`$, where $`\mathrm{\Sigma }_0\mathrm{\Sigma }_0^{}`$ and $`c_j\mathrm{\Sigma }_0^{}`$, such that the following inequality for multiplicities $`m^h=mult_{B_0}Z^h`$ and $`m^v=mult_{B_0}Z_j^v`$ of horizontal and vertical cycles at $`B_0=\phi (E_j)`$ holds:
(2.2)
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)(\mathrm{\Sigma }_0m^h+\mathrm{\Sigma }_0^{}m^v)(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+b_j)^2.$$
Notice that $`\mathrm{}b_jc_j^1\mathrm{}\mathrm{\Sigma }_0b_j+\mathrm{\Sigma }_0`$. By lemma 2.8, $`𝒟`$ has an infinitely near singularity over the point $`B_0F_p`$. Consider two cases.
1. $`B_1F_p^1=\mathrm{}`$, i.e., the center of the second blow up in the chain realizing the maximal singularity, does not intersect the strict transform of the fiber. In this case, since $`B_1`$ is a point, we have $`\mathrm{\Sigma }_0^{}\frac{1}{2}\mathrm{\Sigma }_0`$. Then
$$\mathrm{\Sigma }_0^{}m^v2\mathrm{deg}Z_p^v\frac{\mathrm{\Sigma }_0}{2}(2n^2+2m^{}n+2nA)\mathrm{\Sigma }_0+2nb_j,$$
and, using (2.2) and an estimation $`m^hA`$, we get
$$n^22m^{}n2n<0,$$
i.e.,
(2.3)
$$\frac{m^{}}{n}>\frac{1}{2}\frac{1}{n}.$$
2. Now let $`B_1S^1\mathrm{}`$. Here we can not state that $`\mathrm{\Sigma }_0^{}\frac{1}{2}\mathrm{\Sigma }_0`$, but lemma 2.8 gives us a good estimation $`\nu _1+\nu _23n`$.
We know that $`m^hA`$ and
$$\mathrm{\Sigma }_0^{}m^v2\mathrm{deg}Z_p^v\mathrm{\Sigma }_0(4n^2+4m^{}n+4n2A)\mathrm{\Sigma }_0+4nb_j.$$
Substituting this in (2.2), we obtain
$$(A4m^{}n4n)\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)+(n\mathrm{\Sigma }_1b_j)^2<0.$$
Notice that a condition $`b_j>0`$ is nothing but the Noether-Fano inequality (): for some set of non-increasing numbers $`\{r_i\}`$, where $`r_i=\mathrm{\Sigma }_0+\mathrm{\Sigma }_1`$, it holds
$$b_j=r_i\nu _i2n\mathrm{\Sigma }_0n\mathrm{\Sigma }_1>0.$$
Applying $`\nu _1+\nu _23n`$, we get
$$(n\mathrm{\Sigma }_1b_j)^2\frac{1}{2}n^2\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1).$$
Thus
$$\frac{1}{2}n^2+A4m^{}n4n<0,$$
and since $`A0`$, than
(2.4)
$$\frac{m^{}}{n}>\frac{1}{8}\frac{1}{n}.$$
Comparing (2.3) and (2.4), we may assume that (2.4) holds always. It only remains to use the condition that $`K_W`$ is nef: we can choose $`𝒴`$ such that $`n`$ is big enough but $`\frac{m^{}}{n}`$ is as small as we want, and then get a contradiction because of (2.4).
Now, substituting $`U`$ for $`W`$, we obtain the statment about the birational automorphisms group. The uniqueness of smooth models of $`U`$ and $`V`$ follows from the description of all smooth (rigid and non-rigid) Mori fibrations on del Pezzo surfaces of degree 1 given above.
Proposition 2.14 is proved. This also completes corollary 2.11 and theorem 2.6.
## 3. Smooth varieties with a pencil of del Pezzo surfaces of degree 2.
### 3.1. The essential construction.
In this section we study smooth Mori fibrations on del Pezzo surfaces of degree 2. Let $`\rho :V^1`$ be such a fibration. Denoting $`F`$ the class of a fiber of $`\rho `$, we have
$$Pic(V)=[K_V][F]$$
and $`(K_V)^2F=2`$. Since $`\rho `$ is flat and $`K_V`$ is $`\rho `$-ample, for some integer $`m`$ we have
$$\rho _{}𝒪(K_V+mF)=,$$
where $`=𝒪𝒪(n_1)𝒪(n_2)`$ is a vector bundle of rank 3 over $`^1`$, $`0n_1n_2`$. Set
$$X=𝐏𝐫𝐨𝐣$$
with a natural projection $`\pi :X^1`$. By $`M`$ we denote the class of the tautological bundle on $`X`$, $`L`$ the class of a fiber. Then
$$Pic(X)=[M][L].$$
It is easy to see that there exists a double cover $`\phi :VX`$ branched along a smooth divisor $`R`$ of the class $`4M+2aL`$ such that
$$\rho =\pi \phi .$$
Further, let $`t_0`$ be the class of a section that corresponds to $`𝒪0`$, $`l`$ the class of a line in a fiber of $`\pi `$, $`s_0=\frac{1}{2}\phi ^{}(t_0)`$, and $`f=\frac{1}{2}\phi ^{}(l)`$. Then
$$\overline{\mathrm{𝐍𝐄}}(X)=\mathrm{𝐍𝐄}(X)=_+[t_0]_+[l],$$
$$\overline{\mathrm{𝐍𝐄}}(V)=\mathrm{𝐍𝐄}(V)=_+[s_0]_+[f].$$
Denote $`b=n_1+n_2`$, $`H=\phi ^{}(M)`$; clearly, $`F=\phi ^{}(L)`$. We have $`M^3=b`$, $`M^2=t_0+bl`$, $`K_X=3M+(b2)L`$, $`Mt_0=Ll=0`$, $`Ml=Lt_0=1`$, $`V`$: $`K_V=H+(a+b2)F`$, $`H^2=2s_0+2bf`$, $`HF=2f`$, $`Hs_0=Ff=0`$, $`Hf=Fs_0=1`$, $`(K_V)^2=2s_0+(84a2b)f`$, $`H^3=2b`$, $`(K_V)^3=126a4b`$ (see ).
Such varieties with $`K_V^2=2s_0+\beta f`$ for $`\beta 0`$ were first studied in . That was an exclusively important step in studying of geometry of Fano-fibered varieties.
### 3.2. Results about rigidity.
###### Theorem 3.1.
Let $`V/^1`$ be a smooth Mori fibration on del Pezzo surfaces of degree 2 over $`^1`$ with $`K_V^2=2s_0+\beta f`$, where $`\beta 2`$ (i.e., $`2a+b3`$). Then $`V/^1`$ is birationally rigid.
Proof. Let $`\chi :VV^{}`$ be a birational map onto a Mori-fibration $`\rho ^{}:V^{}S^{}`$, $`𝒟^{}=|n^{}K_V^{}+\rho ^{}{}_{}{}^{}A_{}^{}|`$ a linear system as in proposition 1.2, and $`𝒟=\chi _{}^1𝒟^{}|nK_V+mF|`$. We will denote $`D`$ a general element of $`𝒟`$. Obviously, $`𝒟`$ has no fixed components. If $`n=n^{}`$, the assertion follows from proposition 1.2. So let $`n>n^{}`$.
###### Lemma 3.2.
$`m0`$, i.e. $`\mu (𝒟)=\alpha (𝒟)`$.
Proof of the lemma. The cases $`\beta 0`$ are proved in (the so-called $`𝐊^\mathrm{𝟐}`$-condition holds in these cases). Let $`\beta =2`$. Then $`2a+b=3`$, $`b`$ is odd, so $`n_1<n_2`$. Clearly, $`\mu (𝒟)=n`$.
Assume the converse, i.e., $`m<0`$. Consider a general curve $`C2s_0+2n_1f`$ on $`V`$ (there is at least 1-dimensional family of such curves). Then for general $`D`$ we have $`DC0`$, so
$$(nH(\frac{n}{2}(b1)m)F)(s_0+n_1f)0,$$
and we get a contradiction:
$$0<\frac{m}{n}n_2\frac{b}{2}+\frac{1}{2}=\frac{n_1n_2+1}{2}0.$$
The lemma is proved.
By the lemma, we may assume $`m0`$. The following argumentation proving the theorem is nearly the same as in .
Obviously, $`𝒟`$ is not canonical, so it has maximal singularities. Let a curve $`B`$ be the center of a maximal singularity. Then $`B`$ is a section of $`\rho `$ not lying on the ramification divisor, and we can ”untwist” such a maximal singularity using the Bertini involution centered at $`B`$.
So we assume that the center of any maximal singularity is a point. Set $`D_1D_2=Z^h+Z^v`$, where $`Z^h`$ and $`Z^v`$ are horizontal and vertical (effective) cycles. We see that
$$Z^h+Z^v=2s_0+(\beta n^2+4mn)f.$$
Denote $`\mathrm{deg}Z^h\stackrel{\mathrm{def}}{=}Z^hF`$ and $`\mathrm{deg}Z^v\stackrel{\mathrm{def}}{=}Z^v(K_V)`$. Then, there exist a discrete valuation $`\nu `$ centered over a point $`B_0`$ in a fiber $`F_t`$ (over a point $`t^1`$), a positive number $`e=\nu (𝒟)n\delta `$, where $`\delta `$ is the canonical multiplicity with respect to $`\nu `$, such that a component $`Z_t^v`$ of $`Z^v`$ lying in $`F_t`$ has the degree
$$\mathrm{deg}Z_t^v<\{\begin{array}{cc}\beta n^2+4n\frac{e}{\nu (F_t)},\hfill & \text{if }\beta >0;\hfill \\ 4n\frac{e}{\nu (F_t)},\hfill & \text{if }\beta 0.\hfill \end{array}$$
Notice that $`m^h\stackrel{\mathrm{def}}{=}mult_{B_0}Z^h\mathrm{deg}Z^h=2n^2`$ and $`m^v\stackrel{\mathrm{def}}{=}mult_{B_0}Z_t^v\mathrm{deg}Z_t^v`$. Then, there exist positive numbers $`\mathrm{\Sigma }_0\mathrm{\Sigma }_0^{}`$ and $`\mathrm{\Sigma }_1`$ suct that the following inequality holds:
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)(m^h\mathrm{\Sigma }_0+m^v\mathrm{\Sigma }_0^{})>(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+e)^2.$$
It only remains to substitute the estimations of $`m^h`$ and $`m^v`$:
$$(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)((2+\beta n^2)\mathrm{\Sigma }_0+4ne)>(2n\mathrm{\Sigma }_0+n\mathrm{\Sigma }_1+e)^2,$$
i.e.,
$$(2\beta )n^2\mathrm{\Sigma }_0(\mathrm{\Sigma }_0+\mathrm{\Sigma }_1)+(n\mathrm{\Sigma }_1e)^2<0,$$
which is impossible, if $`\beta 2`$. Theorem 3.1 is proved.
Now we shall consider smooth Mori fibrations on del Pezzo surfaces of degree 2 for $`2a+b2`$. First, let $`2a+b=2`$, i.e., $`K_V^2=2s_0+4f`$.
###### Lemma 3.3.
For $`2a+b=2`$, the only following cases may occur:
* $`a=1`$, $`n_1=n_2=2`$;
* $`a=1`$, $`n_1=n_2=0`$;
* $`a=0`$, $`n_1=n_2=1`$;
* $`a=0`$, $`n_1=0`$, $`n_2=2`$;
* $`a=2`$, $`n_1=2`$, $`n_2=4`$;
* $`a=3`$, $`n_1=2`$, $`n_2=6`$;
* $`a=4`$, $`n_1=2`$, $`n_2=8`$;
The first three cases are non-rigid.
Proof. Since $`b0`$, the only cases 2), 3), and 4) are possible if $`a0`$.
Let $`a<0`$. Then $`b>0`$, and since $`Rt_0<0`$, any curve of the class $`s_0`$ lies on $`R`$. So $`n_1>0`$ because of irreducibility of $`R`$, and such a curve is unique on $`X`$. Further, let $`\psi :\stackrel{~}{X}X`$ be the blow-up of $`t_0`$ with $`E`$ as an exceptional divisor. $`E`$ is a ruled surface of the type $`𝐅_{𝐧_\mathrm{𝟐}𝐧_\mathrm{𝟏}}`$. By $`t_E`$ and $`l_E`$ we denote the classes of an exceptional section and a fiber of $`E`$. Suppose $`\stackrel{~}{R}`$ is the strict transform of $`R`$, then
$$\stackrel{~}{R}|_Et_E+(2a+n_2)l_E$$
is an irreducible curve on $`V`$ because of the smoothness of $`R`$. So either $`2a+n_2=0`$, or $`2a+n_2n_2n_1`$. It is easy to check that the second case is possible only if $`n_1=n_2=2`$ and $`a=1`$. Further, if $`2a+n_2=0`$, then $`n_1=2`$. Moreover, curves of the class $`t_0+n_1l`$ lie on a unique effective divisor of the class $`Mn_2L`$, and since $`R`$ is irreducible, we have $`R(t_0+n_1l)0`$, whence $`a4`$.
Now we show that the first three cases are non-rigid. Let $`a=1`$, $`n_1=n_2=0`$. We see that $`X^2\times ^1`$, and $`V`$ is a conic bundle with respect to the projection onto $`^2`$: double covers of curves of the class $`t_0`$ are conics since $`Rt_0=2`$. So $`V`$ is non-rigid. Notice that this projection is given by a linear system $`|K_VF|`$, which is free from base points.
Let $`a=1`$, $`n_1=n_2=2`$. A unique curve of the class $`t_0`$ on $`X`$ lies on the ramification divisor $`R`$, so $`s_0`$ is unique on $`V`$, too. Note that
$$𝒩_{s_0|V}=𝒪(1)𝒪(2)$$
and $`|nH|`$ for $`n2`$ gives a birational morphism contracting $`s_0`$. Then, the base set of the pencil $`|ML|`$ is exactly $`s_0`$; all elements of this pencil are smooth and isomorphic to $`𝐅_\mathrm{𝟐}`$. For a general $`S|ML|`$ the restriction $`R|_S`$ is composed from $`s_0`$ and some three-section that does not intersect $`s_0`$, so after taking the double cover and contracting the pre-image of $`s_0`$ the surface $`S`$ becomes a del Pezzo surface of degree 1. This means that the anti-flip $`V/^1V^{}/^1`$ (centered at $`s_0`$) gives us a Mori fibration on del Pezzo surfaces of degree 1 with a terminal singular point lying on the exceptional curve of $`V^{}`$. Thus, $`V/^1`$ is not rigid. Note that the pencil $`|K_VF|`$ has no fixed components.
Finally, for $`a=0`$, $`n_1=n_2=1`$ the linear system $`|2K_V|`$ gives a small contraction onto the canonical model of $`V`$, which can be realized as double covering of a non-degenerated quadratic cone in $`^4`$ branched along a quartic section. If a curve of the class $`s_0`$ is unique on $`V`$ (this means that $`t_0`$ lies on the ramification divisor), then $`s_0`$ is -2-curve of the width 2 (in the notions of ). Otherwise, there are two curves of the class $`s_0`$, which are disjoint and -2-curves of the width 1. In both the cases we obtain another structure of a smooth Mori fibration on del Pezzo surfaces of degree 2 after making a flop centered at these curves. Notice again that $`|K_VF|`$ has no fixed components. The case of two curves was studied in detail in .
Lemma 3.3 is proved.
###### Remark 3.4.
It is highly likely that cases 4) – 7) are all rigid. Unfortunately, the author can not prove it yet. As it often occur in the practice of the maximal singularity method, the problems are related to excluding of infinitely near singularities over points on some rational curves of a special kind.
Now we consider cases when $`2a+b=1`$, i.e., $`K_V^2=2s_0+6f`$.
###### Lemma 3.5.
If $`2a+b=1`$, the only following three cases may occur:
* $`a=0`$, $`n_1=0`$, $`n_2=1`$;
* $`a=1`$, $`n_1=1`$, $`n_2=2`$;
* $`a=2`$, $`n_1=1`$, $`n_2=4`$.
The first two are non-rigid.
Proof. Arguing as before, it is easy to show that $`a2`$, so the only 1) – 3) are possible.
In order to show that the case 1) is non-rigid, note that $`V`$ can be obtained by blowing up of an elliptic curve of degree 2 on a smooth Fano variety $`U`$ of genus 9 and index 2 (this is so-called the double space of index 2, i.e., the double cover of $`^3`$ branched along a quartic). Observe that the birational morphism $`VU`$ is defined by the linear system $`|2K_V2F|`$. Some results (but very incomplete) about these varieties are contained in .
Consider the case $`a=1`$, $`n_1=1`$, $`n_2=2`$. There exists a unique curve of the class $`s_0`$, and it is a -2-curve of the width 1. Let $`VV^+`$ be a flop centered at this curve. A linear system $`|H2F|`$ contains a unique element, which we denote $`G_V`$. Its strict transform $`G_V^+`$ on $`V^+`$ is a surface that is isomorphic to the double cover of $`^2`$ branched along either a smooth conic or a couple of different lines. So $`G_V^+`$ is either $`^1\times ^1`$, or a quadratic cone in $`^3`$. Then, there exists an extremal contraction $`V^+U`$ of $`G_V^+`$. U is a double cone over the Veronese surface, but with a quadratic singularity (arising from a quadratic singularity of the ramification divisor). We see also that $`U`$ has (birationally) structures of fibrations on del Pezzo surfaces of degree 1 (see section 2, the case $`n_1=0`$, $`n_2=1`$, $`n_3=2`$). Notice that $`|K_VF|`$ has no fixed components.
The lemma is proved.
As after lemma 3.3, the author can say the same thing about the case 3): it should be rigid, but I can not prove it yet.
We complete this survey of del Pezzo fibrations by the following lemma:
###### Lemma 3.6.
Cases $`2a+b0`$ can not occur.
Proof. The same reasons as above, except the case $`a=0`$. If $`a=0`$, then $`b=0`$, so $`V`$ is isomorphic to the direct product of $`^1`$ and a smooth del Pezzo surface of degree 2. But in this case $`V`$ is not a Mori fibration because of the relative Picard number (it is equal to 8). The lemma is proved. |
warning/0002/cond-mat0002062.html | ar5iv | text | # Many-body solitons in a one-dimensional condensate of hard core bosons
## Abstract
A mapping theorem leading to exact many-body dynamics of impenetrable bosons in one dimension reveals dark and gray soliton-like structures in a toroidal trap which is phase-imprinted. On long time scales revivals appear that are beyond the usual mean-field theory.
Dark and gray solitons are a generic feature of the nonlinear Schrödinger equation with repulsive interactions, and several calculations of their dynamics based on the Gross-Pitaevskii (GP) equation have appeared , as well as experiments demonstrating their existence in atomic BECs . Since the GP equation is a nonlinear approximation to the more exact linear many-body Schrödinger equation, this raises the question of how observed solitonic behavior arises in a theory which is linear at a fundamental level. Here this issue will be examined with the aid of exact many-body solutions. It has been shown by Olshanii that at sufficiently low temperatures, densities, and large positive scattering length, a BEC in a thin atom waveguide has dynamics which approach those of a one-dimensional (1D) gas of impenetrable point bosons. This is a model for which the exact many-body energy eigensolutions were found in 1960 using an exact mapping from the Hilbert space of energy eigenstates of an ideal gas of fictitious spinless fermions to that of many-body eigenstates of impenetrable, and therefore strongly interacting, bosons . The term “Bose-Einstein condensation” is used here in a generalized sense; it was shown by Lenard and by Yang and Yang that for the many-boson ground state of this system, the occupation of the lowest single-particle state is of order $`\sqrt{N}`$ where $`N`$ is the total number of atoms, in contrast to $`N`$ for usual BEC. Nevertheless, since $`N1`$ and the momentum distribution has a sharp peak in the neighborhood of zero momentum , this system shows strong coherence effects typical of BEC. The response of a BEC of this type to application of a delta-pulsed optical lattice was recently calculated by Rojo et al. , using the Fermi-Bose mapping theorem , as an exactly calculable model of dynamical optical lattice behavior. They found spatial focussing and periodic self-imaging (Talbot effect), which decay as a result of interactions. This decay is absent in the GP approximation and therefore serves as a signature of many-body interaction effects omitted in GP.
In this Letter we examine the appearence of dark soliton-like structures using the model of a 1D hard-core Bose gas in a toroidal trap, or ring, with cross section so small that motion is essentially circumferencial . The Fermi-Bose mapping is employed to generate exact solutions for this problem. We identify stationary solutions which reflect some properties of dark solitons from the GP theory when the ring is pierced at a point by an intense blue-detuned laser. We also present dynamical solutions when half of an initially homogeneous ring BEC is phase-imprinted via the light-shift potential of an applied laser, leading to gray soliton-like structures whose velocity depends on the imposed phase-shift . Such structures are apparent for times less than the echo time $`\tau _e=L/c`$, with $`L`$ the ring circumference and $`c`$ the speed of sound in the BEC. On longer time scales the dynamics becomes very complex showing Talbot recurrences which are beyond the GP theory.
Time-dependent Fermi-Bose mapping theorem: The original proof was restricted to energy eigenstates, but the generalization to the time-dependent case is almost trivial. The Schrödinger Hamiltonian is assumed to have the structure
$$\widehat{H}=\underset{j=1}{\overset{N}{}}\frac{\mathrm{}^2}{2m}\frac{^2}{x_j^2}+V(x_1,\mathrm{},x_N;t),$$
(1)
where $`x_j`$ is the one-dimensional position of the $`j\mathrm{𝑡ℎ}`$ particle and $`V`$ is symmetric (invariant) under permutations of the particles. The two-particle interaction potential is assumed to contain a hard core of 1D diameter $`a`$. This is conveniently treated as a constraint on allowed wave functions $`\psi (x_1,\mathrm{},x_N;t)`$:
$$\psi =0\text{if}|x_jx_k|<a,1j<kN,$$
(2)
rather than as an infinite contribution to $`V`$, which then consists of all other (finite) interactions and external potentials. The time-dependent version starts from fermionic solutions $`\psi _F(x_1,\mathrm{},x_N;t)`$ of the time-dependent many-body Schrödinger equation (TDMBSE) $`\widehat{H}\psi =i\mathrm{}\psi /t`$ which are antisymmetric under all particle pair exchanges $`x_jx_k`$, hence all permutations. As in the original theorem , define a “unit antisymmetric function”
$$A(x_1,\mathrm{},x_N)=\underset{1j<kN}{}\text{sgn}(x_kx_j),$$
(3)
where $`\text{sgn}(x)`$ is the algebraic sign of the coordinate difference $`x=x_kx_j`$, i.e., it is +1(-1) if $`x>0`$($`x<0`$). For given antisymmetric $`\psi _F`$, define a bosonic wave function $`\psi _B`$ by
$$\psi _B(x_1,\mathrm{},x_N;t)=A(x_1,\mathrm{},x_N)\psi _F(x_1,\mathrm{},x_N;t)$$
(4)
which defines the Fermi-Bose mapping. $`\psi _B`$ satisfies the hard core constraint (2) if $`\psi _F`$ does, is totally symmetric (bosonic) under permutations, obeys the same boundary conditions as $`\psi _F`$, e.g. periodic boundary conditions on a ring, and $`\widehat{H}\psi _B=i\mathrm{}\psi _B/t`$ follows from $`\widehat{H}\psi _F=i\mathrm{}\psi _F/t`$ .
Exact solutions for impenetrable point bosons: The mapping theorem leads to explicit expressions for all many-body energy eigenstates and eigenvalues of a 1D scalar condensate (bosons all of the same spin) under the assumption that the only two-particle interaction is a zero-range hard core repulsion, represented by the $`a0`$ limit of the hard-core constraint. Such solutions were obtained in Sec. 3 of the original work for periodic boundary conditions and no external potential. The exact many body ground state was found to be a pair product of Bijl-Jastrow form: $`\psi _0=\text{const.}_{i>j}|\mathrm{sin}[\pi L^1(x_ix_j)]|`$. In spite of the very long range of the individual pair correlation factors $`|\mathrm{sin}[\pi L^1(x_ix_j)]|`$, the pair distribution function $`D(x_{ij})`$, the integral of $`|\psi _0|^2`$ over all but two coordinates, was found to be of short range: $`D(x_{ij})=1j_0^2(\pi \rho x_{ij})`$, with $`j_0(\xi )=\mathrm{sin}\xi /\xi `$, the spherical Bessel function of order zero. The system was found to support propagation of sound with speed $`c=\pi \mathrm{}\rho /m`$ where $`\rho =N/L`$, the 1D atom number density.
To generalize to the time-dependent case, assuming that the many-body potential of Eq. (1) is a sum of one-body external potentials $`V(x_j,t)`$, one generalizes the time-independent determinantal many-fermion wavefunction to a determinant
$$\psi _F(x_1,\mathrm{},x_N;t)=C\underset{i,j=1}{\overset{N}{det}}\varphi _i(x_j,t),$$
(5)
of solutions $`\varphi _i(x,t)`$ of the one-body TDSE in the external potential $`V(x,t)`$. It then follows that $`\psi _F`$ satisfies the TDMBSE, and it satisfies the impenetrability constraint (vanishing when any $`x_j=x_{\mathrm{}}`$) trivially due to antisymmetry. Then by the mapping theorem $`\psi _B`$ of Eq.(4) satisfies the same TDMBSE.
Dark solitons on a ring: Consider $`N`$ bosons in a tight toroidal trap, and denote their 1D positions measured around the circumference by $`x_j`$. This is equivalent to the exactly-solved model of $`N`$ impenetrable point bosons in 1D with wave functions satisfying periodic boundary conditions with period $`L`$ equal to the torus circumference, and the fundamental periodicity cell may be chosen as $`L/2<x_j<L/2`$. However, the rotationally invariant quantum states of this problem do not reveal any dark soliton-like structures. To proceed we therefore consider the case that a blue-detuned laser field pierces the ring at $`x=0`$ by virtue of the associated repulsive dipole force: The light sheet then provides a reference position for the null of the dark soliton. Assume that the light sheet is so intense and narrow that it may be replaced by a constraint that the many-body wave function (hence the orbitals $`\varphi _i`$) must vanish whenever any $`x_j=0`$. Then the appropriate orbitals $`\varphi _i(x)`$ are free-particle energy eigenstates vanishing at $`x=0`$ and periodic with period $`L`$. The complete orthonormal set of even-parity eigenstates $`\varphi _n^{(+)}`$ and odd-parity eigenstates $`\varphi _n^{()}`$ are
$`\varphi _n^{(+)}(x)`$ $`=`$ $`\sqrt{2/L}\mathrm{sin}[(2n1)\pi |x|/L],`$ (6)
$`\varphi _n^{()}(x)`$ $`=`$ $`\sqrt{2/L}\mathrm{sin}(2n\pi x/L),`$ (7)
with $`n`$ running from $`1`$ to $`\mathrm{}`$. The odd eigenstates are the same as those of free particles with no $`x=0`$ constraint, since these already vanish at $`x=0`$. However, the even ones are strongly affected by the constraint, their cusp at $`x=0`$ being a result of the impenetrable light sheet at that point. If one bends a 1D box $`L/2<x<L/2`$ with impenetrable walls into a ring, identifying the walls at $`\pm L/2`$, then those particle-in-a-box eigenfunctions which are even about the box center become identical with the $`\varphi _n^{(+)}`$, and their cusp results from the nonzero slope of these functions at the walls. The $`N`$-fermion ground state is obtained by inserting the lowest $`N`$ orbitals (6) into the determinant (5) (filled Fermi sea). Assume that $`N`$ is odd. Since $`\varphi _1^{(+)}`$ is lower than $`\varphi _1^{()}`$, this Fermi sea consists of the first $`(N+1)/2`$ of the $`\varphi _n^{(+)}`$ and the first $`(N1)/2`$ of the $`\varphi _n^{()}`$. The $`N`$-boson ground state is then given by (4). Since $`A^2=1`$, its one-particle density $`\rho (x)`$ is the same as that of the $`N`$-fermion ground state, the sum of partial densities contributed by all one-particle states in the Fermi sea. Thus it is the sum of
$$\rho ^{(+)}(x)=\frac{N+1}{2L}\frac{\mathrm{sin}[2(N+1)\pi x/L]}{2L\mathrm{sin}(2\pi x/L)},$$
(8)
and
$$\rho ^{()}(x)=\frac{N1}{2L}\frac{\mathrm{sin}[(N1)\pi x/L]\mathrm{cos}[(N3)\pi x/L]}{L\mathrm{sin}(2\pi x/L)}$$
(9)
In the thermodynamic limit $`N\mathrm{}`$, $`L\mathrm{}`$, $`N/L\rho `$ for fixed $`x`$, $`\rho ^{(\pm )}`$ each contribute half of the total density $`\rho (x)`$:
$$\rho (x)\rho [1j_0(2\pi \rho x)].$$
(10)
Since $`j_0(0)=1`$, $`\rho (x)`$ vanishes at $`x=0`$ and approaches the mean density $`\rho `$ over a healing length $`L_h=1/2\rho `$ with damped spatial oscillations about its limiting value. This differs in detail from the density $`\rho _{\mathrm{}}\mathrm{tanh}^2(x/w)`$ of a GP dark soliton , with $`\rho _{\mathrm{}}`$ the background density and $`w`$ the corresponding healing length, but has some qualitative similarity. However, it is only the odd component $`\rho ^{()}(x)\rho (x)/2`$ which has the feature of a dark soliton that the corresponding odd orbitals have a $`\pi `$ phase-jump at $`x=0`$ (and also at $`x=\pm L/2`$ to obey the periodic boundary conditions). But the odd and even components can never be separated physically, so the odd dark soliton-like component is always accompanied by the even non-soliton component.
Next, suppose that the light-sheet is turned off at $`t=0`$ by removing the constraint that the wave function vanish at $`x=0`$. The solution of the TDMBSB for the many-boson system is then given by (4) where the Slater determinant (5) is built from the first $`(N+1)/2`$ of the $`\varphi _n^{(+)}(x,t)`$ and the first $`(N1)/2`$ of the $`\varphi _n^{()}(x,t)`$, where these time-dependent orbitals are solutions of the single-free-particle TDSE which (a) reduce to the orbitals (6) at $`t=0`$, and (b) satisfy periodic boundary conditions with periodicity cell $`L/2<x<L/2`$. The odd solutions are trivial: Since these never “see” the $`x=0`$ constraint even for $`t<0`$, they differ from the odd orbitals (6) only by time-dependent phase shifts: $`\varphi _n^{()}(x,t)=\varphi _n^{()}(x)e^{i\omega _nt}`$ with $`\omega _n=\mathrm{}k_n^2/2m`$ and $`k_n=2n\pi /L`$. It follows that $`\rho ^{()}(x,t)`$ is time-independent, and given in the thermodynamic limit by
$$\rho ^{()}(x,t)(\rho /2)[1j_0(2\pi \rho x)].$$
(11)
This further reinforces our view that the odd component of the density shares features of a dark soliton. The even-parity orbitals $`\varphi _n^{(+)}(x,t)`$ are complicated since the removal of the light sheet constitutes a large, sudden perturbation. Indeed, the periodic even-parity solutions of the free-particle Schrödinger equation are $`\chi _p^{(+)}(x)=\sqrt{(2\delta _{p0})/L}\mathrm{cos}(2p\pi x/L)`$ with $`p=0,1,2,\mathrm{}`$, and these are very different from the solutions $`\varphi _n^{(+)}(x)`$ with the $`x=0`$ constraint \[Eq. (6)\]. Nevertheless, since the $`\chi _p^{(+)}(x)`$ are complete for the subspace of even-parity, spatially periodic functions, one can expand the $`\varphi _n^{(+)}(x,t)`$ in terms of the $`\chi _p^{(+)}(x)`$, which evolve with time-dependent phases $`e^{i\omega _pt}`$ with $`\omega _p=\mathrm{}k_p^2/2m`$ and $`k_p=2p\pi /L`$. One finds
$$\varphi _n^{(+)}(x,t)=\frac{2(2n1)}{\pi }\sqrt{\frac{2}{L}}\underset{p=0}{\overset{\mathrm{}}{}}\frac{(2\delta _{p0})\mathrm{cos}(k_px)e^{i\omega _pt}}{(2n1)^24p^2}$$
(12)
$`\rho ^{(+)}(x,t)`$ is the sum of absolute squares of the first $`(N+1)/2`$ of the sums (11), generalizing (7). Adding the time-independent expression $`\rho ^{()}(x,t)`$, given in the thermodynamic limit by (10) or exactly by (8), one finds the time-dependent total density $`\rho (x,t)`$. There are two important time scales: One is the Poincaré recurrence time $`\tau _r`$. Noting that $`\omega _p`$ in (11) is proportional to $`p^2`$, one finds that all terms in the sum are time-periodic with period $`\tau _r=mL^2/\pi \mathrm{}`$, which is therefore the recurrence time for the density and in fact all properties of our model . The other important time is the echo time $`\tau _e`$, the time for sound to make one circuit around the torus. Recalling that the speed of sound in this system is $`c=\pi \mathrm{}\rho /m`$ , one finds $`\tau _e=\tau _r/N`$. For $`t<<\tau _e`$ after the constraint is removed, the initial density develops sound waves that propagate around the ring, and that we examine below in the context of phase-imprinting. For $`t>\tau _e`$ the evolution is very complex, but complete recurrences occur for times $`t=n\tau _r`$ with fractional revivals in between.
Gray soliton formation by phase-imprinting: Consider next a toroidal BEC in its ground state to which a phase-imprinting laser is applied over half the ring at $`t=0`$. This is described by the single-particle Hamiltonian
$$\widehat{H}=\underset{j=1}{\overset{N}{}}\left[\frac{\mathrm{}^2}{2m}\frac{^2}{x_j^2}\mathrm{}\mathrm{\Delta }\theta \delta (t)S(x_j)\right]$$
(13)
where $`S(x)=\theta (L/4|x|)`$, i.e., unity for $`L/4<x<L/4`$ and zero elsewhere. This is the technique used in recent experiments , here idealized to a delta-function in time and to sharp spatial edges. Before the pulse the most convenient free-particle orbitals in (5) are plane waves $`\varphi _n(x)=\sqrt{(1/L)}e^{ik_nx}`$ where $`k_n=2n\pi /L`$ and $`n=n_F,n_F+1,\mathrm{},n_F1,n_F`$ with $`n_F=(N1)/2`$. Let $`\varphi _n(x,t)`$ be the solution of the TDSE with the Hamiltonian (12) reducing to the above $`\varphi _n(x)`$ just before the pulse. Then the solutions just after the pulse are $`\varphi _n(x,0+)=\varphi _n(x)e^{iS(x)\mathrm{\Delta }\theta }`$. The potential gradients at the pulse edges impart momentum kicks to the particles there which induce both compressional waves propagating at the speed, $`c`$, of sound and density dips (gray solitons) moving at speeds $`|v|<c`$. The expansion of $`\varphi _n(x,t)`$ in terms of the unperturbed plane waves is evaluated as
$`\varphi _n(x,t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+e^{i\mathrm{\Delta }\theta }\right){\displaystyle \frac{1e^{i\mathrm{\Delta }\theta }}{\pi }}{\displaystyle \underset{\mathrm{}=\mathrm{}}{\overset{\mathrm{}}{}}}`$ (14)
$`\times `$ $`{\displaystyle \frac{(1)^{\mathrm{}}\varphi _{n2\mathrm{}1}(x)e^{i\omega _{n2\mathrm{}1}t}}{2\mathrm{}+1}}`$ (15)
and the time-dependent density is the sum of the absolute squares of the lowest $`N`$ of these. Figure 1 shows numerical simulations obtained using Eq. (13) for $`N=51`$, $`t/\tau _e=0.051`$, and $`\mathrm{\Delta }\theta =\pi `$ (solid line), and $`\mathrm{\Delta }\theta =0.5\pi `$ (dashed line): due to symmetry we show only half of the ring $`L/2<x<0`$, the phase-shift being imposed at $`x=L/4`$. Considering times short compared to the echo time means that the corresponding results are not very sensitive to the periodic boundary conditions, and also therefore apply to a linear geometry. The initial density profile is flat with a value $`\rho _0L=51`$. For both phase-shifts two distinct maxima are seen, which travel at close to the speed of sound $`c`$, and two distinct minima, which are analogous to gray solitons and travel at velocities $`|v|/c<1`$.
In addition, there are also high wavevector oscillations which radiate at velocities greater than $`c`$. In the case of a phase-shift $`\mathrm{\Delta }\theta =\pi `$, the density is symmetric about $`x=L/4`$, whereas for a phase-shift other than a multiple of $`\pi `$ the evolution is not symmetric, see the dashed line where the global minimum moves to the right in reponse to the phase-shift.
In Fig. 2 we plot the calculated velocity of the global density minimum relative to the speed of sound for a variety of phase-shifts $`\mathrm{\Delta }\theta `$. The basic trend is that larger phase-shift means lower velocity, in qualitative agreement with recent experiments , but there is a sharp velocity peak at $`\mathrm{\Delta }\theta 0.83\pi `$: This peak results from the cross-over between two local minima in the density. These general features, the generation of gray solitons and density waves, agree with those of the GP theory, but here arise out of the exact many-body calculation.
In conclusion, using our exactly-soluble 1D model we hope to have shown that the dark solitonic features of atomic BECs normally described within the mean-field GP theory arise naturally from consideration of the exact linear many-body theory for times less than the echo time. An advantage of this approach is that it is number-conserving and does not rely on any symmetry-breaking approximation. In addition, long time dynamics such as collapses and revivals are accounted for . A detailed comparison between our results and current experiments is not possible as they do not conform to the conditions for a 1D system. However, some estimates are in order to set the appropriate time scales: If we consider <sup>87</sup>Rb with a ring of circumference $`L=100`$ $`\mu `$m, and a high transverse trapping frequency $`\omega _{}=2\pi \times 10^5`$ Hz, then we are limited to atom numbers $`N<300`$ , so these are small condensates. We then obtain $`\tau _r=4.6`$ s, and $`\tau _e=90`$ ms for $`N=51`$. Finally, we remark that since our approach relied on the mapping between the strongly-interacting Bose system and a non-interacting “spinless Fermi gas” model, this suggests that dark and gray solitons should also manifest themselves in the density for the 1D Fermi system. Although real fermions have spin, the interactions used here to generate solitons were spin-independent.
This work was supported by the Office of Naval Research Contract No. N00014-99-1-0806. |
warning/0002/math0002029.html | ar5iv | text | # On the Weyl tensor of a self-dual complex 4-manifold
## 1. Introduction
Twistor theory, created by Penrose , establishes a close relationship between conformal Riemannian geometry in dimension 4, and (almost) complex geometry in dimension 3. In particular, to a Riemannian manifold $`M`$ for which the part $`W^{}`$ of the Weyl tensor vanishes identically (self-dual), one associates its twistor space $`Z`$, a complex 3-manifold containing rational curves with normal bundle $`𝒪(1)𝒪(1)`$, and admitting a real structure with no fixed points , , .
The space of such curves is a complex 4-manifold $`M^{}`$ with a holomorphic conformal structure and is, therefore, a conformal complexification of $`M`$ , , .
As the conformal geometry of $`M`$ is encoded by the complex geometry of $`Z`$, we ask ourselves what holomorphic object on $`Z`$ corresponds to $`W^+`$, the Weyl tensor of the self-dual manifold $`M`$. It seems that this question, although natural, has not been considered in the literature, and maybe a reason for that is that the answer appears to be a highly non-linear object.
This object is more easily understood in the framework of complex-Riemannian geometry (see Section 2): following LeBrun , we (locally) introduce the space $`B`$ of complex null-geodesics of M (ambitwistor space). For a self-dual (complex) 4-manifold M, its (local) twistor space is then defined as the 3-manifold of $`\beta `$-surfaces (some totally geodesic isotropic surfaces, see Section 2).
The ambitwistor space $`B`$ and (in the self-dual case) the twistor space $`Z`$ completely describe the conformal structure of M. In particular, a null-geodesic $`\gamma `$ in M corresponds to the set of rational curves in $`Z`$ tangent to a 2-plane. The union of these curves, called the integral $`\alpha `$-cone of $`\gamma `$ (see Section 3), is lifted to a (linearized) $`\alpha `$-cone in $`T_\gamma B`$. Our first result (Theorem 1) is that the Weyl tensor of M is equivalent to the projective curvature (see Section 4) of the field of $`\alpha `$-cones on $`B`$. In particular, if such a cone is flat, then $`W^+`$ vanishes on certain isotropic planes in M.
We use Theorem 1 to investigate global properties of a self-dual manifold M: If the integral $`\alpha `$-cone of $`\gamma `$ is part of a smooth surface in $`Z`$, then the linearized $`\alpha `$-cone is flat (Theorems 2, 2). In particular, the space $`𝐌_0`$ of rational curves of $`Z`$ with normal bundle $`𝒪(1)𝒪(1)`$ is compact iff $`Z^3`$. On the other hand, it is known, from a theorem of Campana , that, for a compact twistor space $`Z`$, $`𝐌_0`$ can be compactified within the space of analytic cycles iff $`Z`$ is Moishezon. It appears then that the conformal structure does not extend smoothly to the compactification.
A good illustration of what happens in the non-flat (self-dual) case is the Kähler-Einstein manifold $`^2`$ whose twistor space is known to be the manifold of flags in $`^3`$ , see Section 8.
Different methods allow us to generalize Theorem 3 to non geodesically-connected self-dual manifolds: We show (Theorem 3) that if a self-dual manifold admits a compact, simply-connected, null-geodesic, then it is conformally flat. We also note that the rational curves in $`Z`$, corresponding to the points of M (see Section 2) are then geodesics of some projective structure of $`Z`$ iff the latter is projectively flat (Corollary 1).
The isotropic, totally geodesic surfaces (called $`\beta `$-surfaces) in a self-dual manifold M appear to have a projective structure, given by the null-geodesics of M contained in it (Section 6). We show that it is flat (i.e. locally equivalent to $`^2`$) (Corollary 3), and we obtain a classification of the compact $`\beta `$-surfaces of a self-dual 4-manifold (Theorem 4).
Theorem 3 can be adapted for conformal 3-manifolds : A conformal 3-manifold admitting a rational curve as a null-geodesic is conformally flat (Theorem 7). The conformal geometry in dimensions 3 and 4 are related, as any geodesically convex 3-manifold $`Q`$ can be realized as the conformal infinity of a self-dual 4-manifold M . In particular, $`Q`$ is umbilic in M, and we relate, in Section 7, the conformal invariants of the 2 manifolds: the Cotton-York tensor of $`Q`$ is identified to the derivative, in the normal direction, of the Weyl tensor of M (Theorem 5). This result can be equally stated in the real framework.
The paper is organized as follows: in Section 2 we recall the classical results of the twistor theory (especially for complex 4-manifolds), in Section 3 we introduce the $`\alpha `$-cones on the (ambi-)twistor space, and, in Section 4, we prove the equivalence between the projective curvature of the latter and the Weyl tensor $`W^+`$ of M. Section 5 is devoted to the proof of some results of the type “compactness implies conformal (projective) flatness”: Theorems 2, 2 and 3, mentioned above. We study the projective structure of $`\beta `$-surfaces in Section 7, and we illustrate the above results on the special case of the self-dual manifold $`^2`$, in Section 8.
Acknowledgements The author is deeply indebted to Paul Gauduchon, for his care in reading the manuscript and for his constant help during the research and redaction.
## 2. Preliminaries
The content of this paper makes use of complex-Riemannian geometry (with the exception of Theorem 5 and Corollary 4, which hold also in the real framework). Complex-Riemannian geometry is obtained from Riemannian geometry by replacing the field $``$ by $``$ (e.g. a complex “metric” becomes a non-degenerate symmetric complex-bilinear form on the tangent space), and all classical results hold, naturally with the exception of those making use of partitions of unity. We will often omit the prefix “complex-”, when referring to geometric objects, and we will always consider them, unless otherwise stated, in the framework of complex-Riemannian geometry.
### 2.1. Conformal complex 4-manifolds
Let M be a 4-dimensional complex manifold. A conformal structure is defined, as in the real case , by a everywhere non-degenerate section $`c`$ of the complex bundle $`S^2(T^{}𝐌)L^2`$, where $`L`$ is a given line bundle of scalars of weight 1, and $`L^4\kappa ^1`$, the anti-canonical bundle of M. (While on an oriented real manifold such a line bundle always exists, being topologically trivial, in the complex case the existence of $`L^2`$, a square root of the anti-canonical bundle, is submitted to some topological restrictions.) From now on, only holomorphic conformal structures will be considered, thus $`L`$ is a holomorphic bundle and $`c`$ a holomorphic section of $`S^2(T^{}𝐌)L^2`$. (In fact, all we need to define the conformal structure $`c`$ on the 4-manifold M is just the holomorphic bundle $`L^2`$; in odd dimensions the situation is different, see Section 7.)
As in the real case, $`c`$ is locally represented by symmetric bilinear forms on $`T𝐌`$, but global representative metrics do not exist, in general.
For each point $`x𝐌`$, there is an isotropy cone $`C_x`$ in the tangent space $`T_x𝐌`$, who uniquely determines the conformal structure $`c`$. In the associated projective space, $`(T_x𝐌)^3`$, the cone $`C_x`$ projects onto the non-degenerate quadratic surface $`(C_x)`$, which is actually a ruled surface isomorphic to $`^1\times ^1`$. We thus get 2 families of complex projective lines contained in $`(C)`$, that is, 2 families of isotropic 2-planes in $`CT𝐌`$, respectively called $`\alpha `$-planes, the other $`\beta `$-planes. This choice corresponds to the choice of an “orientation” of M. On a real 4-manifold an orientation is chosen by picking a class of volume forms (which is not possible in this complex framework) or by choosing one of the two possible Hodge operators compatible with the conformal structure $`:\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^2𝐌`$ (which can also be done in our complex case, ). As $``$ is a symmetric involution, $`\mathrm{\Lambda }^2𝐌`$ decomposes in $`\mathrm{\Lambda }^+𝐌\mathrm{\Lambda }^{}𝐌`$ consisting in $`\pm 1`$-eigenvectors of $``$, respectively called self-dual and anti-self-dual 2-forms; the isotropic vectors in $`\mathrm{\Lambda }^+𝐌`$ and $`\mathrm{\Lambda }^{}𝐌`$ are then exactly the decomposable elements $`uv\mathrm{\Lambda }^\pm 𝐌`$, with $`u,v𝐌`$.
###### Definition 1.
An $`\alpha `$-plane $`F^\alpha `$ (resp. a $`\beta `$-plane $`F^\beta `$ in $`T𝐌`$ is a 2-plane such that $`\mathrm{\Lambda }^2F^\alpha `$ (resp. $`\mathrm{\Lambda }^2F^\beta `$) is a is a self-dual (resp. anti-self-dual) isotropic line in $`\mathrm{\Lambda }^2𝐌`$.
Remark.
The $`\alpha `$\- and $`\beta `$-planes can be interpreted in terms of spinors. The structure group of the tangent bundle $`T𝐌`$ is restricted to the conformal orthogonal complex group, $`CO(4,):=(O(4,)\times ^{})/\{\pm \mathrm{𝟏}\}`$, where $`O(4,):=\{AGL(4,)|A^tA=\text{Id}\}`$, by the given conformal structure of M. The choice of an orientation is the further restriction of this group to the connected component of the identity, $`CO_0(4,):=SO(4,)\times ^{}`$, where $`SO(4,):=O(4,)SL(4,)`$. Consider a local metric $`g`$ in the conformal class $`c`$. We have then locally defined $`Spin`$ structures, and associated $`Spin`$ bundles $`V_+,V_{}`$, as in the real case ,. They are rank 2 complex vector bundles, and for each local section of $`L`$ (i.e. a metric in $`c`$), each of them is equipped with a (complex) symplectic structure $`\omega _+\mathrm{\Lambda }^2V_+,\omega _{}\mathrm{\Lambda }^2V_{}`$, respectively. Then we locally have $`T𝐌V_+V_{}`$, and $`g=\omega _+\omega _{}`$, for the fixed metric $`gc`$. $`\alpha `$-(resp. $`\beta `$-) planes are then nothing but the isotropic 2-planes obtained by fixing the first (resp. the second) factor in the above tensor product:
###### Proposition 1.
An $`\alpha `$-plane, resp. $`\beta `$-plane $`FT_x𝐌`$ is a complex plane $`\psi _+V_{}`$, resp $`V_+\psi _{}`$, where $`\psi _+V_+\{0\}`$, resp. $`\psi _{}V_{}\{0\}`$.
The $`\alpha `$-planes in $`T_x𝐌`$ are, thus indexed by $`(V_+)_x`$, and $`\beta `$-planes by $`(V_{})_x`$, and these projective bundles are globally well-defined on M, .
Remark.
It is obvious that a change of orientation interchanges the $`\alpha `$ and $`\beta `$-planes; the same is true for self-duality and anti-self-duality, to be defined below.
For a local metric $`g`$ in $`c`$, we denote by $`R^g`$ its Riemannian curvature, and by $`W`$ the Weyl tensor, i.e. the trace-free component of $`R^g`$, which is known to be independent of the chosen metric within the conformal class . It splits into two components $`W^+,W^{}`$, and the easiest way to see that is the spinorial decomposition of the space of the curvature tensors $`\mathrm{\Lambda }^2\mathrm{\Lambda }^2`$, ,,, obtained from the relation $`T𝐌=V_+V_{}`$ and some of the Clebsch-Gordan identities .
$$=𝒮𝒲^+𝒲^{},$$
where $`𝒮`$ is the complex line of scalar curvature tensors, included in $`\mathrm{\Lambda }^2V_+\mathrm{\Lambda }^2V_{}`$, $`=S^2V_+S^2V_{}`$ is the space of trace-free Ricci tensors, and $`𝒲^+=S^4V_+`$, $`𝒲^{}=S^4V_{}`$ are the spaces of self-dual, resp. anti-self-dual Weyl tensors (where $`S^pV_\pm `$ denotes the $`p`$-symmetric power of $`V_\pm `$).
The curvature $`R^g`$ restricted to any $`\alpha `$-plane $`F`$ yields a weighted bilinear symmetric form $`R^F`$ on $`\mathrm{\Lambda }^2F`$, i.e. a section in $`L^2(\mathrm{\Lambda }^2F\mathrm{\Lambda }^2F)^{}`$:
$$(g,XY)g(R^g(X,Y)X,Y).$$
###### Proposition 2.
The (weighted) bilinear form $`R^F`$ depends only on the self-dual Weyl tensor, and this one is completely determined by the (weighted) values of $`R^F`$ for all $`\alpha `$-planes $`F`$.
We have the same result for $`\beta `$-planes.
###### Proof.
Let $`F=V_+\psi _{}`$ be an $`\alpha `$-plane, and let $`X=\psi _+\phi _1,Y=\psi _+\phi _2F`$, and suppose, for simplicity, that $`\omega _{}(\phi _1,\phi _2)=1`$, so $`XY\mathrm{\Lambda }^2F`$ is identified to the element $`\psi _+\psi _+S^2V_+`$. Then it is easy to see that $`R^F`$, evaluated on $`XY`$, is nothing but the evaluation of $`RS^2(\mathrm{\Lambda }^2𝐌)`$ on $`(XY)(XY)\psi _+\psi _+\psi _+\psi _+S^4V_+`$, which depends only on the positive (or self-dual) part of the Weyl tensor. To prove the second assertion, we remark that $`W^+`$, being a quadrilinear symmetric form on $`V_+`$, can be identified with a polynomial of degree 4 on $`V_+`$, which is determined by its values. ∎
###### Definition 2.
A conformal structure $`c`$ on a 4-manifold M is called self-dual (resp. anti-self-dual) iff $`W^{}=0`$ (resp. $`W^+=0`$).
Remark.
In general, geodesics on a conformal manifold depend on the chosen metric, with the exception of the isotropic ones (or null-geodesics). Therefore the existence of totally geodesic surfaces tangent to $`\alpha `$\- (resp. $`\beta `$-) planes is a property of the conformal structure alone.
### 2.2. Twistor spaces
###### Definition 3.
An $`\alpha `$-surface (resp. $`\beta `$-surface) $`\alpha 𝐌`$ is a maximal, totally geodesic surface in M, whose tangent space in any point is an $`\alpha `$-plane (resp. $`\beta `$-plane).
On the other hand, any totally geodesic, isotropic surface in M is included in an $`\alpha `$\- or in a $`\beta `$-surface.
###### Definition 4.
, If, in any point $`x𝐌`$, and for any $`\alpha `$\- (resp. $`\beta `$-) plane $`FT_x𝐌`$, there is a $`\alpha `$\- (resp. $`\beta `$-) surface tangent to $`F`$ at $`x`$, we say that the family of $`\alpha `$\- (resp. $`\beta `$-) planes is integrable.
Theorem
. , The family of $`\alpha `$\- (resp. $`\beta `$-) planes of a conformal 4-manifold $`(𝐌,c)`$ integrable if and only if the conformal structure $`c`$ is anti-self-dual (resp. self-dual).
The integrability of $`\alpha `$-planes is equivalent to the integrability (in the sense of Frobenius) of a distribution $`H^\alpha `$ of 2-planes on the total space of the projective bundle $`(V_+)`$. More precisely, let $`g`$ be a local metric in the conformal class $`c`$, and let $``$ be its Levi-Civita connection. $``$ induces a connection in the bundle $`(V_+)`$, thus a horizontal distribution $`H`$, isomorphic to $`T𝐌`$ via the bundle projection. Let $`H^\alpha `$ be the 2-dimensional subspace of $`H_F`$ — where $`F(V_+)`$ is an $`\alpha `$-plane in $`T_x𝐌`$ — which projects onto $`FT_x𝐌`$. It can be easily shown (as in , see also ) that the “tautological” 2-plane distribution $`H^\alpha `$ is independent of the metric $`g`$. Then $`\alpha `$-surfaces are canonically lifted as integrable manifolds of the distribution $`H^\alpha `$. For a geodesically convex open set of M, one can prove (see ) that the space of these integrable leaves is a complex 3-manifold. (This point of view is closely related to the one of , about the integrability of the canonical almost complex structure of the real twistor space.)
The same remark can be made about $`\beta `$-surfaces.
Remark.
The existence, for any point $`x𝐌`$, of an $`\alpha `$-surface containing $`x`$ does not imply, in general, the integrability of the family of $`\alpha `$-planes : in the conformal self-dual (but not anti-self-dual) manifold $`𝐌=^2\times (^2)^{}`$ (the complexification of $`^2`$, ), the surfaces $`(\{x\}\times (^2)^{})𝐌`$ and $`(^2\times \{y\})𝐌`$ are all $`\alpha `$-surfaces, see Section 8.
Remark.
In the real framework, the twistor space of a real Riemannian 4-manifold $`M^{}`$ is the total space $`Z^{}`$ of the $`S^2`$-bundle of almost-complex structures on $`TM^{}`$, compatible with the conformal structure and the (opposite) orientation; it admits a natural almost-complex structure $`𝒥`$, equal, in $`JZ^{}`$, to the complex structure of the fibers on the vertical space $`T_J^{}Z^{}`$, and to $`J`$ itself on the horizontal space (induced by the Levi-Civita connection). Such a complex structure $`J`$ is equivalent to an isotropic complex 2-plane in $`TM`$, thus to an $`\alpha `$\- or $`\beta `$-surface (depending on the conventions), which becomes then the space of vectors of type $`(1,0)`$ for $`J`$; as the integrability of the almost-complex structure $`𝒥`$ can be expressed as the Frobenius condition applied to $`T^{(1,0)}Z^{}`$, it is equivalent to the integrability of the family of $`\alpha `$-, resp. $`\beta `$-planes.
The Penrose construction associates to an (anti-)self-dual manifold M the space $`Z`$ of $`\alpha `$\- (resp. $`\beta `$-)surfaces of M; we have seen above that $`Z`$ admits complex-analytic maps, but it may be non-Hausdorff. This is why we need to introduce the following condition , see also :
###### Definition 5.
An (anti-) self-dual manifold M is called civilized iff the space $`Z^\alpha `$ (resp. $`Z^\beta `$) of integral leaves of the distribution $`H^\alpha `$ (resp. $`H^\beta `$) in $`(V_+)`$ (resp. $`(V_{})`$) is a complex 3-manifold, and the projection $`p^+:(V_+)Z^\alpha `$ (resp. $`p^{}:(V_{})Z^\beta `$) is a submersion.
In this case, the manifold $`Z^\alpha `$ (resp. $`Z^\beta `$) — which is the space of $`\alpha `$-surfaces (resp. $`\beta `$-surface) of M — is called the $`\alpha `$\- (resp. the $`\beta `$-)twistor space of M.
From now on, we suppose that $`(𝐌,c)`$ is a self-dual complex analytic 4-manifold. As any point $`x𝐌`$ has a geodesically convex neighborhood $`U`$ (which is, therefore, civilized), we can construct $`Z^U`$, the $`\beta `$-twistor space (for short, twistor space) of $`U`$. For the infinitesimal results of this paper (from Sections 3,4,6 and 7), we will suppose (with no loss of generality) that M is civilized (for example, by replacing M by $`U`$).
We recall now the correspondence between differential geometric objects on M and complex analytic objects on its twistor space, $`Z`$, ,, see also ,,.
$`\beta `$-surfaces $`\beta 𝐌`$ correspond to points $`\overline{\beta }Z`$, by definition, and the set of $`\beta `$-surfaces passing through a point $`x𝐌`$ is a complex projective line $`Z_x`$, with normal bundle isomorphic (non-canonically) to $`𝒪(1)𝒪(1)`$ (where $`𝒪(1)`$ is the dual of the tautological bundle $`𝒪(1)`$ on $`^1`$) ,, see also .
In fact, this family of complex projective lines in $`Z`$ permits us to recover M and its conformal structure, at least locally, by the reverse Penrose construction: The normal bundle $`N_x`$ of a line $`Z_x`$ in $`Z`$ has the property $`H^1(N_x,𝒪)=0`$, thus, by a theorem of Kodaira , the space $`𝐌_0`$ of projective lines in $`Z`$ having the above normal bundle is a smooth complex manifold, whose tangent space at a point $`xZ_xZ`$ is canonically isomorphic to the space of global sections of the normal bundle $`N_x`$ of $`Z_x`$ (thus $`𝐌_0`$ has dimension 4). The conformal structure of $`𝐌_0`$ is described by its tangent cone, which corresponds to the sections of $`N_x`$ having at least one zero (as such a section decomposes as 2 sections of $`𝒪(1)`$, the vanishing condition means that they both vanish at the same point, which is a quadratic condition on the sections of $`N_x`$). We thus get a conformal diffeomorphism from M to an open set of $`𝐌_0`$.
### 2.3. Ambitwistor spaces
We remark that $`(V_{})`$ is an open set of the projective tangent bundle of $`Z`$, as $`Z`$ is the space of leaves of $`(V^{})`$, but it is important to note that, in general, the reverse inclusion is not true (i.e. not any direction in $`Z`$ is tangent to a line corresponding to a point in M, or, equivalently, $`\beta `$-surfaces are not compact $`^2`$’s, in general, see Section 5).
For example, if $`𝐌=^2\times _{}^{2}{}_{}{}^{}`$ (with the notations in Section 8), $`(V_{})`$ is an open subset in the $`^2`$-bundle $`(TZ)Z`$, consisting in the set of directions transverse to the contact structure of $`Z`$ (see Section 8.4). $`(V_{})`$ is, thus, in this case, a rank 2 affine bundle over $`Z`$.
Another canonical $`^2`$-bundle on $`Z`$, that is $`(T^{}Z)Z`$, leads to the ambitwistor space $`B`$, by definition the space of null-geodesics of M . It is an open set of the projective cotangent bundle of $`Z`$ (or, equivalently, the Grassmannian of 2-planes in $`TZ`$) (more precisely, a plane $`FT_{\overline{\beta }}Z`$ corresponds to a null-geodesic $`\gamma 𝐌`$ (contained in $`\beta `$) if it is tangent to at least one projective line $`Z_x`$, corresponding to a point $`x𝐌`$).
To see that, let $`x`$ be a point in $`𝐌`$, $`\beta `$ a $`\beta `$-surface passing through $`x`$, i.e. $`\overline{\beta }Z`$ and $`Z_x`$ contains $`\overline{\beta }`$; let $`FT_{\overline{\beta }}Z`$ be a plane tangent to $`Z_x`$. As small deformations of $`Z_x`$ still correspond to points of M, we consider those rational curves which are tangent to $`F`$. They correspond to a (continuous) set of points on a curve $`\gamma \beta `$, that will turn out to be a null-geodesic. Indeed, all we have to prove is $`\ddot{\gamma }=0(\text{mod}\dot{\gamma })`$, and $`\dot{\gamma }_x`$ corresponds to a section $`\eta `$ of $`N_x`$, vanishing at $`\overline{\beta }Z_x`$; as $`N_x𝒪(1)`$, $`\eta `$ is determined by its derivative at $`\overline{\beta }`$, which is a linear map $`T_{\overline{\beta }}F/T_{\overline{\beta }}`$ (the infinitesimal deformation of the direction of $`Z_x`$ within $`F`$). As the points of $`\gamma `$ correspond to lines tangent to $`F`$, we have that $`\ddot{\gamma }_x`$ corresponds to a section of $`N_x`$ collinear to $`\eta `$, thus $`\gamma `$ verifies the equation of a (non-parameterized) geodesic. See , , and Section 4 for details.
Example. The space of null-geodesics of $`𝐌=(E)\times (E)^{}`$ is the total space of a $`\times ^1`$-bundle over $`Z=`$, the flag manifold (see Section 8); a 2-plane $`FT_{(L,l)}`$ which corresponds to a null-geodesic in M is identified either to a projective diffeomorphism $`\phi :(l)(L^o)`$ (Section 8.4, case 3), or to a point $`Al,AL`$, resp. a plane $`a`$ containing $`L`$, and different from $`l`$ (Section 8.4, cases 2 and $`\mathrm{𝟐}^{}`$).
## 3. The structure of the ambitwistor space and the field of $`\alpha `$-cones
Conventions. Except for some results in Section 5, we will consider M to be a self-dual civilized 4-manifold, i.e. the (twistor) space $`Z`$ of $`\beta `$-surfaces of M is a Hausdorff smooth complex 3-manifold, and the projection $`(V_{})Z`$ is a submersion (e.g. M is geodesically convex), see .
We will frequently identify, following the deformation theory of Kodaira (see ), the vectors in $`T_x𝐌`$ with sections in the normal bundle $`N(Z_x)`$ of the projective line $`Z_x`$ in $`Z`$.
We also consider the space of null-geodesics $`B`$, as an open subsetset of $`(T^{}Z)`$.
For a null-geodesic $`\gamma `$, resp. a $`\beta `$-surface $`\beta 𝐌`$, we denote by $`\overline{\gamma }`$, resp. $`\overline{\beta }`$, the corresponding point in $`B`$, resp. $`Z`$.
### 3.1. $`\alpha `$\- and $`\beta `$-cones on the ambitwistor space
The vectors on $`B`$ can be expressed in terms of infinitesimal deformations of geodesics of M (Jacobi fields). More precisely:
$$T_{\overline{\gamma }}B𝒥_\gamma ^{}/𝒥_\gamma ^\gamma ,$$
where, for a null-geodesic $`\gamma `$, $`𝒥_\gamma ^{}`$ is the space of Jacobi fields $`J`$ such that $`_{\dot{\gamma }}J\dot{\gamma }`$, and $`𝒥_\gamma ^\gamma `$ is its subspace of Jacobi fields “along” $`\gamma `$, i.e. $`J\dot{\gamma }`$ in any point of the geodesic.
Remark.
A class in $`𝒥_\gamma ^{}/𝒥_\gamma ^\gamma `$ is represented by Jacobi fields yielding the same local section of the normal bundle $`N(\gamma )`$ of $`\gamma `$ in M. This is equivalent to the following obvious fact:
###### Lemma 1.
The kernel of the natural application $`𝒥_\gamma ^{}N(\gamma )`$ is $`𝒥_\gamma ^\gamma `$.
As a consequence, Jacobi fields on $`\gamma `$ induce particular local sections in $`N(\gamma )`$, which turn out to be solutions of a differential operator of order 1 on $`N(\gamma )`$, see Section 4.
The conformal geometry of M induces a particular structure on $`B`$: we describe it in order to obtain an expression of $`W^+`$ in terms of the geometry of the (ambi-)twistor space.
We have a canonical hyperplane $`V_{\overline{\gamma }}`$ in $`T_{\overline{\gamma }}B`$, defined by
$$V\overline{\gamma }:=𝒥_\gamma ^{}/𝒥_\gamma ^\gamma ,$$
where $`𝒥_\gamma ^{}`$ is the set of Jacobi fields $`J`$ everywhere orthogonal to $`\dot{\gamma }`$ (i.e. $`_{\dot{\gamma }}J\dot{\gamma }`$ and $`J\dot{\gamma }`$).
We deine now two fiels of cones in $`TB`$, both contained in $`V\overline{\gamma }`$:
###### Definition 6.
Let $`\gamma `$ be a null-geodesic in M, and, for each point $`x\gamma `$, let $`F_x^\beta `$ be the $`\beta `$-plane containing $`\dot{\gamma }_x`$. The (infinitesimal) $`\beta `$-cone $`V_{\overline{\gamma }}^\beta `$ at $`\overline{\gamma }B`$ is defined as follows:
$$V_{\overline{\gamma }}^\beta :=𝒥_\gamma ^\beta /𝒥_\gamma ^\gamma 𝒥_\gamma ^{}/𝒥_\gamma ^\gamma V_{\overline{\gamma }}T_{\overline{\gamma }}B,$$
where $`𝒥_\gamma ^\beta `$ is the set of Jacobi fields $`J`$ on $`\gamma `$ satisfying the condition
$$x\gamma \text{ such that }J_x=0\text{ and }(_{\dot{\gamma }}J)_xF_x^\beta .$$
###### Proposition 3.
The $`\beta `$-cone $`V_{\overline{\gamma }}^\beta `$ is flat, i.e. it is included in the 2-plane $`F_{\overline{\gamma }}^\beta `$ consisting of Jacobi fields contained in the $`\beta `$-plane defined by $`\dot{\gamma }`$ in each point of it.
###### Proof.
We have to prove that $`𝒥_\gamma ^\beta `$ is included in $`\overline{𝒥}_\gamma ^\beta `$, defined as follows:
$$\overline{𝒥}_\gamma ^\beta :=\{J\text{ Jacobi field on }\gamma |J_x,\dot{J}_xF_x^\beta ,x\gamma \}.$$
We will prove that $`𝒥_\gamma ^\beta \overline{𝒥}_\gamma ^\beta `$, therefore it will follow that the latter is non-empty, and is a linear space of dimension 2.
We denote by $`J^0`$ the parallel displacement, along $`\gamma `$, of a non-zero vector in $`F_x^\beta `$, transverse to $`\dot{\gamma }`$. Then $`J^0T\beta |_\gamma T\gamma `$, because $`\gamma `$ is included in the is totally geodesic surface $`\beta `$, thus we can characterize $`F_y^\beta `$ as the set $`\{XT_yM|X\dot{\gamma },XJ^0\}`$, for any $`y\gamma `$. We then observe that
$$\dot{\gamma }.\dot{J},J^0=R(\dot{\gamma },J)\dot{\gamma },J^0=R(\dot{\gamma },J^0)\dot{\gamma },J=kJ^0,J,$$
because $`R(\dot{\gamma },J^0)\dot{\gamma }`$ is in $`F^\beta `$, thus $`R(\dot{\gamma },J^0)\dot{\gamma }=h\dot{\gamma }+kJ^0.`$ So the scalar function $`J,J^0`$ satisfies to a linear second order equation, hence it it determined by its initial value and derivative. It follows then that it is identically zero, thus $`JF^\beta `$ everywhere, as claimed. ∎
Another subset in $`T_{\overline{\gamma }}B`$ is the $`\alpha `$-cone $`V_{\overline{\gamma }}^\alpha `$, defined as follows:
###### Definition 7.
Let $`\gamma `$ be a null-geodesic in M, and, for each point $`x\gamma `$, let $`F_x^\alpha `$ be the $`\alpha `$-plane containing $`\dot{\gamma }_x`$. The (infinitesimal) $`\alpha `$-cone $`V_{\overline{\gamma }}^\alpha `$ at $`\overline{\gamma }B`$ is defined as follows:
$$V_{\overline{\gamma }}^\alpha :=𝒥_\gamma ^\alpha /𝒥_\gamma ^\gamma 𝒥_\gamma ^{}/𝒥_\gamma ^\gamma V_{\overline{\gamma }}T_{\overline{\gamma }}B,$$
where $`𝒥_\gamma ^\alpha `$ is the set of Jacobi fields $`J`$ on $`\gamma `$ satisfying the condition
$$x\gamma \text{ such that }J_x=0\text{ and }(_{\dot{\gamma }}J)_xF_x^\alpha .$$
It is important to note that, in general, the projective curves $`(V_{\overline{\gamma }}^\alpha )`$ and $`(V_{\overline{\gamma }}^\beta )`$ are non compact, as each of them corresponds to the set of points on $`\gamma `$, which is non-compact, in general. The field of $`\alpha `$-cones on $`B`$ is the object of main interest in this paper. We may already guess that its flatness (i.e. the situation when $`V_{\overline{\gamma }}^\alpha `$ is a subset in a 2-plane) can be related to some special property of the conformal structure of M.
Remark.
We have seen that $`V_{\overline{\gamma }}^\beta `$ is included in the 2-plane $`F_{\overline{\gamma }}^\beta `$, i.e. the condition $`J_x=0,\dot{J}_xF_x`$ can be generalised to the linear condition $`J,\dot{J}F^\beta `$, but there is no canonical way of supplying the “missing” points of $`\gamma `$ with some appropriate Jacobi fields in order to “complete” $`V_{\overline{\gamma }}^\alpha `$ as in the $`\beta `$-cones case. This would be possible, for example, if $`(V_{\overline{\gamma }}^\alpha )`$ would be an open subset in a projective line. But the defect of $`V_{\overline{\gamma }}^\alpha `$ to be part of a 2-plane is measured by its projective curvature, and we will see in Section 4 that the vanishing of the latter implies the vanishing of $`W^+`$ (Theorem 1). $`B_{\overline{\gamma },x}^\alpha `$
### 3.2. Integral $`\alpha `$-cones in $`Z`$ and $`B`$
We study now the field of $`\alpha `$-cones of $`B`$ in relation with $`Z`$ and the canonical projection $`\pi :BZ`$. First, we note that there are complex projective lines in $`B`$ tangent to the directions in $`V_{\overline{\gamma }}^\alpha `$:
###### Definition 8.
Let $`\overline{\gamma }B`$, $`x\gamma `$ a point on the null-geodesic $`\gamma `$; let $`F_x^\alpha `$ be the $`\alpha `$-plane tangent to $`T_x\gamma `$. The rational curve $`B_{\overline{\gamma },x}^\alpha `$ in $`B`$ (containing $`\overline{\gamma }`$), is by definition the set of null-geodesics passing through $`x`$ and tangent to $`F_x^\alpha `$.
The curves $`B_{\overline{\gamma },x}^\alpha ,x\gamma `$ are projected by $`\pi `$ onto the complex lines $`Z_x`$ through $`\overline{\beta }`$ (corresponding to the $`\beta `$-surface $`\beta `$ containing $`\gamma `$) tangent to the 2-plane $`F^\gamma `$.
On the other hand, it is easy to see that the complex projective lines $`B_{\overline{\gamma },x}^\beta `$ (defined in an analogous way to $`B_{\overline{\gamma },x}^\alpha `$), which are tangent to (an open set of the directions of) $`V_{\overline{\gamma }}^\beta `$, are contained in the fibers of $`\pi `$. In fact, they coincide with some of the projective lines passing through the point $`\gamma (T_{\overline{\beta }}^{}Z)^2`$.
###### Definition 9.
The integral $`\alpha `$-cones in $`B`$, resp. $`Z`$ are defined by:
$$B_{\overline{\gamma }}^\alpha :=\underset{x\gamma }{}B_{\overline{\gamma },x}^\alpha \text{(}\beta \text{-cone in }B\text{)};Z^\gamma :=\underset{x\gamma }{}Z_x\text{(}\beta \text{-cone in }Z\text{)}.$$
We intend to prove that $`B_{\overline{\gamma }}^\alpha `$ is the canonical lift of $`Z^\gamma `$ (see Proposition 5). We know that $`\pi (B_{\overline{\gamma }}^\alpha )=Z^\gamma `$. We have then the following:
###### Proposition 4.
Except for the vertices $`\overline{\gamma }B_{\overline{\gamma }}^\alpha `$ and $`\overline{\beta }Z^\gamma `$, the two integral cones $`B_{\overline{\gamma }}^\alpha `$ and $`Z^\gamma `$ are smooth, immersed, surfaces of $`B`$, resp. $`Z`$.
###### Proof.
The open set of $`B`$ which is the space of null-geodesics of M can be viewed as the space of integral curves of the geodesic distribution $`G`$ of lines in $`(C)`$, the total space of the fibre bundle of isotropic directions in $`T𝐌`$. $`G_v`$ is defined as the horizontal lift (for the Levi-Civita connection on M) of $`v`$, which is an isotropic line in $`T_x𝐌`$. This definition is independent of the chosen metric and connection , and, by integrating this distribution (as M is civilized), we get a holomorphic map $`p:(C)B`$, where (an open set of) $`B`$ is the space of leaves of this foliation. This map can be used to compute the normal bundle of $`B_{\overline{\gamma },x}^\alpha `$, $`N(B_{\overline{\gamma },x}^\alpha )`$, see ,,.
Indeed, we have lines $`C_{\gamma ,x}^\alpha (C)_x`$, such that $`\dot{\gamma }_xC_{\gamma ,x}^\alpha `$, which project onto $`B_{\overline{\gamma },x}^\alpha `$, thus we get the following exact sequence of normal bundles:
$$0N(C_{\gamma ,x}^\alpha ;p^1(B_{\overline{\gamma },x}^\alpha ))N(C_{\gamma ,x}^\alpha ;(C))N(B_{\overline{\gamma },x}^\alpha ;B)0,$$
where we have written the ambient spaces of the normal bundles on the second position. The central bundle is trivial ($`C_{\gamma ,x}^\alpha `$ is trivially embedded in $`(C)_x^1\times ^1`$, which is trivially embedded in $`(C)`$ as a fibre), and it is easy to check that the left hand bundle is isomorphic to the tautological bundle over $`^1`$, $`𝒪(1)`$. This proves that $`N(B_{\overline{\gamma },x}^\alpha ;B)𝒪(0)𝒪(0)𝒪(1)`$, in particular the conditions in the completeness theorem of Kodaira are satisfied. Thus the lines in the integral $`\alpha `$-cone $`B_{\overline{\gamma }}^\alpha `$ form an analytic subfamily of the family $`\{B_{\overline{\gamma },x}^\alpha \}_{\overline{\gamma }B,x\gamma M}`$, that correspond to the sections of the normal bundle of $`B_{\overline{\gamma },x}^\alpha `$, vanishing at $`\overline{\gamma }B`$, or, equivalently, to the points $`x`$ of $`\gamma 𝐌`$.
But, in order to prove the smoothness of $`B_{\overline{\gamma }}^\alpha \{\overline{\gamma }\}`$, we first remark that the surface $`C_\gamma ^\alpha (C)`$, defined as follows, is smooth:
$$C_\gamma ^\alpha :=\{v(C)_x|x\gamma ,vF_\gamma ^\alpha \},$$
where $`F_\gamma ^\alpha `$ is the $`\alpha `$-plane containing $`\dot{\gamma }`$. $`C_\gamma ^\alpha `$ is smooth, and $`p(C_\gamma ^\alpha )=B_{\overline{\gamma }}^\alpha `$. We note now that $`C_\gamma ^\alpha `$ is everywhere, with the exception of the points of $`p^1(\overline{\gamma })`$, transverse to the fibers of the submersion $`p:(C)B`$. We may conclude that $`B_{\overline{\gamma }}^\alpha \{\gamma \}`$ is a smooth analytic submanifold of $`B`$ (not closed).
We can use similar methods to prove that $`Z^\gamma \{\overline{\beta }\}`$ is an immersed submanifold of $`Z`$ (by using the projection $`\pi :BZ`$). ∎
There is another argument for this latter claim, which gives the tangent space to $`Z^\gamma `$ in any point:
We see $`Z^\gamma `$ as the “trajectory” of a 1-parameter deformation of $`Z_x`$: we fix $`\overline{\beta }`$ and we “turn” $`Z_x`$ around $`\overline{\beta }`$ by keeping it tangent to $`F^\gamma `$. The trajectory of this deformation is smooth in $`\zeta Z^\gamma \beta `$ iff any non-identically zero section $`\nu `$ of the normal bundle $`N(Z_x)`$ corresponding to this 1-parameter deformation) does not vanish at $`\zeta `$. In particular, the tangent space $`T_\zeta Z^\gamma `$ is spanned by $`T_\zeta Z_x`$ and $`\nu (\zeta )`$.
But the sections $`\nu `$ generating this deformation are the sections of $`N(Z_x)`$ vanishing at $`\overline{\beta }`$, and they vanish at only one point (and even there, only at order 0) unless they are identically zero, because $`N(Z_x)𝒪(1)𝒪(1)`$.
Remark.
The values of these sections in the points of $`Z_x`$ other than $`\overline{\beta }`$, plus their derivatives in $`\overline{\beta }`$ (well-defined as they all vanish at $`\overline{\beta }`$), define a 1-dimensional subbundle of $`N(Z_x)`$, which is isomorphic to $`𝒪(1)`$. In fact, we have a 1-1 correspondence between the subbundles of $`N(Z_x)`$ isomorphic to $`𝒪(1)`$ and the 2-planes in $`T_{\overline{\beta }}Z`$. Then, the space of holomorphic sections of such a bundle is a linear space of dimension 2, consisting in a family of sections of $`N(Z_x)`$ vanishing on different points of $`Z_x`$. Thus we get a 2-plane $`F^\alpha `$ of isotropic vectors in $`T_x𝐌`$, which is easily seen to be an $`\alpha `$-plane, as the $`\beta `$-plane $`F_x^\beta =T_x\beta `$ consists in the set of all sections of $`N(Z_x)`$ vanishing at $`\overline{\beta }`$ (we have $`F_x^\alpha T_x\beta =T_x\gamma `$). The tangent space to $`Z^\gamma `$, in a point $`\zeta Z_x`$, is spanned by the subbundle of $`N(Z_x)`$ (isomorphic to $`𝒪(1)`$ — see above), defined by the isotropic vectors $`vF_x^\alpha `$. If $`\gamma ^\zeta `$ is the null-geodesic generated by $`v^\zeta `$, we conclude that $`T_\zeta Z^\gamma `$ is the 2-plane determined by $`\gamma ^\zeta `$, and that $`\zeta =\pi (\overline{\gamma ^\zeta })`$.
Example. If $`𝐌=(E)\times (E)^{}`$, then the integral $`\alpha `$-cone $`Z^\gamma `$ in $`Z`$, for $`\gamma F^\gamma =F^\phi T_{(L,l)}Z`$ (where $`\phi :(l)(L^o)`$ is a projective diffeomorphism), is the (smooth away from the vertex $`(L,l)`$) surface $`\{(S,\phi (s))|SL,sl,\phi (sl)=S\}`$. Its compactification (by adding the special cycle $`\overline{Z}_{(L,l)}`$) is singular (Section 8.4).
As any smooth surface in $`Z`$ has a canonical lift in $`B=(T^{}Z)`$, we get:
###### Proposition 5.
The integral $`\alpha `$-cone $`B_{\overline{\gamma }}^\alpha `$ is the canonical lift of the integral $`\alpha `$-cone on $`Z`$, $`Z^\gamma `$.
Remark.
Basically, this lift can only be defined for $`Z^\gamma \{\overline{\beta }\}`$, but, in this special case, it can be extended by continuity to $`\overline{\beta }`$. Of course, the smoothness of the lifted surface can only be deduced away from the vertex $`\overline{\gamma }`$ (from the smoothness of $`Z^\gamma \{\overline{\beta }\}`$).
## 4. The projective curvature of the $`\alpha `$-cone $`V_\gamma ^\alpha `$ and the self-dual Weyl tensor $`W^+`$ on M
As noted in Section 3, we intend to find a relation between the “curvature” of the $`\alpha `$-cone $`V_{\overline{\gamma }}^\alpha `$ (its non-flatness) and the Weyl tensor $`W^+`$ of $`(𝐌,c)`$. We begin by defining the projective curvature of $`V_\gamma ^\alpha `$: A projective structure on a manifold $`X`$ is an equivalence class of linear connections yielding the same geodesics. In such a space, we can define the projective curvature of a curve $`S`$ in a point $`\sigma `$ as the linear application $`k:T_\sigma ST_\sigma SN(S)_\sigma =T_\sigma X/T_\sigma S`$, with $`k(Y):=_YY`$(modulo $`T_\sigma S`$), for $``$ any connection in the projective structure of $`X`$. In particular, we take for $`X`$ the projective space $`(T_{\overline{\gamma }}B)`$, with its canonical projective structure, and for $`S`$ we take $`(V_{\overline{\gamma }}^\alpha )`$, the projectivized $`\alpha `$-cone in $`\overline{\gamma }`$.
###### Definition 10.
The projective curvature of the $`\alpha `$-cone $`V_{\overline{\gamma }}^\alpha `$ at the generating line $`\sigma V_{\overline{\gamma }}^\alpha `$ is the projective curvature of $`(V_{\overline{\gamma }}^\alpha )`$ in $`\sigma `$, and is identified to a linear application
$$K_{\gamma ,x}^\alpha :T_\sigma ST_\sigma SN(S)_\sigma ,$$
where $`\sigma `$ is the tangent direction to $`B_{\overline{\gamma },x}^\alpha `$ in $`\overline{\gamma }`$, and $`S:=(V_{\overline{\gamma }}^\alpha )`$.
In order to compute the projective curvature of $`V_{\overline{\gamma }}^\alpha `$, we establish first some canonical isomorphisms between the spaces appearing in the above definition and some linear subspaces of $`T_x𝐌`$. We will fix now the geodesic $`\gamma `$, the point $`x\gamma `$ (therefore also $`\sigma =T_{\overline{\gamma }}B_{\overline{\gamma },x}^\alpha (T_{\overline{\gamma }}B)`$), and, thus, the $`\alpha `$-plane $`F_x^\alpha T_x𝐌`$ containing $`\dot{\gamma }_x`$, as well as $`\dot{\gamma }_x^{}T_x𝐌`$, the orthogonal space to $`\dot{\gamma }_x`$.
For simplicity, in the following lemmas we will currently omit some indices referring to these fixed objects.
###### Lemma 2.
There is a canonical isomorphism $`\tau `$ between the tangent space $`T_\sigma S`$ to the projective cone $`S=(V_{\overline{\gamma }}^\alpha )`$ and the tangent space $`T_x\gamma `$ to the geodesic $`\gamma `$ in the point $`x`$ corresponding to the direction $`\sigma (T_{\overline{\gamma }}B)`$.
###### Proof.
Let $`YT_x\gamma `$. We will define $`\tau ^1(Y)`$ as follows: Recall that $`T_\sigma S\text{Hom}(\sigma ,E/\sigma )`$, where $`E(=E_x):=T_\sigma V_{\overline{\gamma }}^\alpha `$ (the tangent space in a point to a cone depends only on the line containing the point). We know that $`\sigma `$ corresponds to $`𝒥_{\gamma ,x}^\alpha `$, the space of Jacobi fields on $`\gamma `$, vanishing at $`x`$, and such that $`\dot{J}_xF^\alpha `$. It will be shown in the proof of the next theorem that $`E`$ consists of classes of Jacobi vector fields such that $`J_x,\dot{J}_xF^\alpha `$, (4).
Then, on a representative Jacobi field $`J𝒥_{\gamma ,x}^\alpha `$, we define $`\tau ^1(Y)`$ to be the class of Jacobi fields in $`E/\sigma `$, represented by the following Jacobi field $`J^Y`$ on $`\gamma `$, which is given by $`J_x^Y:=_YJ,\dot{J}_x^Y:=0`$. We remark that $`_YJ`$ is what we usually note $`\dot{J}`$, when the parameter on $`\gamma `$ is understood.
It is straightforward to check that $`JJ^Y`$ induces an isomorphism $`\tau ^1(Y):\sigma E/\sigma `$ for each non-zero $`J\sigma =𝒥_{\gamma ,x}^\alpha /𝒥_\gamma ^\gamma `$. ∎
We remark that $`V_{\overline{\gamma }}^\alpha V_{\overline{\gamma }}`$, the 4-dimensional subspace represented by Jacobi fields $`J`$, such that $`J,\dot{J}\dot{\gamma }`$. We further introduce the subspace $`H_{\overline{\gamma },x}^\alpha V_{\overline{\gamma }}`$, represented by Jacobi fields $`J`$ as before, with the additional condition $`J_xF_x^\alpha `$. It is a 3-dimensional subspace, and it contains $`E_x`$. The curvature of $`V_{\overline{\gamma }}^\alpha `$ will take values in $`\text{Hom}(TSTS,N^V(S))`$, and we will show (6) in the proof of the next theorem that it takes values in a smaller space, $`\text{Hom}(TSTS,N^H(S))`$. $`N_\sigma ^V(S)`$ is just the normal space of $`S`$ in $`(V_\gamma )`$ at $`\sigma `$, and $`N^H(S)`$ is the subspace of $`N_\sigma ^V(S)`$ consisting in elements represented by $`\xi \text{Hom}(\sigma ,H_{\overline{\gamma },x}^\alpha )\text{Hom}(\sigma ,V_{\overline{\gamma }})`$.
###### Lemma 3.
There is a canonical isomorphism
$$\rho :N^H(S)\mathrm{Hom}(F^\alpha /T\gamma ,\gamma ^{}/F^\alpha ).$$
###### Proof.
As $`H`$ is a subbundle of the normal bundle $`N(S)`$, $`N^H(S)`$ is isomorphic to$`\text{Hom}(\sigma ,H/E)`$. As in Lemma 2, we will construct the inverse isomorphism $`\rho ^1`$: Let $`\xi :F^\alpha /T\gamma \gamma ^{}/F^\alpha `$ be a linear application. Let $`\xi _0:F^\alpha \gamma ^{}`$ be a representant of $`\xi `$ (it involves a choice of a complementary space to $`F^\alpha `$ in $`\gamma ^{}`$). We define $`\rho ^1(\xi )\text{Hom}(\sigma ,H/E)`$ as being induced by the following linear application between spaces of Jacobi fields on $`\gamma `$:
$`\rho ^1(\xi ):𝒥_{\gamma ,x}^\alpha 𝒥_{\gamma ,x}^{\alpha ,}`$, where the second space corresponds to $`H_x`$, i.e. it contains Jacobi fields $`J`$ such that $`J_xF^\alpha ,\dot{J}_x\dot{\gamma }_x`$. Consider a parameterization of $`\gamma `$ around $`x`$, and let $`J𝒥_{\gamma ,x}^\alpha `$. We define $`J^\xi :=\rho ^1(\xi )(J)`$ by $`J_x^\xi :=0,\dot{J}_x^\xi :=\xi _0(\dot{J}_x)`$, and it is easy to check that the class of $`J^\xi `$ in $`H/E`$ is independent of the representant $`\xi _0`$, such that $`\rho ^1`$ is well-defined. It is also obviously invertible. ∎
We are now in position to translate the projective curvature of $`V_\gamma ^\alpha `$ in terms of conformal invariants of $`(𝐌,c)`$:
###### Theorem 1.
Let $`x`$ be a point in a null-geodesic $`\gamma `$. Then the projective curvature $`K`$ of the $`\alpha `$-cone $`V_{\overline{\gamma }}^\alpha `$ at $`\sigma `$ (corresponding to $`x`$, see Definition 10), which is a linear map
$$K:T_\sigma ST_\sigma SN^V(S)_\sigma ,$$
takes values in $`N^H(S)_\sigma `$ (see above), and is canonically identified to the linear map
$$K^{}:T_x\gamma T_x\gamma \mathrm{Hom}(F_x^\alpha /T_x\gamma ,\gamma _x^{}/F_x^\alpha ),$$
defined by the self-dual Weyl tensor of M:
$$K^{}(Y,Y)(X)=W^+(Y,X)Y,YT_x\gamma ,XF_x^\alpha .$$
###### Proof.
Consider the following analytic map, which parameterizes, locally around $`x\gamma `$, the deformations of the geodesic $`\gamma `$ that correspond to points contained in the integral $`\alpha `$-cone $`B_{\overline{\gamma }}^\alpha `$:
$$f:U𝐌,f(t,s,u)=\gamma ^{t,s}(u),$$
where $`U`$ is a neighborhood of the origin in $`^3`$, and $`\gamma ^{t,s}`$ is a deformation of the null-geodesic $`\gamma `$, such that
$$\gamma ^{t,s}(t)=\gamma (t),\dot{\gamma }^{t,s}(t)F_{\gamma (t)}^\alpha ,$$
where the parameterization of the geodesic $`\gamma `$ satisfies $`\gamma (0)=x`$, and $`F_{\gamma (u)}^\alpha `$ is the $`\alpha `$-plane in $`T_{\gamma (u)}𝐌`$ containing $`\dot{\gamma }(u)`$.
Convention We know that $`f`$ is defined around the origin in $`^3`$, so there exists a polydisc centered in the origin included in $`U`$, therefore all the relations that we will use are true for values of the variables $`t,s,u`$ sufficiently close to 0. For simplicity, we will omit to mention these domains.
The geodesics $`\gamma ^{t,s}`$ correspond to points in $`B_{\overline{\gamma },\gamma (t)}^\alpha `$, and the Jacobi fields $`J^t`$ on $`\gamma `$, defined as
$$J^t(u):=_sf(t,0,u)T_{\gamma (u)}𝐌,$$
correspond to vectors in $`V_{\overline{\gamma }}^\alpha `$ tangent to the above mentioned lines. We suppose that the deformation $`f`$ is effective, i.e. $`_u\gamma ^{t,s}(u)0`$ and $`J^t𝒥_\gamma ^\gamma `$, which is equivalent to $`\dot{J}^t(t)T_{\gamma (t)}\gamma `$. In order to compute the projective curvature of $`V_{\overline{\gamma }}^\alpha `$, we need thus to study the (second order) infinitesimal variation of these Jacobi fields on $`\gamma `$. As they are determined by their value and first order derivative in $`\gamma (0)=x`$, we need to evaluate $`_tJ^t(0)|_{t=o},_t\dot{J}^t(0)|_{t=0}`$ for the first derivative of $`J^t`$ at $`t=0`$, and $`_t^2J^t(0)|_{t=0},_t^2\dot{J}^t(0)|_{t=0}`$ for the second. Dots mean, as before, covariant differentiation with respect to the “speed” vector $`\dot{\gamma }`$, thus correspond to the operator $`_u`$.
As the covariant derivation $``$ has no torsion, we can apply the usual commutativity relations between the operators $`_t,_s,_t`$ and use them to differentiate the following equation, which follows directly from the definition of $`f`$ and $`J^t`$:
(1)
$$J^t(t)=0t.$$
We get then
(2)
$$_tJ^t(t)+\dot{J}^t(t)=0,$$
We recall now that, besides (1), we have $`\dot{J}^t(t)F_{\gamma (t)}^\alpha `$, thus $`\dot{J}^t(t)`$ is isotropic, which implies that:
(3)
$$_t\dot{J}^t(t),\dot{J}^t(t)=0,$$
as $`\ddot{J}^t(t)=R(\dot{\gamma }(t),J^t(t))\dot{\gamma }(t)=0`$. Equations (2) and (3) prove that
(4)
$$_tJ^t|_{t=0}𝒥_{\gamma ,x}^\alpha ,$$
which completes the proof of Lemma 2. From (3), it equally follows that $`_t\dot{J}^t(t)`$ is isotropic, and, by differentiating (3), we get
(5)
$$_t^2\dot{J}^t(t),\dot{J}^t(t)=_t\ddot{J}^t(t),\dot{J}^t(t).$$
From, (2) we have that $`_tJ^t(t)`$ is isotropic, and also
$$_t^2J^t(t)+2_t\dot{J}^t(t)=0,$$
which, together with (3), implies that $`_t^2J^t(0)|_{t=0}F_x^\alpha `$. We have then
(6)
$$_t^2J^t|_{t=0}𝒥_{\gamma ,x}^{\alpha ,},$$
which proves that the curvature $`K`$ of the $`\alpha `$-cone takes values in $`N^H(S)`$, as it is represented by $`_t^2J^t|_{t=0}`$.
In view of the Lemmas 2 and 3, it is clear now that the projective curvature $`K`$ is represented by the following application:
$$(\dot{\gamma },\dot{\gamma },\dot{J})_x_t^2J^t(0)|_{t=0}.$$
From (5), as $`_t\ddot{J}^t(t)=R(\dot{\gamma },_tJ^t)\dot{\gamma }`$ and $`_tJ^t(t)=\dot{J}^t(t)`$, we get
$$K(\dot{\gamma },\dot{\gamma })(\dot{J}),\dot{J}=R(\dot{\gamma },\dot{J})\dot{\gamma },\dot{J}.$$
The right hand side actually involves only $`W^+`$, as the other components of the Riemannian curvature vanish on this combination of vectors, thus we can replace $`R`$ with $`W^+`$ inn the above relation. On the other hand, the class of $`W^+(\dot{\gamma },\dot{J})\dot{\gamma }`$ modulo $`F^\alpha `$ is determined by its scalar product with $`\dot{J}`$, which represents a non-zero generator of $`F^\alpha /T\gamma `$.
The proof of the Theorem is now complete. ∎
Remark.
We may ask whether the projective lines in $`Z`$ are the geodesics of some projective structure. Indeed, in the conformally flat case, when M is the Grassmannian of 2-planes in $`^4`$ (the complexification of the Moebius 4-sphere), $`Z^3`$, and the complex lines are given by the standard (flat) projective structure. But there are two reasons (related to each other, as we will soon see) why $`Z`$ cannot carry a canonical projective structure: First, we do not necessarily have projective lines $`Z_x\overline{\beta }`$ in every direction of $`T_{\overline{\beta }}Z`$ (this would mean that $`\beta ^2`$, see next Section for a treatment of this problem), and second, the lift of a 2-plane $`F^\gamma T_{\overline{\beta }}Z`$ would be a 2-plane in $`T_{\overline{\gamma }}B`$, so $`V_{\overline{\gamma }}^\alpha `$ would be a flat cone:
###### Corollary 1.
The projective lines $`Z_x`$ in the twistor space $`Z`$ are geodesics of a projective structure iff it is projectively flat, and M is conformally flat.
###### Proof.
If $`Z`$ admits a projective structure, some of whose geodesics are the lines $`Z_x`$, then we have, for a fixed $`\overline{\beta }Z`$, a linear connection around $`\overline{\beta }`$, whose geodesics in the directions of $`Z_x,\overline{\beta }Z_x(x\beta 𝐌)`$ coincide, locally, with $`Z_x`$. This means that the integral $`\alpha `$-cone $`Z^\gamma `$, for $`\gamma \beta `$ a null-geodesic, is part of a complex surface (namely $`\mathrm{exp}(F^\gamma )`$, where $`F^\gamma T_{\overline{\beta }}Z`$ is the 2-plane corresponding to $`\gamma `$). Then the integral $`\alpha `$-cone $`B_{\overline{\gamma }}^\alpha `$, the lift to $`B`$ of $`Z^\gamma `$, is also a complex surface, thus $`V_{\overline{\gamma }}^\alpha `$ is a subset of the tangent space $`T_{\overline{\gamma }}B_{\overline{\gamma }}^\alpha `$, thus a flat cone. As this is true for all points of $`Z`$ and for all null-geodesics $`\gamma `$, Theorem 1 implies that M is flat.
On the other hand, it is well-known that the twistor space of a conformally flat manifold admits a flat projective structure, for which the projective lines $`Z_x`$ are geodesics, . ∎
## 5. Compactness of null-geodesics and conformal flatness
### 5.1. Complete $`\alpha `$-cones in $`Z`$
We have given, in the preceding Section, a way to measure the projective curvature of the $`\alpha `$-cone in $`B`$; we shall see now what happens in the special case when this cone is complete in a point $`\overline{\gamma }`$, i.e. when $`(V_{\overline{\gamma }}^\alpha )`$ is a compact submanifold in $`(T_\gamma B)`$.
This situation appears for example if, for any direction in $`F^\gamma T_{\overline{\beta }}Z`$, there are projective lines in $`Z`$ tangent to it.
###### Theorem 2.
Let $`Z`$ be the twistor space of the connected civilized self-dual 4-manifold $`(𝐌,c)`$, and suppose that, for a point $`\beta Z`$ and for a 2-plane $`F^\alpha T_{\overline{\beta }}Z`$, there are projective lines $`Z_x`$ tangent to each direction of $`F^\alpha `$. Then $`(𝐌,c)`$ is conformally flat.
###### Proof.
The idea is to prove that the integral $`\alpha `$-cone $`Z^\gamma `$ is a smooth surface. We know that this holds in all its points except for the vertex $`\overline{\beta }`$ (Proposition 4). The fact that all direction in $`F^\gamma `$ admits a tangent line is a necessary condition for this cone to be a smooth surface, as it needs to be well-defined around $`\overline{\beta }`$.
We choose an auxiliary hermitian (real) metric $`h`$ on $`Z`$. Its restrictions $`h_x`$ to the lines $`Z_xZ^\gamma `$ yield Kählerian metrics on these lines; in fact these metrics are deformations of one another, just like the lines $`Z_x`$ are. This means that the metrics $`h_x`$ depend continuously on $`x(F^\alpha )`$, a parameter in a compact set. We can therefore find a lower bound $`r_0>0`$ for the injectivity radius of all $`(Z_x,h_x)`$ at $`\overline{\beta }`$, and a finite upper bound $`R`$ for the norm of all the second fundamental forms $`H_x:TZ_xTZ_x(TZ_x)^{}(TZ).`$ We can also suppose that $`r_0`$ is smaller than the injectivity radius of $`(Z,h)`$ at $`\overline{\beta }`$.
The first step is to prove that $`Z^\gamma `$ is a submanifold of class $`𝒞^1`$. As its tangent space is everywhere a complex subspace of $`TZ`$, it will follow that it is a complex analytic submanifold.
Consider now the exponential map $`\mathrm{exp}_{\overline{\beta }}:T_{\overline{\beta }}ZZ`$, defined for the metric $`h`$; If we restrict it to a ball of radius less than $`r_0`$, it is a diffeomorphism into $`Z`$. The image of the complex plane $`F^\alpha `$ is then a smooth 4-dimensional real submanifold $`S`$ of $`Z`$, and there exists a positive number $`r_1`$ such that the exponential map in the directions normal to $`S`$,
$$\mathrm{exp}_S:TS^{}Z,\mathrm{exp}(Y):=\mathrm{exp}_y(Y),\text{ for }YT_yS^{},$$
restricted to the vectors of length less than $`r_1`$, is a diffeomorphism.
The image of this diffeomorphism is a tubular neighborhood of $`S`$, and we will denote by $`N(S,r)`$ such a tubular neighborhood of “width” $`r`$, for $`r<r_1`$.
The existence of an upper bound $`R`$ for the second fundamental forms of $`Z_x,x\gamma `$ implies the following fact:
###### Lemma 4.
For any $`r<r_1`$, there is a neighborhood $`UT_{\overline{\beta }}Z`$ of the origin such that $`\mathrm{exp}(U)Z^\gamma `$ is contained in $`N(S,r)`$, and is transverse to the fibers of the orthogonal projection $`p^S:N(S,r)S,p^S(\mathrm{exp}(Y)):=y`$, where $`YT_yS`$.
This is standard if $`Z^\gamma `$ is a submanifold; but it is also true in our case, where $`Z^\gamma `$ is a union of submanifolds $`Z_x`$.
Now it is easy to prove that $`Z^\gamma `$ is a $`𝒞^1`$ submanifold of $`Z`$ (the projection $`p^S`$ yields a local $`𝒞^1`$ diffeomorphism from a neighborhood of $`\overline{\beta }`$ in $`S`$ to a neighborhood of $`\overline{\beta }`$ in $`Z^\gamma `$; it is $`𝒞^1`$ in $`\overline{\beta }`$ because $`S`$ is tangent to $`Z^\gamma `$ at $`\overline{\beta }`$).
So $`Z^\gamma `$ is a $`𝒞^1`$ submanifold of $`Z`$; Its tangent space is complex in each point, thus $`Z^\gamma `$ is a complex-analytic surface immersed in $`Z`$.
We have then that $`B_{\overline{\gamma }}^\alpha B=(T^{}Z)`$, being the lift of $`Z^\gamma `$, is a smooth analytic surface immersed in $`B`$, in particular the $`\alpha `$-cone $`V_{\overline{\gamma }}^\alpha `$ is a complex plane.
Theorem 1 implies that $`W^+`$ vanishes on the $`\alpha `$-plane $`F_x^\alpha T_x𝐌`$ which contains $`\dot{\gamma }_x`$, for every point $`x\gamma `$. Now, the plane $`F^\gamma T_{\overline{\beta }}Z`$ is not the only one admitting projective lines $`Z_x`$ tangent to any of its directions: all planes “close” to $`F^\gamma `$ have the same property. Then $`W^+`$ vanishes on a neighborhood of $`\gamma `$, hence on the whole connected manifold M. ∎
Remark.
There is a more general situation where the integral $`\alpha `$-cone $`Z^\gamma `$ through $`\beta `$ is smooth in $`\beta `$:
###### Theorem 2.
Suppose that, for each direction $`\sigma (T_\beta Z)`$, there is a smooth (non-necessarily compact) curve $`Z_\sigma `$ tangent to $`\sigma `$, such that :
(i) if $`\sigma `$ is tangent to a projective line $`Z_x`$, then $`Z_\sigma =Z_x`$;
(ii) $`Z_\sigma `$ varies smoothly with $`\sigma (F^\gamma )`$.
Then
$$\overline{Z}_\beta ^\gamma :=\underset{\sigma (F^\gamma )}{}Z_\sigma $$
is a smooth surface around $`\beta `$, containing the $`\alpha `$-cone $`Z^\gamma `$ and $`W^+(F_x^\gamma )=0,x\gamma `$, where $`F_x^\gamma T_x𝐌`$ is the $`\alpha `$-plane containing $`\dot{\gamma }`$.
The proof is similar to the one of the previous theorem. Note that, if there is a direction $`\sigma `$ which is not tangent to a projective line $`Z_x`$, we cannot apply the deformation argument in Theorem 2 to conclude that $`W^+`$ vanishes everywhere.
Example. If $`𝐌=(E)\times (E)^{}`$, then $`Z=`$ and there are some particular planes for which the conditions in Theorem 2 are satisfied, although Theorem 2 never applies to $`Z`$: for a generic 2-plane $`F^\gamma `$, the $`\alpha `$-cone $`V_\gamma ^\alpha `$ is not flat. The above mentioned particular planes in $`TZ`$ correspond to the vanishing of $`W^+`$ on some particular $`\alpha `$-planes, but M is not anti-self-dual (see Sections 8.3, and also 8.7, 8.8).
### 5.2. Compact, simply-connected null-geodesics in M
Theorem 2 suggests that the existence of a compact null-geodesic diffeomorphic to $`^1`$ yields strong constraints upon the conformal structure of M. In fact, we have:
###### Theorem 3.
If a connected complex self-dual 4-manifold $`(𝐌,c)`$ admits a compact null-geodesic diffeomorphic to $`^1`$, the conformal structure of M is flat.
Remark.
A similar result has been proven by Y.-G. Ye using algebraic geometry techniques : if a projective complex manifoldadmits a conformal structure having a compact null-geodesic diffeomorphic to $`^1`$, it is conformally flat . Note that we do not need M to be compact in Theorem 3 ; on the other hand, we assume it to be self-dual.
###### Proof.
We first remark that the main difficulty is the definition of $`B`$, the space of null-geodesics , and of $`Z`$, the twistor space of $`(𝐌,c)`$, as M is not necessarily civilized. This is only possible on small open sets, but, in general, we can not expect to have any global construction of this kind. Thus, things that were almost obvious in the twistorial framework (like the existence of compact deformations of the null-geodesic $`\gamma `$), seem much more difficult to prove directly. The idea is to prove that all null-geodesics close to $`\gamma `$ are diffeomorphic to $`^1`$. Then, we show that, conversely, every projective line which is a deformation of $`\gamma `$ as a compact curve is a null-geodesic. In particular, sections in the normal bundle $`N(\gamma )`$ are induced by (local) Jacobi fields. We obtain then directly that $`W^+=0`$.
###### Proposition 6.
Let $`\gamma `$ be an immersed null-geodesic, diffeomorphic to a projective line $`^1`$. Then any local Jacobi field $`J`$ with $`\dot{J}\dot{\gamma }`$ induces a global normal field $`\nu ^J`$ on $`\gamma `$.
###### Proof.
If $`\gamma `$ is just a compact geodesic, it may have points of self-intersection, but it is always an immersed curve. It is more convenient then to think of $`\gamma `$ as a projective line immersed in M rather than the image of this immersion. The tangent, normal bundles, etc. are also to be thought as bundles over this projective line, still denoted by $`\gamma `$. Tubular neighborhoods of $`\gamma `$ are then neighborhoods of the zero section in the normal bundle $`N(\gamma )`$, small enough to be immersed (non-injectively) in M as a neighborhood of the image of $`\gamma `$.
We first notice that $`\gamma `$ may be decomposed in the union of a two open sets $`U_1U_2`$, both biholomorphic to the unit disk in $``$, and such that $`U_1U_2`$ is connected.
Then, for any local metric in $`c`$, we have a Jacobi equation around a point $`x\gamma `$, and a Jacobi field $`J`$ corresponding to prescribed $`J_x,\dot{J}_x`$. It is easy to prove that the (local) normal field induced by $`J`$ is independent of the chosen metric. Moreover, this normal field is the unique solution, for the prescribed 1-jet in $`x`$ induced by $`J_x,\dot{J}_x`$, of a second order differential equation on $`N(\gamma )`$:
###### Lemma 5.
The Jacobi equations for null-geodesic induce a second order linear differential operator $`P`$ on $`N(\gamma )`$, depending only on the conformal structure $`c`$ of M.
###### Proof.
For a Levi-Civita connection $``$ of a local metric on M, we locally define the following differential operator on $`T𝐌|_\gamma `$:
$$P:\mathrm{\Gamma }(T𝐌|_\gamma S^2(T\gamma ))\mathrm{\Gamma }(T𝐌|_\gamma ),$$
by $`P(Y,X,X):=_X_XY_{_XX}YR(X,Y)X`$. It obviously induces a (local) differential operator on $`N(\gamma )`$, and all we need to show is that, for a different connection $`^{}`$, the corresponding operator $`P^{}`$ induces the same one on $`N(\gamma )`$. First we write
$$P(Y,X,X)=_X[X,Y]+_{[X,Y]}X[_XX,Y],$$
then we recall that another Levi-Civita connection $`^{}`$ is related to $``$ by the formula :
$$_A^{}B=_AB+\theta (A)B+\theta (B)Ag(A,B)\theta ^{\mathrm{}},\text{ for }\theta \mathrm{\Lambda }^1𝐌,$$
so we directly obtain:
$$P^{}(Y,X,X)P(Y,X,X)=2(_Y\theta )(X)X+2\theta (_XY),$$
thus they induce the same operator on $`N(\gamma )`$. This one is, therefore, globally defined (the topology of $`\gamma `$ is not important). ∎
Now, for any $`x\gamma `$, we have an unique solution $`\nu `$ of $`P`$, for a prescribed 1-jet $`j^1(\nu )_x`$ (which consists in the values in $`x`$ of $`\nu `$ and of his first-order derivative), globally defined on every contractible open set $`Ux`$ in $`\gamma `$. This is because on any such contractible set the equation $`P\nu =0`$ becomes a second order ordinary linear equation on a disk in $``$, which admits global holomorphic solutions (unique if we fix the initial conditions).
Take now $`xU_1U_2`$. Then, the two solutions $`\nu _1`$ and $`\nu _2`$, defined on $`U_1`$, resp. $`U_2`$, coincide on the connected $`U_1U_2`$, so they yield a global solution $`\nu `$ with the prescribed initial conditions in $`x`$.
In particular, this solution is a global section of the normal bundle $`N(\gamma )`$. ∎
After the infinitesimal result, the local one:
###### Proposition 7.
Small deformations of $`\gamma `$ are also compact immersed projective lines.
Remark.
The tubular neighborhoods considered below are always seen as images, by a local diffeomorphism, of subsets — which are, generally, fiber bundles over $`^1`$ — of the normal bundle of $`\gamma `$, resp. $`\stackrel{~}{\gamma }`$, the lift of $`\gamma `$ to $`(C)`$. We need this because of the possible self-intersections of $`\gamma `$; $`\stackrel{~}{\gamma }`$ is always embedded.
###### Proof.
Consider an auxiliary hermitian metric $`h`$ on M.Then $`h|_\gamma `$ induces the same topology like a round metric $`h_0`$ on the sphere $`S^2`$. We can define a tubular neighborhood $`N(r_1)`$ of $`\gamma `$, as the open set of points $`y𝐌`$ with $`d(y,\gamma )<r_1`$. We choose $`r_1`$ small enough for $`N(r_1)`$ to be a fiber bundle over $`\gamma `$ (the fiber $`N(r_1)_x`$, for $`x\gamma `$, being the image of the real 4-plane of $`T_x𝐌`$, $`h`$-orthogonal to $`T_x\gamma `$, by the exponential of $`h`$).
Take now a finite number of contractible open sets whose union covers $`\gamma `$, and we choose holomorphic metrics in $`c`$ on each of this open sets. Then, on these sets we have connections, and it is well-known that any point in such a set has a basis of geodesically connected neighborhoods , . We choose a finite number of such geodesically convex open sets, that cover $`\gamma `$, and such that they are all included in $`N(r_1)`$. We have then
$$N(r_1)\underset{i=1}{\overset{n}{}}U_i\gamma .$$
(Of course, they are geodesically convex only with respect to some particular metric in $`c`$, but we are interested only in the implications involving the null-geodesics, which are independent of the metric.)
It is immediate, , that a geodesically convex set $`U_i`$ has the following property: all maximal geodesics are closed submanifolds of $`U_i`$ and are contractible as topological spaces. This is particularly true for null-geodesics included in $`U_i`$.
We refine now the covering by another one, $`U_i^{}\overline{U}_i^{}U_i`$, such that $`U_i^{}=N(r_2)|_{V_i}`$, where $`V_i:=U_i^{}\gamma `$, such that $`\gamma _{i=1}^nU_i^{}`$. In fact, we ask for $`U_i^{}`$ to be restrictions to $`V_i`$ of a tubular neighborhood $`N(r_2)`$, for $`r_2`$ sufficiently small ($`U_i^{}`$ are “cylindrical” neighborhoods). We can easily imagine how to find such a refinement of the initial covering.
The proof of the proposition now follows two ideas: first, we consider a very special covering by disks (for the round metric $`h_0`$) of $`\gamma `$; second, we define a neighborhood of $`\stackrel{~}{\gamma }`$ in $`(C)`$ of isotropic directions for which, using the special covering of $`\gamma `$, we can extend the associated null-geodesics and eventually get compact ones. We call $`\stackrel{~}{\gamma }`$ the canonical lift of a null-geodesic in $`(C)`$; it is an integral curve of the geodesic distribution of lines in $`(C)`$.
The first idea has nothing to do with complex analysis; it is just a matter of metric topology on the round sphere $`(S^2,\text{can})`$.
###### Lemma 6.
Let $`\{V_i\}_{i=\overline{1,n}}`$ be a covering of $`(S^2,\text{can})(\gamma ,h_0)`$ by open sets. Then there exists a positive number $`r_0`$ and a finite set of points $`x_j\gamma `$, for $`j=\overline{1,N}`$, such that:
(i) All disks $`D(x,10r_0)`$ are contained in at least one of the open sets $`V_i`$;
(ii) The disks $`D(x_j,r_0)_{j=\overline{1,N}}`$ cover $`\gamma S^2`$;
(iii) $`100r_0<l`$, where $`l`$ is the diameter of $`\gamma S^2`$.
An important property of this covering is that all sets, as well as the intersections of a finite number of them, are convex for the round metric, thus contractible.
We intend to extend null-geodesics which can be projected diffeomorphically onto $`\gamma `$ by means of the fiber projection $`\rho :N(r_1)\gamma `$. We will do that step-by-step, extending it over disks $`D(x,R)`$ of increasing radius. But before that, we need to restrict ourselves to some particularly “close to $`\gamma `$” null directions.
We need two things: the extensions need to remain within $`N(r_2)`$, and they should also be transverse to the fibers of $`\rho `$, otherwise the projection into $`\gamma `$ would not be an invertible diffeomorphism.
First, we consider the following compact subset of $`(T𝐌)`$:
$$S:=\{LT_yM|y\overline{N(r_1)},L\text{ tangent to the fibers of }\rho \},$$
where we say that a complex line $`L`$ is tangent to a real manifold $`\rho ^1(x),x\gamma `$ if it contains a non-zero (real) vector tangent to this real submanifold.
The hermitian metric $`h`$ on M induces a metric on $`(T𝐌)`$, and also one on $`(C)`$. We can, then, evaluate the distance between $`\stackrel{~}{\gamma }`$ and $`S`$:
$$\mu _0:=d(\stackrel{~}{\gamma },S)>0,$$
as they are disjoint compact sets.
Following LeBrun , we can define the complex 5-manifolds $`B^i`$ as the spaces of null-geodesics of $`(U_i,c)`$, equivalently the space of integral curves of the geodesic distribution in $`(C)|_{U_i}`$. The projections $`p_i:(C)|_{U_i}B^i`$, which send an isotropic direction to the null-geodesic tangent to it, is a submersion and the (closed) fibers are precisely the lifts of the null-geodesics of $`U_i`$. This construction is possible because $`U_i`$ are geodesically convex (for a particular local holomorphic metric in $`c`$), see for details.
We first consider a tubular neighborhood $`\stackrel{~}{N}(r^0)`$ of $`\stackrel{~}{\gamma }`$ in $`(C)`$ which projects, by $`\pi :(C)𝐌`$ inside $`N(r_2)`$, and such that $`100r^0<\mu _0`$. This second condition ensures that all directions in $`\stackrel{~}{N}(r^0)`$ are transverse to the fibers of $`\rho `$.
Consider then the following neighborhoods of $`\stackrel{~}{\gamma }|_{\overline{V}_i}`$: $`p_i:(C)|_{U_i}B^i`$ is an open application, so we define $`C_i`$ to be $`p_i^1(C_B^i)`$, where $`C_B^i`$ is an open neighborhood of $`\overline{\gamma }B^i`$ contained in $`p_i(\stackrel{~}{N}(r_0))`$.
It is important to note that $`C_i`$ have the following property:
###### Lemma 7.
for any point $`YC_i`$, the null-geodesic $`\gamma ^Y`$ tangent to $`Y`$, contained in $`U_i`$, lies into $`N(r_2)`$ and is always transverse to the fibers of $`\rho :N(r_1)\gamma `$. Hence, its restriction to the points of $`U_i^{}`$ projects diffeomorphically onto $`V_i\gamma `$. Moreover, all the points of $`\stackrel{~}{\gamma }^Y`$ that lie over $`\overline{V}_i`$ are in $`C_i`$.
Obviously, the crucial property of these open sets is that every null-geodesic starting there is totally contained in $`C_i`$, at least the part that “lies over” (in the sense of the projection $`\rho `$) $`\overline{V}_i`$. After constructing these sets $`C_i,i`$, we define $`r^1>0`$ small enough for the tubular neighborhood $`\stackrel{~}{N}(r^1)`$, restricted to $`\overline{V}_i`$, to be contained into $`C_i`$, for all $`i`$. For each $`i`$, this means that $`r^1`$ has to be less than the minimum of the following continuous functions defined on the compact set $`\overline{V}_i`$:
$$\overline{V}_ixd(T_x\gamma ,(C)_xC_i).$$
This neighborhood $`\stackrel{~}{N}(r_1)`$ of $`\stackrel{~}{\gamma }`$ has the following property:
###### Lemma 8.
For each $`Y\stackrel{~}{N}(r_1)`$, and for each $`i`$ such that $`\pi (Y)=xV_i`$, we have $`YC_i`$, and thus the whole null-geodesic $`\gamma ^{Y,i}`$, contained in $`U_i`$, is lifted to $`\stackrel{~}{\gamma }^{Y,i}C_i`$.
The disadvantage of $`\stackrel{~}{N}(r_1)`$ is that it does not necessarily contain $`\stackrel{~}{\gamma }^{Y,i}`$. But we know that the latter is contained in $`\stackrel{~}{N}(r_0)`$, which contains the union of all $`C_i`$.
We recall now that the idea of proof is to extend a null-geodesic $`\gamma ^Y`$ close to $`\gamma `$ over the disks $`D_j:=D(x_j,r_0)`$. Every extension over a disk brings $`\gamma ^Y`$ from $`\stackrel{~}{N}(r^1)`$ to the larger set $`\stackrel{~}{N}(r^0)`$. As we have a finite, well-determined, number of disks $`N`$, all we need now to apply our extending idea is a sequence of open sets $`\stackrel{~}{N}(r^k)`$ such that
(7)
$$Y\stackrel{~}{N}(r^k),\text{ such that }x:=\pi (Y)V_i,\stackrel{~}{\gamma }^{Y,i}\stackrel{~}{N}(r^{k1}).$$
To do that, we construct $`C_i^1\stackrel{~}{N}(r^1)`$ as we have done for $`C_i`$, and then $`\stackrel{~}{N}(r^2)`$ by repeating the same procedure. We stop after $`N`$ (the number of disks covering $`(\gamma ,h_0)`$, see Lemma 6) steps and claim:
###### Proposition 8.
$`Y\stackrel{~}{N}(r^N)`$, the null-geodesic $`\gamma ^Y`$ extends to a compact curve which projects (via $`\rho `$) diffeomorphically onto $`\gamma `$.
###### Proof.
Fix $`x:=\pi (Y)`$. We can define $`\gamma ^Y`$ inside $`U_i`$, where $`U_i`$ is an open set containing $`x`$. In particular, $`\gamma ^Y`$ is well-defined over $`D(x_0,10r_0)`$, where $`x_0:=\rho (x)`$, see Lemma 6. (Of course, this is because $`\gamma ^Y|_{U_i}`$ is transverse to the fibers of $`\rho `$. We intend to extend it over disks centered in $`x_0`$.
Consider the domains $`D_i^x:=D_i^{x_0}`$ which are “quadrilaterals” contained in $`D(x_i,10r_i)`$ and containing $`D(x_i,r_0)`$, as in the following picture (the “vertical” parts of the border of $`D_i^x`$ are segments of circles centered in $`x_0`$):
We change, if necessary, the order of the indices $`i`$ of $`x_i`$ and $`D_i^x`$ such that it coincides with the ordering of increasing distances $`d(x_0,x_i)`$. We define than the open sets
$$\mathrm{\Delta }_k:=\{\begin{array}{cc}D(x_0,r_0)\hfill & \text{ if }d(x_0,x_k)<9r_0\text{,}\hfill \\ \mathrm{\Delta }_{k1}\hfill & \text{ if }d(x_0,x_k)>l10r_0\text{,}\hfill \\ \mathrm{\Delta }_{k1}D_k^x\hfill & \text{ otherwise}.\hfill \end{array}$$
Remark.
The closed disk $`\overline{D}(x_0,d(x_0,x_{k+1}r_0)`$ is included in $`\mathrm{\Delta }_k`$ as soon as $`x_kD(\overline{x}_0,10r_0)`$, where $`\overline{x}_0`$ is the point of $`\gamma =S^2`$ opposed to $`x_0`$. Then, because of the specific geometry of the domains $`D_i^x`$ (see Lemma 6), we easily conclude that the domains $`\mathrm{\Delta }_k`$ are contractible (along the geodesics of the sphere passing through $`x_0`$) and so are the intersections $`\mathrm{\Delta }_kD_{k+1}^x`$, too.
We prove, then, by induction, that $`\gamma ^Y`$ can be extended over $`\mathrm{\Delta }_k`$, and that all the corresponding points of $`\stackrel{~}{\gamma }^Y`$ are contained in $`\stackrel{~}{N}(r^{Nk})`$.
This is obvious for small values of $`k`$. When we add $`D_k^x`$ to $`\mathrm{\Delta }_{k1}`$, we consider a point $`z`$ in the connected (see above) intersection $`\mathrm{\Delta }_{k1}D_k^x`$. It is contained in the $`U_i^{}`$ that contains $`D(x_k,10r_0)`$ and, as $`\stackrel{~}{\gamma }^Y|_{\mathrm{\Delta }_{k1}}`$ is contained in $`\stackrel{~}{N}(r^{Nk+1})`$, the connected piece $`\stackrel{~}{\gamma }^Y|_{\mathrm{\Delta }_{k1}D_k^x}`$ is contained in $`C_i^{Nk}`$, and thus the whole extension $`\stackrel{~}{\gamma }^{Y,i}`$ in $`U_i`$ is contained in $`C_i^{Nk}`$, hence in $`\stackrel{~}{N}(r^{Nk})`$. The connectedness of the considered piece implies that $`\gamma ^{Y,i}|_{\mathrm{\Delta }_{k1}D_k^x}`$ coincides with $`\stackrel{~}{\gamma }^Y|_{\mathrm{\Delta }_{k1}D_k^x}`$. Thus we have obtained an extension of $`\gamma ^Y`$ to $`\mathrm{\Delta }_k`$, such that the corresponding lift lies in $`\stackrel{~}{N}(r^{Nk})`$, as claimed.
For large values of $`k`$, the $`\mathrm{\Delta }_k`$ are all identical to $`\mathrm{\Delta }_{N1}`$, which contains $`D(x_0,l10r_0)`$. We have thus proven that there is an extension of $`\gamma ^Y`$ over this disk, such that the corresponding points of $`\stackrel{~}{\gamma }^Y`$ are in $`\stackrel{~}{N}(r^1)`$. Consider then the disk $`D(\overline{x}_0,10r_0)`$; it is contained in some $`V_i`$:
The intersection $`\mathrm{\Delta }^x`$ of this disk with $`\mathrm{\Delta }_{N1}`$ is a connected open subset of $`V_i`$, and we know that $`\stackrel{~}{\gamma }^Y|_{\mathrm{\Delta }^x}`$ is contained in $`\stackrel{~}{N}(r^1)`$. This implies that we can extend in a unique way $`\stackrel{~}{\gamma }^Y|_{\mathrm{\Delta }^x}`$ to $`V_i`$, in particular to $`D(\overline{x}_0,10r_0)`$, and the corresponding points in $`\stackrel{~}{\gamma }^Y|_{V_i}`$ are in $`C_i\stackrel{~}{N}(r^0)`$.
We have proven that $`\gamma ^Y`$ extends over $`\gamma `$, i.e. there is a maximal extension (obviously unique) of the null-geodesic tangent to $`Y`$, such that it projects (via $`\rho `$) diffeomorphically onto $`\gamma `$. The projection is $`C^{\mathrm{}}`$, but the extended null-geodesic is clearly an analytic submanifold of M. ∎
The proof of Proposition 7 is now complete. ∎
Remark.
Generic deformations of $`\gamma `$ are embedded projective lines. Indeed, let $`\{x_1,\mathrm{},x_k\}`$ be the nodes (self-intersection points) of $`\gamma `$, and consider the manifold $`\stackrel{~}{𝐌}`$, obtained by blowing-up the points $`\{x_1,\mathrm{},x_k\}`$. Then, generically, any null-geodesic $`\gamma ^{}`$ close to $`\gamma `$ avoids these points, hence is diffeomorphic to its lift to $`\stackrel{~}{𝐌}`$, which is a deformation of the lift of $`\gamma `$, thus it is an immersed projective line. But the lift of $`\gamma `$ is embedded, and so must be its deformations, hence $`\gamma ^{}`$ is embedded.
The next step is to prove that all deformations of $`\gamma `$ as a compact curve are null-geodesics , by a dimension-counting argument; we need to compute the normal bundle of $`\gamma `$.
We ask now if the family of projective lines in M defined as the deformations of $`\gamma `$ is locally complete in the sense of Kodaira . For this, we need to prove that the dimension of the space of global sections in $`N(\gamma )`$ is equal to 5, i.e. to the dimension of $`B`$. The extensions of the null-geodesics close to $`\gamma `$ yield local diffeomorphisms between neighborhoods of $`\gamma `$ in $`B^i`$, resp. $`B^j`$. In fact, we have a projection $`p:\stackrel{~}{N}(r^N)W^B`$ onto the space of integral curves of the geodesic distribution in $`\stackrel{~}{N}(r^N)(C)`$. $`W^B`$ is the space of complex null-geodesic close to $`\gamma `$ in M. But essential for us is that $`p`$ is a submersion, fact that has important consequences for the normal bundle of $`\gamma `$ in M.
###### Proposition 9.
The normal bundle of $`\gamma `$ in M is isomorphic to $`𝒪(1)𝒪(1)𝒪(0)`$.
###### Proof.
It is well-known that all holomorphic bundles over $`^1`$ are direct sums of line bundles, all of which are isomorphic to $`𝒪(k),k`$.
We have the subbundle $`N^{}(\gamma )`$ of $`N(\gamma )`$, represented by vectors orthogonal to $`\dot{\gamma }`$. We have the following exact sequence:
(8)
$$0N^{}(\gamma )N(\gamma )N(\gamma )/N^{}(\gamma )0.$$
The right hand term of this sequence is a line bundle, and it admits global non-zero sections (extensions of Jacobi fields $`J`$ such that $`J,\dot{\gamma }`$ is a non-zero constant, see Proposition 6). $`N(\gamma )/N^{}(\gamma )`$ is then isomorphic to $`𝒪(a)`$, with $`a`$.
We denote by $`N^\beta (\gamma )`$ the subbundle of the normal bundle represented by vectors in $`T\beta `$. It admits global sections, namely the extensions of Jacobi fields contained in $`T_\beta `$, see Proposition 6. It is a line bundle, thus isomorphic to $`𝒪(c_1),c_11`$, as it contains global sections with prescribed 1-jet in a point.
Remark.
In general, $`𝒪(k)`$ is the line bundle over $`^1`$ admitting global sections for any prescribed $`k`$-jet in a point $`x`$. This section is unique, and it gives a unique value of the $`k+1`$-jet in $`x`$.
We have the following exact sequence:
(9)
$$0N^\beta (\gamma )N^{}(\gamma )N^{}(\gamma )/N^\beta (\gamma )0.$$
It is easy to check that the right hand term admits local sections represented by Jacobi fields for any prescribed 1-jet in a point $`x`$. Hence, $`N^{}(\gamma )/N^\beta (\gamma )𝒪(c_2),c_21`$.
All we can obtain now is that $`N^{}(\gamma )𝒪(b_1)𝒪(b_2)`$, with $`b_1+b_2=c_1+c_22`$. Actually $`b_1,b_21`$, otherwise all sections of $`N^{}`$ that vanish somewhere would be contained in a line subbundle, which would contradict Proposition 6.
We have then $`N(\gamma )𝒪(a)𝒪(b_1)𝒪(b_2),a0,b_1,b_21`$. We want to prove that there is equality in all these three inequalities.
We know, from Proposition 7, that there is a tubular neighborhood $`\stackrel{~}{N}(r^N)`$ of $`\stackrel{~}{\gamma }`$ in $`(C)`$ such that the null-geodesic distribution yields a foliation with compact leaves, and such that the projection onto the space $`W^B`$ of these compact curves is a submersion. (Of course, this space is nothing but the space of complex null-geodesics close to $`\gamma `$.)
It is obvious then that the normal bundle of $`\stackrel{~}{\gamma }`$ in $`(C)`$ is trivial, as $`\stackrel{~}{\gamma }`$ is e fiber of a submersion.
We have now the following exact sequence of bundles, related to the projection $`\pi :(C)𝐌`$:
(10)
$$0N^\pi (\stackrel{~}{\gamma })N(\stackrel{~}{\gamma })\pi ^{}N(\gamma )0,$$
where $`N^\pi (\stackrel{~}{\gamma })`$ is the normal subbundle of $`\stackrel{~}{\gamma }`$ represented by vectors tangent to the fibers of $`\pi `$. In a point $`T_x\gamma \stackrel{~}{\gamma }(C)`$, the fiber of $`\pi `$ is equal to $`(C)_x`$, so the tangent space to it is isomorphic to $`\text{Hom}(T_x\gamma ,T_x\gamma ^{}/T_x\gamma )`$, for the projective variety $`(C)_x(T_x𝐌)`$. Thus
$$N^\pi (\stackrel{~}{\gamma })\text{Hom}(T\gamma ,N^{}(\gamma ))𝒪(2)(𝒪(b_1)𝒪(b_2)).$$
The central bundle in the exact sequence (10) is trivial. The equation above then implies that the Chern number of $`N(\gamma )`$ is subject to the following constraint:
$$a+b_1+b_2+(b_12+b_22)=0,$$
thus, as $`b_1,b_21`$ and $`a0`$, we have $`a=0`$ and $`b_1=b_2=1`$.
As observed above, generic, compact, simply-connected null-geodesics are embedded. From now on, we suppose $`\gamma `$ is one of them.
Then $`H^1(N(\gamma ))=0`$, so we can apply the theory of Kodaira to deform $`\gamma `$, , so the dimension of the space of global sections of $`N(\gamma )`$ is $`1+2+2=5`$, the same as the space of complex null-geodesic close to $`\gamma `$, hence
###### Corollary 2.
The deformations of $`\gamma `$ as a compact curve are null-geodesics in $`(𝐌,c)`$.
This means that any global section in $`N(\gamma )`$ can be represented, locally, by the Jacobi fields that yield the same element in $`J^1N(\gamma )`$, the space of jets of order 1 in $`N(\gamma )`$.
Recall now the exact sequence (9); we conclude that
$$N^\beta (\gamma )N^{}(\gamma )/N^\beta (\gamma )𝒪(1).$$
We have the canonical isomorphism $`N^\alpha (\gamma )N^{}(\gamma )/N^\beta (\gamma )`$, coming from the restriction of the projection $`N^{}(\gamma )N^{}(\gamma )/N^\beta (\gamma )`$ (we denote by $`N^\alpha (\gamma )`$, resp. $`N^\beta (\gamma )`$, the subbundle of the normal bundle of $`\gamma `$, such that its fiber at $`x\gamma `$ is $`F_x^\alpha /T_x\gamma `$, resp. $`F_x^\beta /T_x\gamma `$).
We have thus $`N^\alpha (\gamma )𝒪(1)`$, and all the 1-jets of $`N^\alpha (\gamma )`$ yield global sections, thus local Jacobi fields.
But the existence of a Jacobi field in $`𝒥_\gamma ^\alpha 𝒥_\gamma ^\gamma `$ implies, by the Jacobi equation $`\ddot{J}=R(\dot{\gamma },J)\dot{\gamma }`$, that $`W^+(F^\alpha )=0`$, the self dual Weyl tensor vanishes on $`F^\alpha `$, the $`\alpha `$-plane generated by $`\dot{\gamma }`$.
We recall now that, for a fixed point $`x`$, $`W_x^+`$ is a polynomial of order 4, and it is zero on $`F^\alpha `$ for the compact null-geodesic $`\gamma `$. The same is true for the compact deformations of $`\gamma `$, which implies that $`W_x^+=0`$. This holds for all the points of $`\gamma `$, and also for all points covered by the deformations of $`\gamma `$. As these deformations cover at least an open set around $`\gamma `$, we conclude that $`W^+`$ vanishes on a non-empty open set, thus, being holomorphic, $`W^+=0`$ on the whole (connected) manifold M.
This completes the proof of Theorem 3
## 6. The projective structure of $`\beta `$-surfaces in a self-dual manifold
The null-geodesics contained in a $`\beta `$-surface $`\beta `$ define a projective structure on the totally-geodesic surface $`\beta `$, which is also given by any connection on $`\beta `$ induced by a Levi-Civita connection on M. We claim that this projective structure is flat, i.e. locally equivalent to $`^2`$.
Example. If $`𝐌=(E)\times (E)^{}`$, then a $`\beta `$-surface indexed by $`(L,l)`$ is $`\beta ^{(L,l)}=\{(A,a)|Al,La,Aa\}^2`$, and the null-geodesics in $`\beta ^{(L,l)}`$ are identified to the affine lines in $`^2`$ (see Section 8.5).
To prove the projective flatness of a 2-dimensional manifold $`\beta `$, we need to prove that the Thomas tensor $`T`$ vanishes identically . This tensor is an analog of the Cotton-York tensor in conformal geometry (there is also a Weyl tensor of a projective structure, but it only appears in dimensions greater than 2).
For a connection $``$ in the projective class of $`\beta `$, the Thomas tensor is defined as follows : For $`X,Y,ZT\beta ,`$
(11)
$$\begin{array}{cc}\hfill T(X,Y,Z):=& 2(_ZK)(Y)X+2(_YK)(Z)X\hfill \\ & (_ZK)(X)Y+(_YK)(X)Z,\hfill \end{array}$$
where the derivation involves only the curvature tensor $`R`$, and the $`K(Y)X:=\text{tr}R(Y,)X`$ is the trace of the endomorphism $`R(Y,)X\text{End}(T\beta )`$.
The Thomas tensor is independent of the connection $``$, therefore we will consider that $``$ is induced by a Levi-Civita connection on M.
###### Proposition 10.
The Thomas tensor of a $`\beta `$-surface can be expressed in terms of the anti-self-dual Cotton-York tensor of M, thus it is identically zero.
###### Proof.
We need first to define the anti-self-dual Cotton-York tensor as an irreducible component of the Cotton-York tensor of M.
Convention We note $`C`$ the Cotton-York tensor of $`(𝐌,c)`$; we will not use this letter for the isotropic cone in this Section, nor in the following one.
The Cotton-York tensor is not conformally invariant; its definition depends on a (local) metric $`g`$ in the conformal structure, which is supposed to be fixed :
(12)
$$C(X,Y)(Z):=(_Xh)(Y,Z)(_Yh)(X,Z),X,Y,ZT𝐌,$$
where $`h`$ is the normalized Ricci tensor of M,
(13)
$$h=\frac{1}{2n(n1)}Scalg+\frac{1}{n2}Ric_0,$$
$`Ric_0,Scal`$ being the trace-free Ricci tensor, resp. the scalar curvature of the metric $`g`$, and $`n:=\text{dim}𝐌`$. In our case, $`n=4`$, but the formula applies in all dimensions greater than 2 .
Remark.
The Cotton-York tensor $`C`$ of M is a 2-form with values in $`T^{}𝐌`$, thus it has two components $`C^+T^{}𝐌\mathrm{\Lambda }^+𝐌`$, and $`C^{}T^{}𝐌\mathrm{\Lambda }^{}𝐌`$. $`C`$ satisfies a first Bianchi identity, as $`h`$ is a symmetric tensor, and also a contracted (second) Bianchi identity, coming from the second Bianchi identity in Riemannian geometry, :
(14) $`{\displaystyle C(X,Y)(Z)}`$ $`=`$ $`0\text{ circular sum};`$
(15) $`{\displaystyle C(X,e_i)(e_i)}`$ $`=`$ $`0\text{ trace over an orthonormal basis}.`$
That means that $`C\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^1𝐌`$, and is orthogonal on $`\mathrm{\Lambda }^3𝐌\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^1𝐌`$ and on $`\mathrm{\Lambda }^1𝐌`$, which is identified with the image in $`\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^1𝐌`$ by the metric adjoint of the contraction (15).
Now, the Hodge operator $`:\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^2𝐌`$ induces a symmetric endomorphism of $`\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^1𝐌`$, which maps the two above spaces isomorphically into each other. This implies that $`C^+`$ and $`C^{}`$ satisfy (14) and (15) (note that these two relations are equivalent in their case).
The Cotton-York tensor is related to the Weyl tensor of M by the formula :
(16)
$$\delta W=C,$$
where $`\delta :\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^1𝐌`$ is induced by the codifferential on the second factor, and by the Levi-Civita connection $``$. Then, $`C^+`$ has to be the component of $`\delta W`$ in $`\mathrm{\Lambda }^1𝐌\mathrm{\Lambda }^+𝐌`$, and we know that the restriction of $`W^{}`$ to $`\mathrm{\Lambda }^2𝐌\mathrm{\Lambda }^+𝐌`$ is identically zero. This means that
(17) $`\delta W^+`$ $`=`$ $`C^+,\text{and also}`$
(18) $`\delta W^{}`$ $`=`$ $`C^{}.`$
Hence, as M is self-dual, $`C^{}`$ vanishes identically.
We can prove now that the Thomas tensor of a $`\beta `$-surface $`\beta `$ is identically zero: first we prove
(19)
$$K(Y)X=\text{tr}|_{T\beta }R(Y,)X=h(X,Y),X,YT\beta .$$
We recall that the suspension $`hI`$, viewed as an endomorphism of $`\mathrm{\Lambda }^2𝐌`$, is defined by :
(20)
$$(hI)(X,Y):=h(X)Yh(Y)X,X,YT𝐌,$$
where $`h`$ is identified with a symmetric endomorphism of $`T𝐌`$.
We have then the following decomposition of the Riemannian curvature :
$$R=hI+W^++W^{}.$$
Of course, if M is self-dual, $`W^{}=0`$ and $`W^+(X,Y)=0`$ if $`X,YT\beta `$ (in fact, the elements in $`\mathrm{\Lambda }^2F^\beta `$, for any $`\beta `$-plane $`F^\beta T_x𝐌`$, correspond to the isotropic vectors in $`\mathrm{\Lambda }^{}𝐌`$), because $`W^+|_{\mathrm{\Lambda }^{}M}=0`$. Then, if we choose the basis $`\{X,Y\}`$ in $`T\beta `$, we get
$`K(Y)X`$ $`=`$ $`\text{tr}|_{T\beta }(hI)(Y,)X=`$
$`=`$ $`\text{the component along }X\text{ of }(hI)(Y,X)X=`$
$`=`$ $`h(Y,X),`$
which proves (19). The Thomas tensor of the projective structure of $`\beta `$ has the following expression (see (11)):
$$T(X,Y,Z)=3(_Zh)(Y,X)+3(_Yh)(Z,X)=3C(Y,Z)(X),X,Y,ZT\beta ,$$
and, as $`C^+(,)(X)`$ vanishes on the anti-self-dual 2-form $`YZ`$, we conclude
(21)
$$T(X,Y,Z)=C^{}(Y,Z)(X)=0.$$
As the flatness of the projective structure on a 2-dimensional manifold is equivalent to the vanishing of its Thomas tensor , we get
###### Corollary 3.
The projective structure of the $`\beta `$-surfaces of a self-dual complex manifold M is flat.
From the classification of projectively flat compact complex surfaces (, see also ), we get then a classification of compact $`\beta `$-surfaces in M:
###### Theorem 4.
A compact $`\beta `$-surface of a self-dual complex 4-manifold belongs (up to finite covering) to the following classes:
1. $`^2`$;
2. a compact quotient of the complex-hyperbolic plane $`𝐇_{}^2/\mathrm{\Gamma }`$;
3. a compact complex surface admitting a (flat) affine structure:
(i) a Kodaira surface;
(ii) a properly elliptic surface with $`b_1`$ odd;
(iii) an affine Hopf surface;
(iv) an Inoue surface;
(v) a complex torus.
See , , for details.
## 7. Umbilic hypersurfaces in self-dual 4-manifolds
It has been shown by LeBrun that, for a given geodesically connected 3-manifold with conformal structure $`(𝐐,c^{})`$, there is a (germ-unique) self-dual 4-manifold $`(𝐌,c)`$, such that $`(𝐐,c^{})`$ is an umbilic hypersurface of $`(𝐌,c)`$. We can therefore note by $`c`$ the conformal structure of $`𝐐`$, as it coincides with the restriction of the conformal structure of M.
The technical tool used in the proof of this result is the twistor space $`Z`$ of $`(𝐐,c)`$, which is the space of complex null-geodesics of this manifold. It has been shown by LeBrun that $`Z`$ is a 3-manifold with a contact structure, and containing projective lines with normal bundle $`𝒪(1)𝒪(1)`$. Conversely, for any such manifold $`Z`$, the space of these lines tangent to the distribution of planes induced by the contact structure, is a conformal 3-manifold. On the other hand, $`Z`$ can be identified with the twistor space of a self-dual 4-manifold M, in which $`𝐐`$ is an umbilic hypersurface (the conformal infinity of a Einstein metric on M, ). $`Z`$ has an additional structure, namely a contact structure, represented by a (non-integrable) distribution of 2-planes $`F_\beta T_\beta Z`$, which corresponds to the space of Jacobi fields $`J\gamma ^\beta `$, where $`\gamma ^\beta =𝐐\beta `$ and $`Z`$ is considered as the twistor space of $`𝐐`$, or, equivalently, to the null-geodesic $`\gamma ^\beta 𝐌`$, if $`Z`$ is considered as the twistor space of M. We also remark that the above contact structure yields a section in the bundle $`BZ`$, which is never tangent to the $`\alpha `$-cones (as $`𝐐`$ is transverse to all $`\alpha `$-planes ).
Remark.
A holomorphic contact structure on a twistor space $`Z`$ does not necessarily determine a conformal 3-manifold: if $`Z=`$ with its contact structure (see Section 8.4), there is no rational curve in $`Z`$, with normal bundle $`𝒪(1)𝒪(1)`$, tangent to the corresponding distribution of planes. On the other hand, this contact structure yields an Einstein metric on $`𝐌=(E)\times (E)^{}`$, which admits the smooth 3-manifold $`\overline{𝐌}=(E)\times (E)^{}`$ as “infinity”. However, this “infinity” is not conformal (see Section 8.6).
In the 3-dimensional case, the conformally invariant tensor “measuring” the non-flatness of $`(𝐐,c)`$ is the Cotton-York tensor, defined in general by
(22)
$$C(X,Y):=_Xh(Y)_Yh(X),X,YT𝐐,$$
where $`h`$ is the normalized Ricci tensor , (13). For the 3-manifold $`𝐐𝐌`$, we have $`h(X)=\frac{1}{12}Scal+Ric_0`$.
It is natural to ask how is the Cotton-York tensor of $`𝐐`$ related to the Weyl tensor of M. We first recall a few facts of conformal geometry in dimensions 3 and 4.
For the 3-dimensional manifold $`𝐐`$, the Riemannian curvature has the following expression:
(23)
$$R^𝐐(X,Y)=(h\text{I})(X,Y):=h(X)Yh(Y)X,X,YT𝐐,$$
as there is no Weyl tensor (in general, $`h\text{I}`$ is the Ricci component of the curvature, ). If we introduce the Hodge operator $`^𝐐:\mathrm{\Lambda }^2𝐐\mathrm{\Lambda }^1𝐐`$, then $`R^𝐐`$ is equivalent to the symmetric 2-tensor $`^𝐐R^𝐐^𝐐`$. A straightforward application of the above formula yield
(24)
$$R^𝐐:=^𝐐R^𝐐^𝐐=h+(\mathrm{tr}h)I.$$
For the 4-dimensional manifold M, the components of the Riemannian curvature can also be expressed as eigenspaces of $``$-type operators. Namely, regarding $`R:=R^𝐌`$ as a symmetric endomorphism of $`\mathrm{\Lambda }^2𝐌=\mathrm{\Lambda }^+𝐌\mathrm{\Lambda }^{}𝐌`$, $`W^+`$ is the trace-free component of $`R`$ in $`\mathrm{End}(\mathrm{\Lambda }^+𝐌)`$, and $`W^{}`$ is the trace-free component of $`R`$ in $`\mathrm{End}(\mathrm{\Lambda }^{}𝐌)`$ .
Let $`𝐐𝐌`$ be a hypersurface, such that the restriction of the conformal structure $`c`$ of M to $`𝐐`$ is non-degenerate (equivalently, $`T𝐐`$ is nowhere tangent to an isotropic cone). We call $`c`$ the induced structure on $`𝐐`$. We suppose that $`𝐐`$ is umbilic.
Remark.
There is no possible choice of an “orientation” in this case. Indeed, the group of conformal transformations of $`^n`$, $`CO(n,)=O(n,)\times ^{}/\{\pm \mathrm{𝟏}\}`$ is non-connected if $`n`$ is even, and a choice of an orientation is a restriction of the frame bundle of a $`n`$-dimensional conformal manifold to the connected component of $`\mathrm{𝟏}CO(n,)`$. But if $`n`$ is odd, $`CO(n,)`$ is connected, so all $`CO(n,)`$-frames can be connected to each other by continuous paths. Therefore, although a complex-Riemannian 3-manifold admits, locally, 2 possible orientations, they are conformally equivalent, fact that makes impossible a canonical way to associate an orientation to a metric in the conformal class.
There is another way to see this difference between the even- and odd-dimensional conformal manifolds: Let $`(X,c)`$ be a $`n`$-dimensional conformal manifold. Then a (local) metric $`g`$ in the conformal class $`c`$ is a global section in $`L^2S^2T^{}X`$, where $`L`$ is the bundle of weighted scalars. We can canonically associate to $`c`$ a global section of $`\kappa ^2L^{2n}`$, which is the induced metric on the canonical bundle. For a given metric, there are only 2 (local) sections of $`\kappa `$ of “norm” $`1`$. We can pick one of them if we have a given section in $`\kappa L^n`$ (an “orientation”), and if we have a canonical way to associate to $`g`$, which is a section in $`L^2`$, a section of $`L^n`$ . This can be done if $`n`$ is even, namely $`g^{\frac{n}{2}}`$. We also remark that, if $`n`$ is even, all we need to define a conformal structure on $`X`$ is just the bundle $`L^2`$ (which is a $`n/2`$-th root of $`\kappa `$), but if $`n=2k+1`$ is odd, $`L^2`$ automatically gives us $`L\kappa L^{2k}`$.
We can canonically identify $`\mathrm{\Lambda }^+𝐌`$ and $`\mathrm{\Lambda }^{}𝐌`$, restricted to $`𝐐𝐌`$, to $`\mathrm{\Lambda }^2𝐐`$, by:
(25)
$$\begin{array}{ccc}\mathrm{\Lambda }^2𝐐\alpha \hfill & & \hfill \alpha +^𝐌\alpha \mathrm{\Lambda }^+𝐌\\ \mathrm{\Lambda }^2𝐐\alpha \hfill & & \hfill \alpha ^𝐌\alpha \mathrm{\Lambda }^{}𝐌.\end{array}$$
###### Theorem 5.
Let $`𝐐`$ be an umbilic hypersurface of a self-dual manifold M. Then:
(i) The Weyl tensor of M, restricted to $`𝐐`$, is identically zero;
(ii) The Cotton-York tensor of $`𝐐`$ is related to the self-dual Weyl tensor of M by the formula:
$$=g(_\nu W^+(A),B)_x=C(A)(^𝐐B)_x,$$
where $`A,B\mathrm{\Lambda }_x^2𝐐`$, $`\nu T_x𝐐`$ is unitary for the metric $`g`$, and the Hodge operator $`^𝐐`$ is induced by $`g`$ and the orientation on $`𝐐`$ admitting $`\nu `$ as an exterior normal vector.
###### Corollary 4.
If $`(𝐌,c)`$ is self-dual and $`𝐐𝐌`$ is an umbilic hypersurface, then the Cotton-York tensor of $`𝐐`$, $`C^𝐐`$, is identified to the restriction to $`𝐐`$ of the self-dual Cotton-York tensor $`C^+`$ of M:
$$C^+(X,Y)(Z)=C^𝐐(X,Y)(Z),X,Y,ZT𝐐.$$
Proof of the Theorem. The claimed identity is conformally invariant: If $`X,Y,Z,\nu `$ is a $`g`$-orthonormal oriented basis of M, then $`X,Y,Z`$ is a $`g`$-orthonormal basis on $`𝐐`$ giving the orientation as above. Then $`^𝐐(ZX)=Y`$, and, if we take $`A:=XY,B:=ZX`$, the claimed identity becomes
(26)
$$_\nu W^+(X,Y)Z,X=C(X,Y)(Y),$$
where angle brackets denote the scalar product induced by $`g`$.
The tensors $`W^+,C`$, in the above form, are independent of the chosen metric $`g`$ , which depends on the normal vector $`\nu `$, supposed to be $`g`$-unitary. If $`\nu ^{}:=\lambda \nu `$, for $`\lambda ^{}`$, then the corresponding metric $`g^{}=\lambda ^2g`$, and also $`_{}^{𝐐}{}_{}{}^{}=\lambda ^1^𝐐`$, thus the identity (26) for $`\nu ^{},g^{}`$ is equivalent to the one for $`\nu ,g`$.
Remark.
As $`W^+`$ is the trace free component of the Riemannian curvature contained in $`\text{End}(\mathrm{\Lambda }^+𝐌)`$, and is symmetric, it is enough to evaluate it on pairs $`A,B\mathrm{\Lambda }^2𝐐\mathrm{\Lambda }^+𝐌`$ which are unitary and orthogonal for the metric $`g`$, therefore the check of the equation (26) will prove the theorem.
As $`W^\pm `$ are $`^𝐌`$-eigenvectors in $`\text{End}_0(\mathrm{\Lambda }^2𝐌)`$ (the space of trace-free endomorphisms of $`\mathrm{\Lambda }^2𝐌`$), they are determined by the following formulas, where $`X,Y,Z`$ is any oriented orthonormal basis of $`T𝐐`$:
$`W^+(X,Y)Z,X`$ $`=`$ $`{\displaystyle \frac{1}{4}}(R(X,Y)Z,X+R(Z,\nu )Y,\nu +`$
$`+R(X,Y)Y,\nu +R(Z,\nu )Z,X)`$
$`W^{}(X,Y)Z,X`$ $`=`$ $`{\displaystyle \frac{1}{4}}(R(X,Y)Z,X+R(Z,\nu )Y,\nu `$
$`R(X,Y)Y,\nu R(Z,\nu )Z,X),`$
where $`X,Y,Z,\nu `$ is supposed to be a local extension, around a region of $`𝐐`$, of the $`g`$-orthonormal frame used in (26). As M is supposed to be self-dual, $`W^{}`$ is identically zero, thus, in the points $`x𝐐`$, we have
(29)
$$W^+(X,Y)Z,X_x=\frac{1}{2}(R(X,Y)Y,\nu +R(Z,X)Z,\nu )_x.$$
It is a standard fact that, if $`𝐐`$ is umbilic, there is a local metric $`g`$ in the conformal class $`c`$ of M, such that, for $`g`$, $`𝐐`$ is totally geodesic. We fix such a metric. Then we have
(30)
$$R(X^{},Y^{})Z^{}=R^𝐐(X^{},Y^{})Z^{},X^{},Y^{},Z^{}T𝐐,$$
which, together with 29, implies that $`W^+|_𝐐0`$.
On the other hand, (30), together with (29) and (7), yield
(31)
$$R(X,Y)Z,X_x+R(Z,\nu )Y,\nu _x=0,x𝐐.$$
Let us compute now the normal derivative of $`W^+`$ in a point $`x𝐐`$; we suppose that $`X,Y,Z,\nu `$ are locally extended by an orthonormal frame, and that they are parallel at $`x`$ (we omit, for simplicity of notation, the point $`x`$ in the following lines:
$$_\nu W^+(X,Y)Z,X=\frac{1}{2}(_\nu R(X,Y)Y,\nu +_\nu R(Z,X)Z,\nu ),$$
from (29). This is then equal to:
$$\begin{array}{ccc}\hfill _\nu W^+(X,Y)Z,X& =& \frac{1}{2}(_XR(Y,\nu )Y,\nu +_YR(\nu ,X)Y,\nu +\hfill \\ & & +_ZR(X,\nu )Z,\nu +_XR(\nu ,Z)Z,\nu ),\hfill \end{array}$$
from the second Bianchi identity. Then we have
$$\begin{array}{ccc}\hfill _\nu W^+(X,Y)Z,X& =& \frac{1}{2}(_XR(Z,X)Z,X+_YR(Z,Y)Z,X+\hfill \\ & & +_ZR(Y,Z)X,Y+_XR(Y,X)X,Y)\hfill \end{array}$$
from analogs of (31). Then
$$\begin{array}{ccc}\hfill _\nu W^+(X,Y)Z,X& =& \frac{1}{2}(_XR^𝐐(Z,X)Z,X+_YR^𝐐(Z,Y)Z,X+\hfill \\ & & +_ZR^𝐐(Y,Z)X,Y+_XR^𝐐(Y,X)X,Y)\hfill \end{array}$$
from (30)
$$\begin{array}{ccc}\hfill _\nu W^+(X,Y)Z,X& =& \frac{1}{2}(_Xh(Z,Z)+_Xh(X,X)+_Yh(Y,X)\hfill \\ & & _Zh(Z,X)_Xh(X,X)_Xh(Y,Y)),\hfill \end{array}$$
from (23). Finally, from (15), we get
$$_\nu W^+(X,Y)Z,X=\frac{1}{2}(C(X,Z)(Z)C(X,Y)(Y))=C(X,Y)(Y)$$
This proves equation (26).
The Corollary 4 now easily follows from the above theorem and (17).
The main results in Section 5 hold also in the case of a conformal (complex) 3-manifold:
###### Theorem 6.
Let $`Z`$ be a twistor space of a conformal 3-manifold $`𝐐`$; let $`F_{\overline{\gamma }}T_{\overline{\gamma }}Z`$ be its contact structure. Suppose there is a point $`\overline{\gamma }Z`$ such that, to any direction in $`F_{\overline{\gamma }}`$, there is a tangent rational curve in $`Z`$ with normal bundle $`𝒪(1)𝒪(1)`$. Then $`Z`$ is projectively flat, and $`𝐐`$ is conformally flat.
This follows directly from Theorem 2, as $`𝐐`$ is umbilic in M, the space of all projective lines in $`Z`$ with the above normal bundle , and from Theorem 5.
###### Theorem 7.
Let $`𝐐`$ be a conformal 3-manifold containing an immersed rational curve as null-geodesic. Then $`𝐐`$ is conformally flat.
###### Proof.
We cannot use directly Theorem 5, as the “ambient” self-dual manifold M can only be defined for a civilized (e.g. geodesically connected) 3-manifold. Hence, we follow the steps in Proposition 7 and prove
###### Proposition 11.
Null-geodesics close to a compact, simply-connected one are also compact and simply-connected, and they are embedded.
Still using the same arguments as in Section 5, we get an embedded null-geodesic $`\gamma 𝐐`$ diffeomorphic to $`^1`$, and we have:
###### Proposition 12.
The deformations of $`\gamma `$ as a compact curve coincide with the null-geodesics close to $`\gamma `$.
We cover $`\gamma `$ with geodesically convex open sets $`U_i,i=\overline{1,n}`$, such that:
(32)
$$ij\text{ s.t. }U_iU_j\gamma \mathrm{},U_{ij}(U_iU_j),$$
where $`U_{ij}`$ is still geodesically convex (with respect to a particular Levi-Civita connection). This is possible by choosing $`U_i,i=\overline{1,n}`$, small enough. Then we choose a relatively compact tubular neighborhood $`N(r_0)`$ of $`\gamma `$, such that its closure is covered by the $`U_i`$’s.
We consider then the twistor spaces $`Z_i`$, the spaces of null-geodesics of $`U_i`$. The compact, simply-connected, null-geodesics close to $`\gamma `$ identify (diffeomorphically) the neighborhoods of $`\overline{\gamma }^iZ_i`$ with the space $`Z`$ of the deformations of $`\gamma `$ as a compact curve. We can see then $`Z`$ as an open set common to all the $`Z_i`$’s:
Following LeBrun, we define the self-dual manifolds $`𝐌_i`$ as the space of projective lines in $`Z_i`$, with normal bundle $`𝒪(1)𝒪(1)`$. Then $`U_i`$ is an umbilic hypersurface in $`𝐌_i`$.
The local twistor spaces $`Z_i`$ admit contact structures, which coincide on $`Z`$, and contain projective lines $`Z_x^i`$ corresponding to points $`x\gamma U_i`$. If we denote by $`Z_{ij}`$ the twistor space of $`U_{ij}`$, then $`Z_i`$ and $`Z_j`$ are identified to open sets in $`Z_{ij}`$, in particular the lines $`Z_x^i`$ and $`Z_x^j`$ are identified, thus their intersections with the common set $`Z`$ coincide. We denote by $`Z_x`$ this (non-compact) curve in $`Z`$, and by $`F`$ the canonical contact structure of $`Z`$ (restricted from the ones of $`Z_i`$).
Remark.
We already have obtained that the $`\alpha `$-cone corresponding to $`F_{\overline{\gamma }}`$ is a part of a smooth surface: the union of the lines $`Z_x`$, $`x\gamma `$, thus, from Theorem 2, the Weyl tensor $`W_i^+`$ of the self-dual manifold $`𝐌_i`$ vanishes on the $`\alpha `$-planes generated by $`T\gamma `$. But this is nothing new: we know, from Theorem 5, that $`W_i^+`$ vanishes on $`U_i`$.
We intend to apply Theorem 2 to prove that $`W_i^+`$ vanishes on points close to $`U_i`$, but in $`𝐌_iU_i`$. We do that by showing that the integral $`\alpha `$-cones corresponding to planes $`F^yT_{\overline{\gamma }}Z`$ are parts of smooth surfaces, then we conclude using Theorem 2.
First we choose hermitian metrics $`h_i`$ on $`Z_i`$, such that they coincide (with $`h`$) on $`Z`$. We have a diffeomorphism between $`\gamma `$ and $`(F_{\overline{\gamma }})`$, so we choose relatively compact open sets in $`(T_{\overline{\gamma }}Z)`$, covering $`(F_{\overline{\gamma }})`$, with the following properties: As the metrics $`h_i`$ induce metrics on $`𝐌_i`$, we first choose a small enough distance $`r_1>0`$ such that
1. $`i`$, there is a sub-covering $`V_iU_i`$ of $`\gamma `$ such that the “tubular neighborhoods” $`W_i:=\{y𝐌_i|\mathrm{d}(y,\overline{V}_i)r_1\pi _i(y)\overline{V}_i\gamma \}`$ are compact ($`\pi _i`$ is the “orthogonal projection” — for the hermitian metric — from $`𝐌_i`$ to $`\gamma U_i`$; it is well defined because of the condition below);
2. $`r_1`$ is less than the bijectivity radius of the (hermitian) exponentials for the points of $`\overline{V}_i`$ in $`𝐌_i`$, and for the points of $`\overline{V_iV_j}`$ in $`𝐌_{ij}`$ (if $`U_iU_j\gamma \mathrm{}.`$).
We have then
###### Lemma 9.
For any $`y_iW_i𝐌_i`$, $`y_jW_j𝐌_j`$ such that the curves $`Z_{y_i}:=Z_{y_i}^iZ`$, $`Z_{y_j}:=Z_{y_j}^jZ`$ are tangent to the same direction in $`\overline{\gamma }Z`$, the respective curves $`Z_{y_i},Z_{y_j}`$ coincide.
###### Proof.
We first note that the projection $`\pi _i`$ from $`𝐌_i`$ is equivalent to the $`h`$\- orthogonal projection of the direction of $`T_{\overline{\gamma }}Z_{y_i}`$ to a direction in $`F_{\overline{\gamma }}`$, so $`\pi _i(y_i)=\pi _j(y_j)=:y\gamma `$; thus $`y`$ belongs to both $`U_i`$ and $`U_j`$, and we use again the twistor space $`Z_{ij}`$ to conclude that $`Z_{y_i}`$ and $`Z_{y_j}`$ are “restrictions” to $`Z`$ of the same projective line (as they both have the same tangent space at $`\overline{\gamma }`$) $`Z_{y_{ij}}^{ij}`$, for a point $`y_{ij}𝐌_{ij}`$. ∎
Now we have a tubular neighborhood $`S(T_{\overline{\gamma }}Z)`$ of $`(F_{\overline{\gamma }})`$, of radius $`r_1/2`$, such that, for any 2-plane $`F^{}S`$, the conditions in Theorem 2 are satisfied (considering any of the local twistor spaces $`Z_i`$).
We conclude that the Weyl tensor $`W_i^+`$ of $`𝐌_i`$ vanishes along all null-geodesics of $`𝐌_i`$, close (in $`W_i`$) to $`\gamma `$ and included in the $`\beta `$-surface $`\beta ^i`$, determined by $`\gamma `$. This means that $`W^+`$ vanishes everywhere on $`\beta ^i`$. By deforming $`\gamma `$, we obtain that $`W_i^+`$ vanishes on a neighborhood of $`U_i`$ in $`𝐌_i`$, hence $`𝐌_i`$, as well as $`U_i`$, are conformally flat (by Theorem 5).
It follows from Theorem 5 that $`𝐐`$ is conformally flat. ∎
## 8. Examples
### 8.1. The flat case
The first example is the “flat” case: $`Z=^3=(^4)`$, with its canonical projective structure, and its space of projective lines $`𝐌=\text{Gr}(2,^4)`$. ($`Z`$ is equally the twistor space of the Riemannian round 4-sphere, which is, therefore, a real part of $`\text{Gr}(2,^4)`$.) If $`\beta Z`$, then the $`\beta `$-surface associated to it is the set $`\{x\text{Gr}(2,^4)|\beta x^4\}`$. In this flat case, we can equally define the $`\alpha `$-twistor space $`Z^{}`$, which is the dual projective 3-space $`(^3)^{}:=((^4)^{})=\text{Gr}(3,^4)`$, and an $`\alpha `$-surface $`\alpha Z^{}`$ is the set $`\{x\text{Gr}(2,^4)|x\alpha ^4\}𝐌`$. A null-geodesic $`\gamma `$ is then determined by a pair of incident $`\alpha `$-, resp. $`\beta `$-surface $`\beta \alpha ^4`$:
$$\gamma =\{x\text{Gr}(2,^4)|\beta x\alpha \}.$$
$`\alpha `$-surfaces and $`\beta `$-surfaces are diffeomorphic to $`^2`$, null-geodesics to $`^1`$, and the ambitwistor space $`B`$ is the “partial flag” manifold
$$B=\{(\alpha ,\beta )(^3)^{}\times ^3|\beta \alpha \}.$$
The flag manifold, of dimension 7, is isomorphic to total space of the projective cone bundle over M, $`(C)`$.
### 8.2. $`^2`$
Another example is when $`Z`$ is the twistor space of the real Riemannian manifold $`^2`$, with the Fubini-Study metric. Then $`Z`$ is the manifold of flags in $`E=^3`$, $`:=\{(L,l)(E)\times (E)^{}|Ll\}`$, ($`(E)`$, resp. $`(E)^{}`$ are viewed as the space of lines, resp. 2-planes, in $`E`$) . A projective line $`Z_x`$ in $`Z`$ is a set
$$Z_x=\{(L,l)|La^x,A^xl\},$$
where $`(A^x,a^x)`$ belongs to $`(E)\times (E)^{}`$, which is, therefore, the space M of such lines, and a conformal self-dual 4-manifold. It can be naturally compactified within the space of analytic cycles of $`Z`$ to $`\overline{𝐌}=(E)\times (E)^{}`$, which is obviously a smooth manifold, but it carries no global conformal structure, as its canonical bundle has no square root. This means that the conformal structure on $`\overline{𝐌}`$ is smooth on M, and singular on $`=\overline{𝐌}𝐌`$. The cycles of $`Z`$ corresponding to a point $`\overline{x}=(A,a)`$ in this subset are pairs of complex projective lines in $`Z`$:
$$Z_{\overline{x}}=\{(A,l)Z=\}\{(L,a)Z=\}.$$
A $`\beta `$-surface in M, corresponding to a point $`\beta =(L,l)Z`$, is the set
$$\beta =\{(A,a)(E)\times (E)^{}|Al,La,AL,al\},$$
and can be naturally compactified to
$$\overline{\beta }=\{(A,l^\beta )\}\times \{(L^\beta ,a)\}^1\times ^1.$$
### 8.3. The tangent space to $``$
In order to describe the null-geodesics of M as 2-planes in $`Z`$, we study first the tangent space of $`Z=`$ at $`\beta =(L,l)`$:
A vector in $`T_{(L,l)}`$ is a pair of vectors $`(V,v)`$, with $`VT_L(E)`$ and $`vT_l(E)^{}`$, which satisfy a linear condition (as $`(E)\times (E)^{}`$). Actually, there is a duality between $`(E)^{}`$, the Grassmannian of 2-planes in $`E`$, and $`(E^{})`$, the projective space of $`E^{}:=\text{Hom}(E,)`$, and an analogous one between $`(E)`$ and $`(E^{})^{}`$:
$`(E)^{}l`$ $`\stackrel{}{}`$ $`l^o(E^{})`$
$`(E)L`$ $`\stackrel{}{}`$ $`L^o(E^{})^{}.`$
Then, the flag manifold $``$ is defined, as a submanifold of $`(E)\times (E)^{}`$, by the equation
$$y(Y)=0,yl^o,YL.$$
The vector $`VT_L(E)`$ is a homomorphism in $`\text{Hom}(L,E/L)`$. By duality, $`vv^0\text{Hom}(l^o,E^{}/l^o)`$. Then the vector $`(V,v)T_{(L,l)}(E)\times (E)^{}`$ lies in $``$ iff:
(33)
$$v^o(y^o)(Y)+y^o(V(Y))=0,YL,y^ol^o,$$
or, equivalently,
(34)
$$v|_L=\pi _lV,$$
where $`\pi _l:E/LE/l`$ is the projection (as $`Ll`$).
The geometry of $``$, as a subset of $`(E)\times (E)^{}`$, can be described in the following figure:
### 8.4. The 2-planes in $``$
Let us consider now a 2-plane $`F`$ in $`T_{(L,l)}`$, and the cycles (corresponding to points in $`\overline{𝐌}`$) tangent to it. We have three cases:
1. $`F=\overline{F}_\beta `$ is the “degenerate” 2-plane tangent to the 2 special curves $`\overline{Z}_L`$, $`\overline{Z}_l`$ whose union is the special cycle $`\overline{F}_(L,l)`$ corresponding to $`(L,l)\overline{𝐌}𝐌`$. There are no projective lines $`Z_x,x𝐌`$, tangent to it; only the special cycles $`\overline{Z}_{(L,a)},La`$ and $`\overline{Z}_{(A,l)},Al`$ are tangent to $`\overline{F}_{(L,l)}`$, actually only to the two privileged directions of $`\overline{Z}_L`$, resp. $`\overline{Z}_l`$.
Remark.
The special curves $`\overline{Z}_L`$, $`\overline{Z}_l`$ have trivial normal bundle, being fibers of the projections from $``$ to $`(E)`$, resp. $`(E)^{}`$, so these special curves form two complete families of analytic cycles in $``$, isomorphic to $`(E)`$, resp. $`(E)^{}`$. Two such curves are incident iff they are of different types ($`\overline{Z}_L`$ is of type $`E`$, $`\overline{Z}_l`$ is of type $`E^{}`$), so they can only form “polygons” with an even number of edges. But there are no quadrilaterals, as one can easily check, using the fact that $`\overline{Z}_L`$ and $`\overline{Z}_l`$ are incident iff $`Ll`$, thus iff $`l`$ is a line in $`(E)`$ containing $`L`$. But there are hexagons, corresponding to the 3 vertices and 3 sides of a triangle in $`(E)^2`$:
The hexagon above is not “flat”, i.e. there is no canonical submanifold of $``$ containing it. This, and the fact that there are no quadrilaterals made of $`\overline{Z}`$-type curves, is just a consequence of the fact that the distribution $`\overline{F}`$ on $`Z=`$ is non integrable; in fact it is the holomorphic contact structure induced by the Fubini-Study Einstein metric on $`^2`$, , see also Section 8.6.
2. $`F=F^a`$, for $`aL`$, $`al`$. This is a 2-plane that is tangent to only one of the special curves $`\overline{Z}_L`$. The projective lines tangent to $`F^a`$ at $`\beta =(L,l)`$ are $`Z_{(A,a)}`$, $`Al,AL`$, hence the corresponding null-geodesic is
(35)
$$\gamma ^a=\{(A,a)(E)\times (E)^{}|Al,Al\},$$
thus it is diffeomorphic to $``$, and its closure is
$$\overline{\gamma }^a=\{(A,a)(E)\times (E)^{}|Al\}^1.$$
Remark.
The “limit” curve is $`\overline{Z}_{(L,a)}`$, so it is non-singular at $`(L,l)`$. Actually, the points of $`Z_{(A,a)}`$ close to $`(L,l)`$ converge, when $`AL`$, to some points in $`\overline{Z}_L`$, which is tangent to $`F^a`$. We can, then, apply the same method as in Theorem 2 to conclude that the integral $`\alpha `$-cone associated to $`F^a`$ is a smooth manifold around $`(L,l)`$, thus, from Theorem 1, the Weyl tensor $`W^+`$ of M vanishes on the $`\beta `$-planes generated, along $`\gamma ^a`$, by its own direction. We will see that the vanishing of $`W^+`$ on these $`\alpha `$-planes leads to the existence of some $`\alpha `$-surfaces, see below. Of course, the deformation argument in Theorem 2 does not hold in the present case, as the normal bundle of $`\overline{Z}_L`$ is trivial, thus different from the one of the rest of the rational curves $`Z_{(A,a)}`$ (as we will see below, generic 2-planes through $`(L,l)`$ do not admit projective lines tangent to all their directions).
$`\mathrm{𝟐}^{}.`$ We have a similar situation for planes $`F=F^A`$$`Al`$, $`AL`$ — tangent to the other special curve $`\overline{Z}_l`$.
3. $`F=F^\phi `$, where $`\phi :(l)(L^o)`$ is a projective diffeomorphism such that $`\phi (L)=l^o`$. Indeed, the tangent spaces $`T_L(E)`$ and $`T_l(E)^{}`$ are isomorphic to $`\text{Hom}(L^o,E^{}/L^o)`$, resp. to $`\text{Hom}(l,E/l)`$, and a generic 2-plane $`F`$ in $`T_{(L,l)}`$ is the graph of a linear isomorphism $`\varphi :T_L(E)T_l(E)^{}`$ satisfying a linear condition (33) or (34). Actually, the graph is determined by the projective application $`\phi `$ induced by $`\varphi `$ from $`(T_L(E))(L^o)`$ to $`(T_l(E)^{})(l)`$:
The condition $`\phi (L)=l^o`$ is implied by (34). The null-geodesic associated to the 2-plane $`F^\phi `$ is
(36)
$$\gamma ^\phi =\{(A,a)(E)\times (E)^{}|Al,a^oL^oa^o=\phi (A)\},$$
and its closure in $`\overline{𝐌}`$ is
(37)
$$\overline{\gamma }^\phi =\{(A,a)(E)\times (E)^{}|Al,a^oL^o\},$$
hence the “limit” point is $`(L,l)\overline{𝐌}`$, corresponding to the special cycle $`\overline{Z}_{(L,l)}`$, none of whose components is tangent to $`F^\phi `$. The integral $`\alpha `$-cone associated to $`F^\phi `$ looks like suggested in the picture below:
### 8.5. The null-geodesics of the complexification of $`^2`$
The application $`\phi `$ has the following interpretation in terms of projective geometry on $`^2=(E)`$: a direction $`v`$ in $`T_l(E)^{}`$ is identified to the point $`\mathrm{ker}vAl/L(E)`$ and a direction $`VT_L(E)`$ is identified to a direction (thus a projective line $`a`$) through $`L(E)`$. $`\phi `$ is, thus, a homography that associates to $`Al`$ (we identify $`l`$ with the projective line $`l/L(E)`$) the line $`aL`$. As $`\phi (L)=l`$, we have, then, that three points $`(A,a),(B,b),(C,c)\beta ^{(L,l)}`$ belong to the same null-geodesic iff
(38)
$$(A,B:C,L)=(a,b:c,l),$$
i.e. the cross-ratio of the points $`A,B,C,Ll`$ equals the cross-ratio of the lines $`a,b,c,l`$ through $`L`$ (the dotted lines, and their intersections with the lines $`a,b,c`$ are the points in the integral $`\alpha `$-cone):
We can now describe the null-geodesics passing through a point $`(A,a)𝐌`$ and contained in a $`\beta `$-surface $`\beta ^(L,l)`$, whose closure $`\overline{\beta }`$ is isomorphic to $`^1\times ^1`$: they coincide with the rational curves in $`\overline{\beta }`$, containing $`(A,a)`$; except the “horizontal” ($`\overline{\gamma }^A`$) and “vertical” ($`\overline{\gamma }^a`$) ones, all these curves contain $`(L,l)`$:
We remark that, in the usual affine coordinates on $`\beta (^1\{L\})\times (^1\{l\})^2`$, these null-geodesics are the affine lines containing $`(A,a)`$, thus the projective structure on $`\beta `$ is (locally) isomorphic to a flat affine structure. We have seen, in Section 6 (Corollary 3), that this is true for all $`\beta `$-surfaces of a self-dual manifold.
### 8.6. The conformal structure of the complexification of $`^2`$
Let us study now the conformal structure of $`𝐌=(E)\times (E)^{}`$ directly; actually M has a complex metric $`g`$. Let $`(A,a)𝐌`$, then $`A`$ is transverse to $`a`$, thus we have the isomorphisms $`E/aA`$ and $`E/Aa`$. Then, a vector $`(V,v)T_{(A,a)}𝐌`$ is identified to a pair of homomorphisms $`V:Aa`$ and $`v:aA`$. Then the metric $`g`$ is given by
(39)
$$g((V,v),(W,w)):=\text{tr}(vW+wV),(V,v),(W,w)T_{(A,a)}𝐌.$$
Remark.
(The real part). Let $`h`$ be a hermitian metric on $`E`$. Then we have a real-analytic embedding of $`𝐌_0(E)`$ into M, given by:
$$(E)A(A,A^{})(E)\times (E)^{}.$$
A vector $`(V,v)T_{(A,A^{})}𝐌`$ is tangent to $`𝐌_0`$ iff
$$h(x,v(y))+h(V(x),y)=0,xA,yA^{}.$$
Then one easily checks that $`g((V,v),(W,w))=2h(V,W),(V,v),(W,w)T_{(A,A^{})}𝐌_0`$, hence, up to a constant, the restriction of $`g`$ to $`𝐌_0^2`$ is the Fubini-Study metric of $`^2S^5/S^1`$.
An isotropic vector in M is $`(V,v)T_{(A,a)}𝐌`$, such that $`vV=0`$, viewed as an endomorphism of $`A`$ (see above), or, equivalently, such that
(40)
$$dim(A+V(A)\mathrm{ker}v)>0.$$
Let us see which is the limit of the isotropic cone in the points of $``$: from the relation above, it follows that the isotropic cone in a point $`x`$ is
$$C_x=\{(0,v)T_x\}\{(V,0)T_x\},$$
so the conformal structure of M is singular at the “infinity” $``$.
Remark.
The situation $`(E)\times (E)^{}`$ is very similar to the one treated in Section 7, see also : $`(E)\times (E)^{}`$ has an Einstein self-dual metric $`g`$, singular at the “infinity”, and this Einstein structure yields a contact structure on the twistor space $`Z=`$; the field of 2-planes determined by this contact structure corresponds to the “infinity” $`(E)\times (E)^{}`$. But these planes do not admit tangent rational curves, with normal bundle $`𝒪(1)𝒪(1)`$: the conformal structure does not extend to the “infinity” (which is, therefore, not a conformal infinity).
### 8.7. $`\alpha `$-planes and $`\beta `$-planes
We consider the isotropic planes in $`T_{(A,a)}𝐌`$ ($`Aa`$): For a fixed isotropic direction, represented by a generic vector $`(V,v)T_{(A,a)}𝐌`$, the line $`\mathrm{ker}va`$ and the plane $`V(A)+AA`$ are fixed. The linear space of all vectors $`(W,v)T_{(A,a)}𝐌`$ satisfying
$$W(A)A+V(A);w|_{\mathrm{ker}v}=0$$
is isotropic and orthogonal to $`(V,v)`$: they form a $`\beta `$-plane. The $`\alpha `$-plane $`F^\alpha `$ containing $`(V,v)`$ corresponds to the isotropic vectors $`(W,w)`$, orthogonal to $`(V,v)`$, with $`\mathrm{ker}w\mathrm{ker}v`$. As a plane transverse to all the $`\beta `$-planes (whose projection onto $`T_A(E)`$ or $`T_a(E)^{}`$ is never injective), $`F^\alpha `$ is determined by a linear isomorphism $`\phi :T_A(E)T_a(E)^{}`$, whose graph in $`T_{(A,a)}(E)\times (E)^{}`$ is $`F^\alpha `$; $`\phi `$ induces the application $`\phi :(a)(E/A)`$ between the projective spaces of $`T_A(E)`$, resp. $`T_a(E)^{}`$. The plane $`F^\alpha =F^\phi `$, the graph of $`\phi `$, is isotropic iff $`V\phi (V)`$, $`V(a)`$, i.e. $`\phi `$ is the homography that sends a point $`X`$ in $`a`$ into the projective line through $`A`$ and $`X`$. We can extend $`\phi `$ to a projective isomorphism $`\phi ^{}:(T_A(E))(T_a(E)^{})`$: for example, $`(T_A(E))`$ contains $`T_A(E)`$ as an affine open set. Then $`\phi ^{}`$ is defined as follows:
$`\phi ^{}|_{T_A(E)}`$ $`:=`$ $`\phi `$
$`\phi ^{}|_{(T_A(E))}`$ $`:=`$ $`\phi .`$
Actually $`(T_A(E))(E)`$ and $`(T_a(E)^{})(E)^{}`$. We have then:
###### Proposition 13.
A generic $`\alpha `$-plane $`F^\alpha =F^\phi `$ in $`T_{(A,a)}𝐌`$ is the graph of a linear isomorphism $`\phi :T_A(E)T_a(E)^{}`$, which is determined by a projective isomorphism
$$\phi ^{}:(E)(E)^{},$$
such that $`\phi ^{}(A)=a`$ and $`\phi ^{}(l)=la`$, for all $`lA`$.
### 8.8. Exponentials of $`\alpha `$-planes
The exponential $`\mathrm{exp}(F^\phi )`$ has an interpretation in terms of projective geometry: Each direction $`(V,v)F^\phi `$ is determined by the point $`\mathrm{ker}v`$ in $`a(E)`$ and the line through $`A`$ and $`\mathrm{ker}v`$, and a homography $`\varphi ^{(V,v)}`$ from the points $`B`$ of the projective line $`A+\mathrm{ker}v`$ to the space of lines $`b`$ through $`\mathrm{ker}v`$ (see next picture and the convention below). As this homography is the restriction of $`\phi ^{}`$ to the appropriate spaces, it follows that it is related to the homography $`\varphi ^{(W,w)}`$, where $`(W,w)`$ is another direction in $`F^\phi `$: the points $`D:=bc`$, $`P:=a(B+C)`$ and $`A`$ are collinear:
Of course, this implies that $`P`$ determines a homography $`\psi ^P`$ between the lines $`A+\mathrm{ker}v`$ and $`A+\mathrm{ker}w`$, such that $`\psi ^P(A)=A`$ and $`\psi ^P(\mathrm{ker}v)=\mathrm{ker}w`$. Then, for any other points $`B^{}(A+\mathrm{ker}v),C^{}=\psi ^P(B)(A+\mathrm{ker}w)`$, the lines $`b^{}=\varphi ^{(V,v)}(B^{}),c^{}=\varphi ^{(W,w)}(C^{})`$ intersect on the line $`(A+P)`$ (see the right hand side of the picture above).
Convention. In the framework of plane projective geometry, we identify a point in $`(E)^{}`$ with a line in $`(E)`$ (we note, for example $`\mathrm{ker}va`$). The lines determined by the distinct points $`B`$ and $`C`$ will be denoted by $`(B+C)`$ (thus $`B,C(B+C)`$).
The null-geodesic tangent to $`(V,v)`$ at $`(A,a)`$ is the set $`\{(B,b)|B(A+\mathrm{ker}v),b=\varphi ^{(V,v)}(B)\}`$, and the null-geodesic tangent to $`(W,w)`$ is the analogous set of the pairs $`(C,c)`$. Thus
$`\mathrm{exp}_{(A,a)}(F^\phi )=\mathrm{exp}_{(A,a)}(F^\alpha )=\{(C,c)|C(E),CA,`$
$`c=((C+A)a)+((A+P)b^C)\}\{(A,a)\},`$
where $`B^C:=(A+\mathrm{ker}v)(P+C)`$, and $`b^C:=\varphi ^{(V,v)}(B^C)`$, as in the picture above (where $`B=B^C`$, $`b=b^C`$). This gives the exponential of the $`\alpha `$-plane determined by the isotropic vector $`(V,v)`$. We remark that the point $`(A,a)`$ has a privileged position in $`\mathrm{exp}_{(A,a)}(F^\alpha )`$: $`(ab)(A+B)`$ $`(B,b)\mathrm{exp}_{(A,a)}(F^\alpha )`$; on the other hand $`(bc)(B+C)`$ in general (see the picture above), which means that the points $`(B,b)`$ and $`(C,c)`$ are not null-separated (i.e. they do not belong to the same null-geodesic). That means that $`\mathrm{exp}_{(A,a)}(F^\alpha )`$ is not totally isotropic, thus there is no $`\alpha `$-surface tangent to a generic $`\alpha `$-plane; not surprising as the corresponding $`\alpha `$-cone is not flat (see Section 8.4).
But there are $`\alpha `$-surfaces tangent to the two $`\alpha `$-planes $`\{(V,0)|VT_A(E)\}`$ and $`\{(0,v)|vT_a(E)^{}\}`$: the “slices” $`\{A\}\times (E)^{}`$ and $`(E)\times \{a\}`$. (It is easy to see that these planes are isotropic, and that they are not $`\beta `$-planes, as these project on lines in $`T_A(E)`$, resp. $`T_a(E)^{}`$.)
Thus $`𝐌=(E)\times (E)^{}`$ is a conformal self-dual manifold, not anti-self-dual, that admits $`\alpha `$-surfaces passing through any point.
Centre de Mathématiques
UMR 7640 CNRS
Ecole Polytechnique
91128 Palaiseau cedex
France
e-mail: belgun@math.polytechnique.fr |
warning/0002/math0002139.html | ar5iv | text | # The distribution of spacings between the fractional parts of 𝑛²𝛼
## 1. Introduction
Fix an irrational number $`\alpha `$. The problem of the distribution of the local spacings between the members of the sequence $`n^2\alpha mod1`$, $`1nN,`$ has received attention recently (see ). It arises for example in the study of the local spacing distributions between the eigenvalues of special Hamiltonians. We order the above numbers in $`[0,1)`$ as
(1.1)
$$0\beta _1\beta _2\mathrm{}\beta _N<1$$
and set $`\beta _{N+j}=\beta _j`$. The $`k`$-th consecutive spacing measure is defined to be the probability measure on $`[0,\mathrm{})`$ given by
(1.2)
$$\mu _k(N,\alpha ):=\frac{1}{N}\underset{j=1}{\overset{N}{}}\delta _{N(\beta _{j+k}\beta _j)}$$
where $`\delta _x`$ is a unit delta mass at $`x`$. The problem is to understand the behavior of these measures as $`N\mathrm{}`$ and in particular their dependence on the diophantine approximations to $`\alpha `$.
We say $`\alpha `$ is of type $`K`$ if there is $`c_\alpha >0`$ such that $`|\alpha \frac{a}{q}|c_\alpha q^K`$ for all relatively prime integers $`a`$ and $`q`$. It is easy to see that if $`\alpha `$ is not of type 3 then there is a subsequence $`N_j\mathrm{}`$ for which the measures $`\mu _k(N_j,\alpha )`$ converge to a measure supported on $`𝐍=\{0,1,2,\mathrm{}\}`$. On the other hand numerical experiments indicate that for $`\alpha =\sqrt{2}`$ these $`k`$-th consecutive spacings behave like what one typically gets for spacings when placing $`N`$ numbers in $`[0,1]`$ uniformly and independently at random . That is $`\mu _k(N,\sqrt{2})`$ appears to converge to $`\mu _k:=\frac{x^k}{k!}e^xdx`$.
A standard approach to the analysis of the consecutive spacing is via local $`m`$-level correlations. These are defined as follows: As test functions we use functions $`f(x_1,x_2,\mathrm{},x_m)`$ which are symmetric in $`(x_1,x_2,\mathrm{},x_m)`$ and which are functions of the difference of the coordinates, that is $`f(x+(t,t,\mathrm{},t))=f(x)`$ for all $`t𝐑`$. We assume further that $`f`$ is local, that is it is compactly supported modulo the diagonal. We will call these admissible test functions. Define the correlations
(1.3)
$$R^{(m)}(N,\alpha ,f):=\frac{1}{N}\underset{1j_1<\mathrm{}<j_mN}{}f(N(\beta _{j_1},\mathrm{},\beta _{j_m})).$$
Note that $`R^{(m)}`$ is not a probability density and it may well tend to infinity as $`N\mathrm{}`$ (think of the case when $`\alpha `$ is rational). In the case that the $`\beta `$’s in (1.1) come from a random choice of points in $`[0,1)`$ these correlations satisfy
(1.4)
$$R^{(m)}(N,f)_{0x_2\mathrm{}x_m}f(0,x_2,\mathrm{},x_m)𝑑x_2\mathrm{}𝑑x_m.$$
We say that $`n^2\alpha mod1`$, $`nN`$, is Poissonian if for all $`m2`$ and $`f`$ as above
(1.5)
$$R^{(m)}(N,\alpha ,f)_{0x_2\mathrm{}x_m}f(0,x_2,\mathrm{},x_m)𝑑x_2\mathrm{}𝑑x_m$$
as $`N\mathrm{}`$.
As with the method of moments in convergence of measures, if the $`m`$-level correlations are Poissonian then the consecutive spacing measures $`\mu _k(N,\alpha )`$ converge to $`\mu _k`$. Thus Poissonian in the sense of (1.5) (i.e. for correlations) implies that as far as local spacings go, the numbers behave randomly. We will also consider cases where (1.5) holds along a subsequence $`N_j\mathrm{}`$, in such a case we say that $`n^2\alpha mod1`$, $`1nN`$ is Poissonian on a subsequence.
The results below lead us to the following
Conjecture: If $`\alpha `$ is of type $`2+ϵ`$ for every $`ϵ>0`$ and the convergents $`\frac{a}{q}`$ to $`\alpha `$ satisfy $`lim_q\mathrm{}\frac{\mathrm{log}\stackrel{~}{q}}{\mathrm{log}q}=1`$, where $`\stackrel{~}{q}`$ is the square free part of $`q`$, then $`n^2\alpha mod1`$ is Poissonian.
We note that almost all $`\alpha `$ (for Lebesgue measure) satisfy the hypothesis in the Conjecture and that assuming some standard conjectures in diophantine analysis any real algebraic irrationality satisfies these hypotheses (see Appendix A).
Unfortunately the methods of this paper appear not to be powerful enough to prove anything for the numbers $`\alpha `$ in the Conjecture. They require that $`\alpha `$ have somewhat better approximations by rationals. One of our main results gives conditions on the diophantine approximations to $`\alpha `$ which ensure that $`n^2\alpha mod1`$ is Poissonian along a subsequence. In particular this allows us to conclude that for the topologically generic $`\alpha `$ (i.e. in the sense of Baire) $`n^2\alpha mod1`$ is Poissonian along a subsequence. On the other hand the naive expectation that for any irrational $`\alpha `$, $`n^2\alpha mod1`$ is Poissonian along a subsequence, fails dramatically. The source of this phenomenon is large square factors in the denominator of the convergents to $`\alpha `$. We will exhibit an $`\alpha `$ for which the $`5`$-level correlations go to infinity as $`N\mathrm{}`$. We also provide an $`\alpha `$ of type less than three and a sequence of integers $`\{N_j\}_{j=1}^{\mathrm{}}`$ along which the 5-level correlations diverge to infinity.
The precise statements are as follows:
###### Theorem 1.
Let $`\alpha 𝐑`$. Suppose there are infinitely many rationals $`b_j/q_j`$, with $`q_j`$ prime, satisfying
$$\left|\alpha \frac{b_j}{q_j}\right|<\frac{1}{q_j^3}.$$
Then there is a subsequence $`N_j\mathrm{}`$ with $`\frac{\mathrm{log}N_j}{\mathrm{log}q_j}1`$ for which (1.5) holds for all $`m2`$ and all $`f`$. That is to say $`n^2\alpha mod1`$ is Poissonian along this subsequence.
With a lot more work concerning the exponential sums discussed in Section 9, for general moduli $`q`$, we can relax the condition that $`q_j`$ be prime in Theorem 1. In fact we can prove the following (we do not go into the proof in this paper) which shows that for such approximants the size of the square free parts $`\stackrel{~}{q}_j`$ of $`q_j`$ is decisive.
Theorem 1’: Let $`\alpha `$ be an irrational for which there are infinitely many rationals $`b_j/q_j`$ satisfying
$$\left|\alpha \frac{b_j}{q_j}\right|<\frac{1}{q_j^3}.$$
Then the following are equivalent :
(i) There is a subsequence $`N_j\mathrm{}`$ with $`\frac{\mathrm{log}N_j}{\mathrm{log}q_j}1`$ such that $`n^2\alpha mod1`$ is Poissonian along $`N_j`$.
(ii)
$$\underset{j\mathrm{}}{lim}\frac{\mathrm{log}\stackrel{~}{q}_j}{\mathrm{log}q_j}=1.$$
As to the divergence of correlations we have:
###### Theorem 2.
(a) There is an irrational $`\alpha `$ and a test function $`f`$ such that
$$R^{(5)}(N,\alpha ,f)N^{\frac{1}{4}}(\mathrm{log}N)^5,\text{ as }N\mathrm{}.$$
(b) For every $`\sigma >23/8`$ there is an $`\alpha `$ of type $`\sigma `$ and a test function $`f`$ such that
$$\overline{\underset{N\mathrm{}}{lim}}R^{(5)}(N,\alpha ,f)=\mathrm{}.$$
The test functions $`f`$ in Theorem 2 are nonnegative and are supported in a neighborhood of $`0`$ (modulo the diagonal) and the source of the divergence is that there are zero density, but non-negligible, clusters among the numbers $`n^2\alpha mod1`$ , $`nN`$.
We note that these clusters which spoil the correlations do not have the same effect on the probability measures $`\mu _k(N,\alpha )`$. So it is quite possible for example that the $`\alpha `$ in part (b) of Theorem 2 has its $`\mu _k(N,\alpha )`$ measures converge to the Poissonian $`\mu _k`$. We have chosen in this paper to call $`n^2\alpha mod1`$ Poissonian if the strongest behavior holds - that is the correlations are Poissonian.
The proofs of the above theorems are based on the following closely related diophantine problem: Consider the spacing distributions (normalized to mean spacing $`1`$ as before) of the numbers $`\{bn^2/q\}`$, $`nN`$, or what is the same, the spacing distribution of the integers
(1.6)
$$n^2bmodq,1nN.$$
Here $`q`$ is prime ($`q\mathrm{}`$), $`b`$ is any number not divisible by $`q`$ and $`N`$ is in the range $`[q^{1/2+ϵ},\frac{q}{\mathrm{log}q}]`$ for some $`ϵ>0`$. The reason for this range for $`N`$ is that if $`N\sqrt{q}`$ and say $`b=1`$ then the spacing distributions may be easily determined (since $`n^2<q`$ for $`nN\sqrt{q}`$) and are certainly nonrandom. Similarly if $`N=q`$ then the sequence in (1.6) consists of all the quadratic residues (or non-residues) and hence the spacings are integers and so cannot follow a Poissonian law. In fact the limiting spacing distributions of $`\mu _k(q,q,b)`$ were determined by Davenport . So it is only in the range $`N[q^{1/2+ϵ},\frac{q}{\mathrm{log}q}]`$ that we can hope for randomization. The following Theorem shows that indeed, to a certain extent, this is the case.
Let $`R^{(m)}(N,b/q,f)`$ denote the scaled $`m`$-level correlations for the sequence (1.6).
###### Theorem 3.
Fix $`m2`$ , $`f`$ and $`\delta >0`$. Then as $`q\mathrm{}`$, $`q`$ prime
$$R^{(m)}(N,b/q,f)_{0x_2\mathrm{}x_m}f(0,x_2,\mathrm{},x_m)𝑑x_2\mathrm{}𝑑x_m$$
uniformly for $`(b,q)=1`$ and $`q^{1\frac{1}{2m}+\delta }N\frac{q}{\mathrm{log}q}`$.
A crucial ingredient in our proof of Theorem 3 is the Riemann Hypothesis for curves (of arbitrary large genus) over finite fields (Weil ).
In the range in which Theorem 3 applies it gives Poisson statistics and Theorem 3 easily yields Theorem 1. For $`m3`$ it is not possible to extend the range of $`N`$ in Theorem 3 much further. The reason is related to the previous divergence of correlations phenomenon. For suitable $`b`$ (depending on $`q`$) there will be large clusters among the numbers $`n^2b(modq)`$ , $`nN`$. This is highlighted by the following Theorem.
###### Theorem 4.
Fix $`m3`$ and $`\delta >0`$. Then there is a test function $`f`$ such that for $`q^{\frac{1}{2}}Nq^{\frac{m}{m+2}\delta }`$,
$$\underset{q\mathrm{}}{lim}\underset{(b,q)=1}{\mathrm{max}}R^{(m)}(N,b/q,f)=\mathrm{}.$$
Acknowledgment: We would like to thank E. Bombieri for his help with the application of the ABC conjecture described in Appendix A. We also thank the referee for suggesting a stronger version of Theorem 2 (a) with a simpler proof than our original one.
## 2. A comparison lemma
We will need to deal with the following situation: We are given two families of sequences $`𝒩=\{x_N(n):nN\}`$ and $`𝒩^{}=\{x_N^{}(n):NN\}`$ in $`[0,1)`$ and we wish to compare the limiting correlation functions of these two families, seeking to show that if the correlations exist for one sequence then they exist for the other, or they diverge for one if they do for the other. We show that it can be done if the two sequences are close in a suitable sense. We define the scaled distance between the sequences to be
$$ϵ(𝒩,𝒩^{}):=N\underset{nN}{\mathrm{max}}|x_N(n)x_N^{}(n)|.$$
A general method for carrying out the comparison is formalized in the following:
###### Lemma 5 (Comparison Lemma).
Assume that $`𝒩,𝒩^{}[0,1)`$ are two families of sequences with $`ϵ(𝒩,𝒩^{})0`$ as $`N\mathrm{}`$. Then for all smooth test functions $`f`$, we have
$$\left|R^{(k)}(𝒩,f)R^{(k)}(𝒩^{},f)\right|R^{(k)}(𝒩,f_+)ϵ(𝒩,𝒩^{})$$
for $`N`$ sufficiently large, where $`f_+0`$ is a smooth admissible test function (depending only on $`f`$).
###### Proof.
For notational simplicity, we will do the case of pair correlation ($`k=2`$). Our test function $`f0`$ can then be written as $`f(x_1,x_2)=g(x_1x_2)`$ for some $`gC_c^{\mathrm{}}(𝐑)`$, say $`g`$ supported inside $`[\rho ,\rho ]`$. Let $`g_+0`$ be smooth, compactly supported and such that $`g_+`$ is constant on $`[2\rho ,2\rho ]`$, where it equals $`\mathrm{max}|g^{}|`$. Set $`f_+(x_1,x_2):=2g_+(x_1x_2)`$. For further notational simplicity also set $`\delta _{m,n}:=x_N(m)x_N(n)`$ and $`\delta _{m,n}^{}:=x_N^{}(m)x_N^{}(n)`$.
By the mean value theorem we have
$$\begin{array}{cc}\hfill R^{(2)}(𝒩,f)R^{(2)}(𝒩^{},f)& =\frac{1}{N}\underset{1m<nN}{}g(N\delta _{m,n})g(N\delta _{m,n}^{})\hfill \\ & =\frac{1}{N}\underset{1m<nN}{}g^{}(N\xi _{m,n})N(\delta _{m,n}\delta _{m,n}^{})\hfill \end{array}$$
where $`\xi _{m,n}`$ lies between $`\delta _{m,n}`$ and $`\delta _{m,n}^{}`$.
For the difference $`R^{(2)}(f,𝒩)R^{(2)}(f,𝒩^{})`$ to contain a nonzero contribution from the term indexed by the pair $`(m,n)`$, we must have at least one of $`N\delta _{m,n}`$ or $`N\delta _{m,n}^{}`$ lying in $`\mathrm{supp}g[\rho ,\rho ]`$. Now $`N\xi _{m,n}`$ is within $`2ϵ(𝒩,𝒩^{})`$ of both $`N\delta _{m,n}`$ and $`N\delta _{m,n}^{}`$, which implies that both lie in $`[2\rho ,2\rho ]`$, as does $`\xi _{m,n}`$ if $`N`$ is sufficiently large so that $`2ϵ(𝒩,𝒩^{})<\rho `$. Since $`g_+`$ is constant on $`[2\rho ,2\rho ]`$ we find that $`g_+(\xi _{m,n})=g_+(N\delta _{m,n})`$.
Thus we get
$$\begin{array}{cc}\hfill |R^{(2)}(𝒩,f)R^{(2)}(𝒩^{},f)|& \frac{1}{N}\underset{1m<nN}{}g_+(N\delta _{m,n})2ϵ(𝒩,𝒩^{})\hfill \\ & =R^{(2)}(𝒩,f_+)ϵ(𝒩,𝒩^{})\hfill \end{array}$$
as required. ∎
## 3. Derivation of Theorem 1
As an immediate application of the comparison lemma, we derive Theorem 1 from Theorem 3.
Fix $`\alpha `$. Suppose there are infinitely many rationals $`b_j/q_j`$ with $`q_j`$ prime, satisfying
$$|\alpha \frac{b_j}{q_j}|<\frac{1}{q_j^3}.$$
We let $`N_j=[\frac{q_j}{\mathrm{log}q_j}]`$, where \[.\] denotes the integer part function. Fix an $`m2`$ and a test function $`f`$ as above. We need to show that
$$\underset{j\mathrm{}}{lim}R^{(m)}(N_j,\alpha ,f)=_{0x_2\mathrm{}x_m}f(0,x_2,\mathrm{},x_m)𝑑x_2\mathrm{}𝑑x_m.$$
By Theorem 3 applied to $`q=q_j,b=b_j`$ and $`N=N_j`$ we know that
$$\underset{j\mathrm{}}{lim}R^{(m)}(N_j,b_j/q_j,f)=_{0x_2\mathrm{}x_m}f(0,x_2,\mathrm{},x_m)𝑑x_2\mathrm{}𝑑x_m$$
We use the comparison principle to estimate the difference:
$$|R^{(m)}(N_j,\alpha ,f)R^{(m)}(N_j,b_j/q_j,f)|.$$
Take $`𝒩_j^{}=\{\{\alpha n^2\}:nN_j\}`$ and $`𝒩_j=\{\{b_jn^2/q_j\}:nN_j\}`$. By lemma 5,
(3.1)
$$|R^{(m)}(𝒩_j^{},f)R^{(m)}(𝒩_j,f)|R^{(m)}(𝒩_j,f_+)ϵ(𝒩_j,𝒩_j^{})$$
for some admissible test function $`f_+0`$. We have
$$|\{\alpha n^2\}\{\frac{b_j}{q_j}n^2\}||\alpha \frac{b_j}{q_j}|n^2\frac{N_j^2}{q_j^3}\frac{1}{N_j(\mathrm{log}N_j)^3}$$
and thus
$$ϵ(𝒩_j,𝒩_j^{})=N_j\underset{nN_j}{\mathrm{max}}|\{\alpha n^2\}\{\frac{b_j}{q_j}n^2\}|\frac{1}{(\mathrm{log}N_j)^3}0.$$
By Theorem 3, $`R^{(m)}(𝒩_j,f_+)`$ is bounded (it converges as $`j\mathrm{}`$). Thus we use (3.1) to deduce that
$$|R^{(m)}(N_j,\alpha ,f)R^{(m)}(N_j,b_j/q_j,f)|0$$
which gives Theorem 1. ∎
## 4. A divergence principle
We present a mechanism that ensure divergence of high correlations of the sequence $`\{bn^2/q\}`$: The presence of larges square factors in $`q`$.
###### Lemma 6.
Let $`q=uv^2`$ with $`v>q^\delta `$ for some $`\delta >0`$, let $`\eta >1\delta `$ and suppose that $`\mathrm{log}N/\mathrm{log}q>\eta `$. Let $`f0`$ be a positive admissible test function which is non-vanishing at the origin. Then for all $`b`$,
(4.1)
$$R^{(m)}(N,\frac{b}{q},f)\frac{1}{N}f(0)(\frac{Nv}{q})^m.$$
In particular $`R^{(m)}(N,b/q,f)`$ will diverge to infinity for $`m`$ sufficiently large in terms of $`\delta `$ and $`\eta `$.
###### Proof.
Write $`f(x_1,\mathrm{},x_m)=g(x_1x_2,\mathrm{}x_{m1}x_m)`$ for $`gC_c(𝐑^{m1})`$, $`g0`$, $`g(0)0`$. Then
$$R^{(m)}(N,\frac{b}{q},f)=\frac{1}{N}\underset{1n_1<\mathrm{}<n_mN}{}g(\mathrm{},N\{\frac{bn_j^2}{q}\}N\{\frac{bn_{j+1}^2}{q}\},\mathrm{})$$
Since $`g0`$, we may count only the contribution of those $`(n_1,\mathrm{},n_m)`$ ($`n_j`$ distinct) for which all the components $`n_1,\mathrm{},n_m`$ are divisible by $`uv`$. There are $`[N/uv]^m=[Nv/q]^m`$ such $`m`$-tuples. If $`n=uvn^{}`$ then since $`q=uv^2`$ we have
$$\{\frac{bn^2}{q}\}=\{bu(n^{})^2\}=0$$
and so we find
$$R^{(m)}(N,\frac{b}{q},f)\frac{1}{N}f(0)(\frac{Nv}{q})^m.$$
Since $`v>q^\delta `$ and $`Nq^\eta `$ with $`\eta >1\delta `$, this gives $`R^{(m)}(N,\frac{b}{q},f)q^s`$ with
$$s=\eta (m1)+m\delta m=m(\eta (1\delta ))\eta $$
which is positive if $`m>\eta /(\eta (1\delta ))>0`$. Thus for $`m`$ sufficiently large, $`R^{(m)}(N,b/q,f)`$ will diverge in these ranges. ∎
## 5. Proof of Theorem 4
Fix $`m2`$, some small $`\delta >0`$, and let $`N`$, $`q`$ be large such that $`q^{1/3}Nq^{m/(m+2)\delta }`$. Let $`f0`$ is an admissible test function, $`f(0)0`$, and $`f_+f`$ the smooth majorant appearing in Lemma 5. We want to show that there exists $`b<q`$ coprime to $`q`$ such that $`R^{(m)}(N,b/q,f_+)`$ is large.
We first produce $`q^{}`$ which is a square, $`q^{}=v^2`$, coprime to $`q`$, such that
(5.1)
$$qq^{}N^3(\mathrm{log}N)^3.$$
To do so, find $`v`$ in the interval
$$J=[\sqrt{\frac{N^3\mathrm{log}^3N}{q}},2\sqrt{\frac{N^3\mathrm{log}^3N}{q}}]$$
which is coprime to $`q`$. Note that $`N^3\mathrm{log}^3N/q(\mathrm{log}q)^3`$ since $`Nq^{1/3}`$ and so the existence of such numbers $`v`$ is assured for any $`q`$ sufficiently large. Indeed, if $`q`$ is sufficiently large then in any interval $`[x,2x]`$ with $`x(\mathrm{log}q)^{3/2}`$ there is a prime $`\mathrm{}`$ not dividing $`q`$, since otherwise $`q`$ would be divisible by all primes in the interval and consequently $`\mathrm{log}q`$ would be at least as large as $`_{xp2x}\mathrm{log}px`$ which contradicts $`x(\mathrm{log}q)^{3/2}`$.
We now put $`q^{}=v^2`$. Thus $`(q^{},q)=1`$ and (5.1) holds. Because $`q^{}=v^2`$ is a square, we may use the divergence principle (4.1) to see that for all $`b^{}`$ we have
$$R^{(m)}(N,\frac{b^{}}{q^{}},f)\frac{1}{N}(\frac{N}{v})^m$$
Since $`v=\sqrt{q^{}}\sqrt{N^3\mathrm{log}^3N/q}`$, we find that for all $`b^{}`$
$$R^{(m)}(N,\frac{b^{}}{q^{}},f)\frac{N^{m1}q^{m/2}}{N^{3m/2}(\mathrm{log}N)^{3m/2}}=\frac{q^{m/2}}{N^{m/2+1}(\mathrm{log}N)^{3m/2}}.$$
Now use $`q^{1/3}Nq^{m/(m+2)\delta }`$ to find that for some $`C>0`$,
(5.2)
$$R^{(m)}(N,\frac{b^{}}{q^{}},f)Cq^{\delta (m/2+1)}(\mathrm{log}q)^{3m/2}$$
uniformly in $`b^{}`$ if $`q>q_0`$. Since $`\delta >0`$, this diverges with $`q`$.
Because $`q`$, $`q^{}`$ are coprime, there are $`0<b<q`$, $`0<b^{}<q^{}`$ so that $`bq^{}b^{}q=1`$ and so
$$|\frac{b}{q}\frac{b^{}}{q^{}}|=\frac{1}{qq^{}}\frac{1}{N^3\mathrm{log}^3N}.$$
By the comparison principle (lemma 5), the two sequences $`𝒩=\{\{bn^2/q\}:nN\}`$ and $`𝒩^{}=\{\{b^{}n^2/q^{}\}:nN\}`$ satisfy
(5.3)
$$|R^{(m)}(N,\frac{b}{q},f)R^{(m)}(N,\frac{b^{}}{q^{}},f)|ϵ(𝒩,𝒩^{})R^{(m)}(N,\frac{b}{q},f_+)$$
where $`f_+`$ is a majorant for $`f`$, and in particular nonvanishing at the origin. Moreover
(5.4)
$$\begin{array}{cc}\hfill ϵ(𝒩,𝒩^{})& =N\underset{nN}{\mathrm{max}}|\{\frac{bn^2}{q}\}\{\frac{b^{}n^2}{q^{}}\}||\frac{b}{q}\frac{b^{}}{q^{}}|N^3\hfill \\ & \frac{1}{(\mathrm{log}N)^3}\frac{1}{(\mathrm{log}q)^3}\hfill \end{array}$$
We claim that
$$R^{(m)}(N,\frac{b^{}}{q^{}},f_+)\frac{C}{3}q^{\delta (m/2+1)}(\mathrm{log}q)^{3m/2}$$
Indeed, assuming otherwise we have from (5.3) and (5.4) that
$$|R^{(m)}(N,\frac{b}{q},f)R^{(m)}(N,\frac{b^{}}{q^{}},f)|=o(q^{\delta (m/2+1)}(\mathrm{log}q)^{3m/2})$$
which together with (5.2) forces $`R^{(m)}(N,b/q,f)>\frac{C}{3}q^{\delta (m/2+1)}(\mathrm{log}q)^{3m/2}`$. However, since $`f_+f0`$ we find that
$$R^{(m)}(N,\frac{b}{q},f_+)R^{(m)}(N,\frac{b}{q},f)>\frac{C}{3}q^{\delta (m/2+1)}(\mathrm{log}q)^{3m/2}$$
contradicting our assumption. ∎
## 6. Preliminaries on continued fractions
We recall the standard notions of the theory of continued fractions (see e.g. ).
Given integers $`a_0𝐙`$, $`a_1,a_2,\mathrm{}1`$, one defines integers $`p_m`$, $`q_m`$ by the recursion ($`m1`$):
$$\begin{array}{cc}\hfill p_m& =a_mp_{m1}+p_{m2}\hfill \\ \hfill q_m& =a_mq_{m1}+q_{m2}\hfill \end{array}$$
with $`p_1=1`$, $`p_0=a_0`$, $`q_1=0`$, $`q_0=1`$. These satisfy the relations
$$p_mq_{m1}p_{m1}q_m=(1)^{m1}$$
and
$$p_mq_{m2}p_{m2}q_m=(1)^ma_m.$$
The finite continued fraction
$$[a_0;a_1,\mathrm{},a_m]:=a_0+\frac{1}{a_1+{\displaystyle \frac{1}{a_2+{\displaystyle \frac{1}{\mathrm{}+{\displaystyle \frac{1}{a_m}}}}}}}$$
is then $`p_m/q_m`$.
The infinite simple continued fraction $`[a_0;a_1,a_2,\mathrm{}]`$ is the limit of the “convergents” $`p_m/q_m`$. Every irrational $`\alpha `$ has a unique continued fraction expansion.
The convergents give very good rational approximations to $`\alpha `$: We have
$$\frac{1}{2}\frac{1}{q_mq_{m+1}}<|\alpha \frac{p_m}{q_m}|<\frac{1}{q_mq_{m+1}}.$$
The convergents $`p_m/q_m`$ are the “best” rational approximations to $`\alpha `$, in the following senses: If $`p/q`$ satisfies $`|\alpha p/q|<1/2q^2`$ then $`p/q=p_m/q_m`$ for some $`m`$. Moreover, for $`m>1`$, if $`0<qq_m`$ and $`p/qp_m/q_m`$ then $`|\alpha p/q|>|\alpha p_m/q_m|`$.
## 7. Proof of Theorem 2(a)
### 7.1. Constructing $`\alpha `$
We want to find an irrational $`\alpha `$ such that
(7.1)
$$R^{(5)}(\alpha ,N)\frac{N^{1/4}}{(\mathrm{log}N)^5}.$$
The construction below is due to the referee, who strengthened and considerably simplified our original argument.
We construct $`\alpha `$ by means of its continued fraction expansion, by inductively finding $`a_0,a_1,\mathrm{},a_m`$ so that the denominators $`q_m`$ of the convergents are squares: $`q_m=v_m^2`$.
To do so, define pairs of integers $`(r_m,v_m)`$ by $`r_1=v_1=0`$, $`r_0=v_0=1`$, $`r_1=v_1=1`$ and for $`m1`$
(7.2)
$$v_{m+1}=r_mv_m^2+v_{m1},r_{m+1}=[\mathrm{log}v_{m+1}]$$
Now set $`a_0=1`$, and for $`m0`$
$$a_{m+1}=r_m^2v_m^2+2r_mv_{m1}$$
Let $`\alpha =[a_0;a_1,a_2,\mathrm{}]=[1;1,3,6,\mathrm{}]`$.
We claim that the denominator $`q_m`$ of convergent to $`\alpha `$ equals $`v_m^2`$. To see this, use induction: By the recursion for the convergents, $`q_{m+1}=a_{m+1}q_m+q_{m1}`$ and by induction
$$\begin{array}{cc}\hfill q_{m+1}& =a_{m+1}v_m^2+v_{m1}^2\hfill \\ & =(r_m^2v_m^2+2r_mv_{m1})v_m^2+v_{m1}^2\hfill \\ & =(r_mv_m^2+v_{m1})^2=v_{m+1}^2\hfill \end{array}$$
as required.
Note also that from the recursion (7.2),
$$q_{m+1}r_m^2q_m^2q_m^2(\mathrm{log}q_m)^2$$
Thus $`\alpha `$ is of type $`3+ϵ`$, for all $`ϵ>0`$.
Now we want to show that $`R^{(5)}(\alpha ,N)N^{1/4}/(\mathrm{log}N)^5`$. Pick $`m`$ so that $`q_mN<q_{m+1}`$. We will replace the sequence of fractional parts $`𝒩=\{\{\alpha n^2\}:nN\}`$ by a different sequence depending on the size of $`N`$ relative to $`q_m`$.
### 7.2. Case 1: Assume that $`q_m^{4/3}N<q_{m+1}`$
Recall that $`q_{m+1}q_m^2\mathrm{log}q_m`$ and so this range is nonempty. Replace $`𝒩`$ by the sequence $`𝒩^{}=\{x_n^{}:nN\}`$ where $`x_n^{}=\{p_{m+1}n^2/q_{m+1}\}`$. These two sequences have asymptotically equal correlations since
$$|x_nx_n^{}||\alpha \frac{p_{m+1}}{q_{m+1}}|n^2<\frac{N^2}{q_{m+1}q_{m+2}}\frac{N^2}{r_{m+1}^2q_{m+1}^3}\frac{1}{N(\mathrm{log}N)^2}$$
since $`q_{m+1}>N`$ and $`r_{m+1}=[\mathrm{log}q_{m+1}]\mathrm{log}N`$. Thus by the comparison principle (lemma 5), it suffices to work with the new sequence $`𝒩^{}`$.
By the divergence principle (see (4.1)), since $`q_{m+1}=v_{m+1}^2`$, if the test function $`f`$ is nonvanishing at the origin we find that
$$R^{(5)}(𝒩^{},f)\frac{1}{N}(\frac{N}{v_{m+1}})^5=\frac{N^4}{q_{m+1}^{5/2}}\frac{N^4}{r_m^5q_m^5}.$$
Since $`q_mN^{3/4}`$ and $`r_m\mathrm{log}N`$ we find that
$$R^{(5)}(𝒩^{},f)N^{1/4}(\mathrm{log}N)^5$$
proving (7.1) when $`q_m^{4/3}<N<q_{m+1}`$.
### 7.3. Case 2: $`q_mN<q_m^{4/3}`$
Set
$$M=\frac{q_m^{3/2}}{N^{1/2}}$$
Note that since $`q_mN<q_m^{4/3}`$, $`M`$ lies between $`N^{5/8}`$ and $`N`$. We replace $`𝒩`$ by the sequence $`𝒩^{\prime \prime }=\{x_n^{\prime \prime }:nN\}`$ where
$$x_n^{\prime \prime }=\{\begin{array}{cc}\{\frac{p_m}{q_m}n^2\},\hfill & nM\hfill \\ x_n=\{\alpha n^2\},\hfill & M<nN\hfill \end{array}$$
To check that correlations of $`𝒩`$ and $`𝒩^{\prime \prime }`$ are asymptotically equal, we need to see that $`|x_nx_n^{\prime \prime }|=o(1/N)`$. For $`n>M`$ this certainly holds, while for $`nM`$ we have
$$|x_nx_n^{\prime \prime }||\alpha \frac{p_m}{q_m}|n^2\frac{M^2}{q_mq_{m+1}}$$
Now use $`q_{m+1}r_m^2q_m^2`$ and since $`q_m(M^2N)^{1/3}`$ and $`r_m\mathrm{log}N`$, we have
$$q_mq_{m+1}>r_m^2q_m^3(\mathrm{log}N)^2M^2N$$
which gives $`|x_nx_n^{\prime \prime }|1/N(\mathrm{log}N)^2`$ as required.
Now we study the sequence $`𝒩^{\prime \prime }`$. The number $`0`$ occurs in $`𝒩^{\prime \prime }`$ if $`nM`$ is a multiple of $`v_m`$: $`n=v_mn^{}`$, since then $`x_n^{\prime \prime }=\{\frac{p_m}{q_m}n^2\}=\{p_m(n^{})^2\}=0`$. Thus $`\stackrel{}{0}`$ occurs as a difference of $`5`$-tuples of elements of $`𝒩^{\prime \prime }`$ at least $`[M/v_m]^5`$ times. Thus if the origin lies in the support of the test function $`f`$ then
$$R^{(5)}(𝒩^{\prime \prime },f)\frac{1}{N}(\frac{M}{v_m})^5=\frac{M^5}{Nq_m^{5/2}}$$
Since $`M^2=q_m^3/N`$, and $`N<q_m^{4/3}`$ we get
$$R^{(5)}(𝒩^{\prime \prime },f)\frac{q_m^5}{N^{7/2}}N^{15/47/2}=N^{1/4}$$
## 8. Proof of Theorem 2(b)
Let $`\sigma >23/8`$. We construct $`\alpha =[a_0;a_1,\mathrm{},a_m,\mathrm{}]`$ which will be of type $`\sigma `$ and for which $`lim\; supR^{(5)}(\alpha ,N)=\mathrm{}`$ by an inductive construction of the partial quotients $`a_m`$.
Suppose we have already found $`a_0,\mathrm{},a_{m1}`$, from which we got the partial convergents $`p_j/q_j`$, $`j=0,\mathrm{},m1`$. Now take an integer $`\mathrm{}q_{m1}^{(\sigma 2)/2}`$, which is coprime to $`q_{m1}`$ (this is certainly possible for $`m1`$, say take $`\mathrm{}`$ a prime between $`q_{m1}^{(\sigma 2)/2}`$ and $`2q_{m1}^{(\sigma 2)/2}`$ which does not divide $`q_{m1}`$). Also set $`v_m=\mathrm{}`$.
Because $`\mathrm{}`$ and $`q_{m1}`$ are coprime, there is a unique solution $`t=a_m`$ of the congruence
(8.1)
$$tq_{m1}+q_{m2}=0mod\mathrm{}^2$$
which lies in $`[\mathrm{}^2,2\mathrm{}^2)`$. Then $`a_m\mathrm{}^2q_m^{\sigma 2}`$, and
$$q_m:=a_mq_{m1}+q_{m2}q_{m1}^{\sigma 1}.$$
Thus $`\alpha `$ is of type $`\sigma +ϵ`$ for all $`ϵ>0`$. Moreover $`q_m`$ is divisible by $`v_m^2`$ by (8.1). Thus
$$q_m=u_mv_m^2$$
for some integer $`u_m`$, and
$$v_mq_{m1}^{(\sigma 2)/2}q_m^{(\sigma 2)/(2\sigma 2)}$$
Now take
$$N_m\frac{q_m^{\sigma /3}}{\mathrm{log}q_m}$$
We will see that $`R^{(5)}(\alpha ,N_m)\mathrm{}`$ as $`m\mathrm{}`$.
To see this, note that in the sequence $`\{\alpha n^2:nN_m\}`$ we may replace $`\alpha `$ by the partial convergent $`p_m/q_m`$ without changing the limiting correlations. To see this, note that by Lemma 5 it suffices to check that $`1/q_mq_{m+1}=o(1/N_m^3)`$. Indeed, we have
$$\frac{1}{q_mq_{m+1}}\frac{1}{q_m^\sigma }\frac{1}{N_m^3(\mathrm{log}N_m)^3}$$
as required.
To see that $`R^{(5)}(N_m,\frac{p_m}{q_m},f)`$ diverges for positive test functions $`f`$ with $`f(0)0`$, use the divergence principle (4.1) to find
$$R^{(5)}(N_m,\frac{p_m}{q_m},f)\frac{1}{N_m}(\frac{N_mv_m}{q_m})^5.$$
Now use $`N_mq_m^{\sigma /3}/\mathrm{log}q_m`$ and $`v_mq_m^{(\sigma 2)/(2\sigma 2)}`$ to find
$$R^{(5)}(N_m,\frac{p_m}{q_m},f)\frac{q_m^E}{(\mathrm{log}q_m)^5}$$
where
$$E=\frac{4\sigma }{3}+5\frac{\sigma 2}{2\sigma 2}5=\frac{\sigma (8\sigma 23)}{6(\sigma 1)}.$$
Since $`\sigma >23/8`$, we have $`E>0`$ which gives divergence of $`R^{(5)}`$. ∎
## 9. Proof of Theorem 3
Fix $`m2`$, $`f`$ and $`\delta >0`$. By approximating $`f(0,x_2,\mathrm{},x_m)`$ from above and below with step functions we see that it is enough to prove the statement for a function $`f`$ symmetric, satisfying $`f(x+(t,t,\mathrm{},t))=f(x)`$ for all $`t𝐑`$ and such that $`f(0,x_2,\mathrm{},x_m)`$ is the characteristic function of a nice compact set $`I𝐑^{m1}`$. In other words, given such an $`I`$ and $`m,\delta `$ as above, it is enough to show that as $`q\mathrm{}`$ one has
(9.1)
$$R^{(m)}(N,b/q,I)Vol(I)$$
uniformly for $`(b,q)=1`$ and $`q^{1\frac{1}{2m}+\delta }N\frac{q}{\mathrm{log}q}`$, where $`NR^{(m)}(N,b/q,I)`$ is the number of tuples $`(x_1,\mathrm{},x_m)`$ with distinct components $`x_1,\mathrm{},x_m`$ in $`\{1,\mathrm{},N\}`$ such that
$$N(\{\frac{bx_1^2}{q}\}\{\frac{bx_2^2}{q}\},\mathrm{},\{\frac{bx_{m1}^2}{q}\}\{\frac{bx_m^2}{q}\})I.$$
Given a large prime number $`q`$ and $`b,N`$ as above, we write $`R^{(m)}(N,b/q,I)`$ in the form
$$R^{(m)}(N,b/q,I)=\frac{1}{N}\underset{\stackrel{}{a}sI}{\overset{}{}}\nu (N,\stackrel{}{a})$$
where $`s=\frac{q}{N}`$ is the dilate factor, and
$$\nu (N,\stackrel{}{a})=\mathrm{\#}\{1x_iN:bx_i^2bx_{i+1}^2=a_i(modq),1im1\}.$$
Here $`^{}`$ means the summation is over the vectors $`\stackrel{}{a}`$ for which the partial sums $`A_i=_{ki}a_k,A_m=0`$, are distinct, a condition which comes from the requirement that the $`m`$-tuples $`x=(x_1,\mathrm{},x_m)`$ to be counted in $`R^{(m)}(N,b/q,I)`$ have distinct components. Let
$$h_\stackrel{}{a}(\stackrel{}{x})=\{\begin{array}{cc}1,\hfill & b(x_i^2x_{i+1}^2)=a_j(modq),i=1,\mathrm{},m1\hfill \\ 0\hfill & \text{else.}\hfill \end{array}$$
Thus:
$$\nu (N,\stackrel{}{a})=\underset{1x_1,\mathrm{},x_mN}{}h_\stackrel{}{a}(\stackrel{}{x}).$$
We now use the Fourier expansion:
$$\nu (N,\stackrel{}{a})=\underset{\stackrel{}{r}(modq)}{}\widehat{h}_\stackrel{}{a}(\stackrel{}{r})\underset{i=1}{\overset{m}{}}F_N(r_i)$$
where
$$\widehat{h}_\stackrel{}{a}(\stackrel{}{r})=\frac{1}{q^m}\underset{\stackrel{}{y}(modq)}{}h_\stackrel{}{a}(\stackrel{}{y})e\left(\frac{\stackrel{}{r}\stackrel{}{y}}{q}\right)$$
and:
$$F_N(r_i)=\underset{1x_iN}{}e\left(\frac{r_ix_i}{q}\right).$$
These last sums are geometric series which can be bounded by:
(9.2)
$$F_N(r_i)\mathrm{min}\{N,\frac{q}{|r_i|}\}$$
where the residues $`r_i`$ are assumed to lie in the interval $`[\frac{q}{2},\frac{q}{2}]`$. In
$$R^{(m)}(N,b/q,I)=\frac{1}{N}\underset{\stackrel{}{a}sI}{}\underset{\stackrel{}{r}(modq)}{\overset{}{}}\widehat{h}_\stackrel{}{a}(\stackrel{}{r})\underset{i=1}{\overset{m}{}}F_N(r_i)$$
we isolate the contribution of $`\stackrel{}{r}=0`$ to get the main term :
(9.3)
$$R^{(m)}(N,b/q,I)=+$$
with
(9.4)
$$=N^{m1}\underset{\stackrel{}{a}sI}{\overset{}{}}\widehat{h}_\stackrel{}{a}(0)$$
and
(9.5)
$$=\frac{1}{N}\underset{0\stackrel{}{r}(modq)}{}\underset{i=1}{\overset{m}{}}F_N(r_i)\underset{\stackrel{}{a}sI}{\overset{}{}}\widehat{h}_a(\stackrel{}{r}).$$
We first estimate the main term. For any $`\stackrel{}{a}`$ let $`C(\stackrel{}{a},q)`$ be the curve mod $`q`$ given by the system of congruences:
$$\begin{array}{c}bx_1^2bx_2^2=a_1(modq)\hfill \\ \mathrm{}\hfill \\ bx_{m1}^2bx_m^2=a_{m1}(modq).\hfill \end{array}$$
One has $`\widehat{h}_\stackrel{}{a}(0)=\frac{1}{q^m}\nu (\stackrel{}{a},q)`$, where $`\nu (\stackrel{}{a},q)`$ is the number of points on the curve $`C(\stackrel{}{a},q)`$. Thus
(9.6)
$$=\frac{N^{m1}}{q^m}\underset{\stackrel{}{a}sI}{\overset{}{}}\nu (\stackrel{}{a},q).$$
We want to show that as $`q\mathrm{}`$ one has:
(9.7)
$$=Vol(I)+o(1).$$
For any $`\stackrel{}{a}=(a_1,\mathrm{},a_{m1})`$ denote by $`r_{eff}(\stackrel{}{a},q)`$ the number of distinct $`y_j`$ satisfying the following system:
(9.8)
$$y_iy_{i+1}=a_i(modq),1im1.$$
Since the solutions of the homogeneous system
$$y_iy_{i+1}=0(modq),1im1$$
are spanned by $`(1,\mathrm{},1)`$, $`r_{eff}(\stackrel{}{a},q)`$ is well-defined (independent of the particular solution $`y`$ of (9.8)). Using the Riemann Hypothesis for curves over finite fields (Weil ) one obtains (see also , Proposition 4):
(9.9)
$$\nu (\stackrel{}{a},q)=2^{mr_{eff}(a,q)}(q+B(\stackrel{}{a},q))$$
with
(9.10)
$$|B(\stackrel{}{a},q)|_mq^{\frac{1}{2}}.$$
We define roots $`\sigma _{ij}(\stackrel{}{a})`$,$`1i<jm`$ by
(9.11)
$$\sigma _{ij}(\stackrel{}{a})=\underset{k=i}{\overset{j1}{}}a_k$$
so that $`\sigma _{i,i+1}(\stackrel{}{a})=a_i`$, $`\sigma _{ij}=_{k=i}^{j1}\sigma _{k,k+1}`$. We set $`D(\stackrel{}{a})=_{1ijm}\sigma _{ij}(\stackrel{}{a})`$. The solutions of (9.8) are all distinct (i.e. $`r_{eff}(\stackrel{}{a},q)=m`$) if and only if $`q`$ does not divide $`D(\stackrel{}{a})`$, since $`y_iy_j=_{k=i}^{j1}y_ky_{k+1}=_{k=i}^{j1}a_k=\sigma _{ij}(\stackrel{}{a})`$. Note that $`D(\stackrel{}{a})`$ is a nonzero integer for any $`\stackrel{}{a}`$ which appears in the above summations $`_{asI}^{}`$. In our case $`q`$ does not divide $`D(\stackrel{}{a})`$, since for $`N`$ large enough in terms of $`I`$ each factor $`\sigma _{i,j}(\stackrel{}{a})`$ of $`D(\stackrel{}{a})`$ is in absolute value smaller than $`q`$. Therefore $`r_{eff}(\stackrel{}{a},q)=m`$ and (9.9) and (9.10) give
(9.12)
$$\nu (\stackrel{}{a},q)=q+O_m(q^{\frac{1}{2}})$$
for all $`\stackrel{}{a}`$ which appear in (9.6). Then (9.6) implies that
(9.13)
$$=\frac{N^{m1}}{q^m}(q+O_m(q^{\frac{1}{2}}))\underset{asI}{\overset{}{}}1=\frac{1}{s^{m1}}(1+O_m(\frac{1}{q^{\frac{1}{2}}}))\underset{asI}{\overset{}{}}1.$$
The number of integer points $`\stackrel{}{a}sI`$ which lie in the union of the hyper-planes $`\sigma _{ij}(\stackrel{}{a})=0`$ is $`O_{m,I}(s^{m2})`$, while by the Lipschitz principle (see Davenport ) it follows that:
$$\mathrm{\#}(sI𝐙^{m1})=s^{m1}Vol(I)+O_{m,I}(s^{m2}).$$
Therefore:
(9.14)
$$\underset{asI}{\overset{}{}}1=\mathrm{\#}(sI𝐙^{m1})\mathrm{\#}\{\stackrel{}{a}sI:D(\stackrel{}{a})=0\}$$
$$=s^{m1}Vol(I)+O_{m,I}(s^{m2})$$
and from (9.13) we get
$$=(1+O_m(\frac{1}{\sqrt{q}}))(1+O_{m,I}(\frac{1}{s}))$$
which proves (9.7).
We now proceed to estimate the remainder $``$. For any $`\stackrel{}{a}`$ and $`\stackrel{}{r}`$ we have:
$$\widehat{h}_\stackrel{}{a}(\stackrel{}{r})=\frac{1}{q^m}\underset{\stackrel{}{y}C(\stackrel{}{a},q)}{}e\left(\frac{\stackrel{}{r}\stackrel{}{y}}{q}\right).$$
Applying Weil’s Riemann Hypothesis for curves over finite fields one has (see , Theorem 6)
(9.15)
$$\left|\underset{yC(a,q)}{}e\left(\frac{\stackrel{}{r}\stackrel{}{y}}{q}\right)\right|_m\sqrt{q}$$
unless the linear form $`\stackrel{}{r}\stackrel{}{y}`$ is constant along the curve. For $`\stackrel{}{a}`$ as in (9.5) this only happens if $`\stackrel{}{r}=0`$. For, let $`\stackrel{}{r}0`$ be such that $`\stackrel{}{r}\stackrel{}{y}`$ is constant along the curve. Then, in the function field $`\overline{𝐅}_q(Y_1,\mathrm{},Y_m)`$ of the curve, where $`\overline{𝐅}_q`$ denotes the algebraic closure of $`𝐅_q=𝐙/q𝐙`$, $`Y_1`$ is a variable and $`Y_2,\mathrm{},Y_m`$ are algebraic functions such that
$$Y_i^2=Y_1^2\frac{a_1+\mathrm{}+a_{i1}}{b}$$
for $`2im`$, we will have an equality $`\stackrel{}{r}\stackrel{}{Y}=c`$, with $`c\overline{𝐅}_q`$. If we choose $`j_0\{1,\mathrm{},m\}`$ such that $`r_{j_0}0`$ then $`Y_{j_0}`$ will lie in $`\overline{𝐅}_q(Y_1,\mathrm{},Y_{j_01},Y_{j_0+1},\mathrm{},Y_m)`$ and hence
(9.16)
$$\overline{𝐅}_q(Y_1,\mathrm{},Y_m)=\overline{𝐅}_q(Y_1,\mathrm{},Y_{j_01},Y_{j_0+1},\mathrm{},Y_m).$$
Now for any unique factorization domain $`D`$ of characteristic $`2`$ and any distinct primes $`p_1,\mathrm{},p_t`$ in $`D`$ one has
$$[K(\sqrt{p_1},\mathrm{},\sqrt{p_t}):K]=2^t$$
where $`K`$ denotes the quotient field of $`D`$ (see Besicovitch )). Applying this with $`D=\overline{𝐅}_q[Y_1]`$ and $`p_i=Y_1^2\frac{a_1+\mathrm{}+a_i}{b}`$ for $`1im1`$ we get:
$$[\overline{𝐅}_q(Y_1,\mathrm{},Y_m):\overline{𝐅}_q(Y_1)]=2^{m1}.$$
By the same argument we see that
$$[\overline{𝐅}_q(Y_1,\mathrm{},Y_{j_01},Y_{j_0+1},\mathrm{},Y_m):\overline{𝐅}_q(Y_1)]=2^{m2}$$
which contradicts (9.16). It follows that for all $`\stackrel{}{r}`$ and $`\stackrel{}{a}`$ which appear in (9.5), the inequality (9.15) holds true and one has:
$$|\widehat{h}_\stackrel{}{a}(\stackrel{}{r})|_m\frac{1}{q^{m\frac{1}{2}}}.$$
This implies that
(9.17)
$$||_m\frac{1}{Nq^{m\frac{1}{2}}}\underset{0\stackrel{}{r}(modq)}{}\left(\underset{i=1}{\overset{m}{}}|F_N(r_i)|\right)\underset{\stackrel{}{a}sI}{\overset{}{}}1.$$
We use (9.2) and (9.14) in (9.17) to conclude that
(9.18)
$$||_{m,I}\frac{s^{m1}}{Nq^{m\frac{1}{2}}}\underset{\stackrel{}{r}(modq)}{}\underset{i=1}{\overset{m}{}}\mathrm{min}\{N,\frac{q}{|r_i|}\}$$
$$_m\frac{q^{m\frac{1}{2}}\mathrm{log}^mq}{N^m}\left(\frac{\mathrm{log}q}{q^\delta }\right)^m.$$
The theorem now follows from (9.3), (9.7) and (9.18).
## Appendix A Square factors of rational approximants
Let $`\alpha `$ be a real number and $`a_n/q_n`$ a sequence of rational approximants of $`\alpha `$: $`|\alpha a_n/q_n|<1/q_n^2`$, and $`q_n\mathrm{}`$. In view of Theorem 1’, we want to investigate the square parts of the denominators $`q_n`$, keeping in mind that large square parts rule out Poisson statistics for the correlation functions.
###### Definition A.1.
A sequence $`\{q_n\}`$ is almost square-free if $`ϵ>0`$, all square divisors $`s_n^2`$ of $`q_n`$ satisfy $`s_n_ϵq_n^ϵ`$.
### A.1. A metric result
We will show that for almost all $`\alpha `$, we have: If $`a_n/q_n`$ is a sequence of rational approximants of $`\alpha `$ (that is $`|\alpha a_n/q_n|<1/q_n^2`$, and $`q_n\mathrm{}`$), then $`\{q_n\}`$ is almost square-free.
In fact, we show more: For an integer $`q1`$, we write $`q=\stackrel{~}{q}s^2`$ with $`\stackrel{~}{q}`$ square-free. Let $``$ be the set of integers $`q`$ whose largest square factor $`s^2`$ satisfies $`s\mathrm{log}^2\stackrel{~}{q}`$. We will show that almost all reals $`\alpha `$ have rational approximants whose denominators are in $``$ except for finitely many exceptions.
###### Proposition 7.
For all reals $`\alpha `$ outside a set of measure zero, there is a $`Q=Q(\alpha )>1`$ so that if $`|\alpha a/q|<1/q^2`$ and $`qQ`$ then $`q`$.
The proof of this follows from a well-known general principle: Given a sequence of integers $`𝒩`$, we say that a real number $`\alpha `$ is $`𝒩`$-approximable if there are infinitely many rationals $`a/q\alpha `$ with denominator $`q𝒩`$ and $`|\alpha a/q|<1/q^2`$. For instance, we may take as $`𝒩`$ the complement of $``$. To prove Proposition 7, we will use
###### Lemma 8.
Suppose that $`𝒩`$ is a sequence such that
$$\underset{q𝒩}{}\frac{1}{q}<\mathrm{}.$$
Then the set of $`𝒩`$-approximable reals has measure zero.
###### Proof.
Without loss of generality we will assume that $`0<\alpha <1`$. For each pair of coprime integers $`(a,q)`$ with $`1a<q`$, denote by $`I_{a,q}`$ the interval
$$I_{a,q}=(\frac{a}{q}\frac{1}{q^2},\frac{a}{q}+\frac{1}{q^2})$$
Then $`\alpha `$ is $`𝒩`$-approximable if and only if it lies in infinitely many of the intervals $`I_{a,q}`$ with $`q𝒩`$. That is for all $`N1`$, $`\alpha `$ lies in
$$M_N:=_{Nq𝒩}_{1a<q}I_{a,q}.$$
Thus we need to compute the measure of $`M:=_{N1}M_N`$. Since $`M_NM_{N+1}\mathrm{}`$, we have
$$\text{meas}(M)=\underset{N}{lim}\text{meas}(M_N)\underset{N}{lim}\underset{Nq𝒩}{}\underset{a=1}{\overset{q}{}}\text{meas}(I_{a,q})\underset{N}{lim}\underset{Nq𝒩}{}\frac{1}{q}$$
(allowing overlap of the intervals). Since $`_{q𝒩}1/q<\mathrm{}`$, the above limit is zero. ∎
Thus to prove Proposition 7, it suffices to show
$$\underset{q}{}\frac{1}{q}<\mathrm{}.$$
We rewrite this sum by grouping together those $`q`$ with the same square-free kernel $`\stackrel{~}{q}`$: Writing $`q=fm^2`$, $`\stackrel{~}{q}=f`$, then
$$\underset{q}{}\frac{1}{q}=\underset{f\text{ square-free}}{}\underset{\begin{array}{c}\stackrel{~}{q}=f\\ q\end{array}}{}\frac{1}{q}=\underset{f\text{ square-free}}{}\frac{1}{f}\underset{q=fm^2}{}\frac{1}{m^2}$$
Now if $`q`$, $`\stackrel{~}{q}=f`$ then $`m>\mathrm{log}^2f`$. Thus for each $`f`$,
$$\underset{q=fm^2}{}\frac{1}{m^2}=\underset{m>\mathrm{log}^2f}{}\frac{1}{m^2}\frac{1}{\mathrm{log}^2f}$$
and so
$$\underset{q}{}\frac{1}{q}\underset{f\text{ square-free}}{}\frac{1}{f}\frac{1}{\mathrm{log}^2f}<\mathrm{}$$
as required. ∎
### A.2. Algebraic $`\alpha `$
For real algebraic $`\alpha `$, the analogue of Proposition 7 follows from a standard belief in diophantine analysis, namely the “ABC Conjecture” of Masser and Oesterle: Define the radical of an integer $`N`$ as the product of all primes dividing it: $`\mathrm{rad}(N):=_{pN}p`$. The ABC conjecture is the assertion that whenever we have an equation in coprime integers $`A+B+C=0`$, then
(A.1)
$$|A|_ϵ\mathrm{rad}(ABC)^{1+ϵ}$$
for all $`ϵ>0`$. This implies a seemingly stronger statement: Suppose that $`G(x,y)𝐙[x,y]`$ is a homogeneous form with integer coefficients and no repeated factors, and $`m,n`$ coprime integers. Then for all $`ϵ>0`$
(A.2)
$$\mathrm{max}(|m|,|n|)^{\mathrm{deg}(G)2\epsilon }_ϵ\mathrm{rad}(G(m,n)),$$
where $`\mathrm{deg}(G)`$ is the degree of $`G`$. The deduction of (A.2) from (A.1) and a theorem of Belyi was noted by Elkies and by Langevin . The ABC-conjecture (A.1) is the special case of the ternary form $`G(x,y)=xy(x+y)`$.
The corollary (A.2) of the ABC conjecture implies the analogue of Proposition 7 for irrational algebraic $`\alpha `$. Indeed, let $`f(x)`$ be the minimal polynomial of $`\alpha `$, of degree $`d>1`$, and write $`f(x/y)=F(x,y)/y^d`$ with $`F(x,y)𝐙[x,y]`$. Suppose that $`p/q`$ is an approximant of $`\alpha `$: $`|\alpha p/q|<1/q^2`$, with $`p`$, $`q`$ coprime. Since $`f(x)`$ is irreducible, $`f^{}(\alpha )0`$ and thus by the mean value theorem, for some $`\xi `$ between $`\alpha `$ and $`p/q`$,
$$|f(\frac{p}{q})|=|f(\frac{p}{q})f(\alpha )|=|\alpha \frac{p}{q}||f^{}(\xi )|\frac{1}{q^2}$$
On the other hand,
$$f(\frac{p}{q})=\frac{F(p,q)}{q^d}$$
and so we find
$$|F(p,q)|q^{d2}.$$
By (A.2), taking $`G(x,y)=xyF(x,y)`$ and noting that $`|p|q`$, we get for all $`ϵ>0`$
$$q^{dϵ}_ϵ\mathrm{rad}(pqF(p,q))|pF(p,q)|\mathrm{rad}(q)q^{d1}\mathrm{rad}(q).$$
Thus if $`q=\stackrel{~}{q}s^2`$ then
$$(\stackrel{~}{q}s^2)^{dϵ}_ϵ\mathrm{rad}(\stackrel{~}{q}s)\stackrel{~}{q}s$$
and so $`s_ϵq^ϵ`$. |
warning/0002/hep-th0002172.html | ar5iv | text | # The Holographic RG flow to conformal and non-conformal theory11footnote 1Talks presented at the TMR conference in Paris, September 99.
## 1 Introduction
The AdS/CFT correspondence has deserved some surprises when extended outside the realm of stric- tly conformally invariant theories. The study of the supergravity dual of RG flows has flourished, both in the concrete application to SYM theories and in a general setting -. Asymptotically AdS<sub>d+1</sub> backgrounds, breaking the full $`O(d,2)`$ invariance but preserving at least $`d`$-dimensional Poincaré invariance, describe RG flows for a $`d`$-dimensional CFT. These supergravity solutions with an asymptotic AdS region have a double QFT interpretation: deformations of an UV fixed point versus the same theory in a different vacuum . Both cases have been extensively studied. Many results have been obtained upon reduction to a $`d+1`$-dimensional effective theory, where the RG flow can be studied in terms of a theory of scalar fields coupled to gravity. In this simple set-up, the RG flows are identified as domain-walls interpolating between AdS<sub>d+1</sub> vacua (or approaching infinity on one side), and general results are very easy to obtain. The correspondence defines a holographic scheme, where beta and $`c`$-functions have a natural definition. A $`c`$-theorem, for example, can be easily proven . Moreover, it is possible to obtain the quantum field theory RG equations from supergravity <sup>2</sup><sup>2</sup>2Notice that the holographic beta and $`c`$-functions do not need to coincide with analogous functions defined in schemes that are more natural from the QFT point of view ..
The study of RG flows between CFTs (at large $`N`$ and strong coupling) can be rigorously performed using supergravity. The phase space of massive deformations of the N=4 SYM theory has been throughly investigated and several IR fixed points have been found . The results are on solid grounds because supergravity is valid all along the RG flow. Still problematic is the precise mapping of some QFT couplings to supergravity quantities. For example, it is still unclear what in supergravity corresponds to the running of the gauge coupling.
Most of the unsolved problems concern the flows to non-conformal theories, where supergravity is invalidated by a (typically naked) singularity in the IR region of the flow. Solutions flowing to infinity for a generic 5$`d`$-Lagrangian are certainly a dense set in the space of solutions. The full recipe for selecting the physical ones is still unclear<sup>3</sup><sup>3</sup>3A criterion for selecting physical solutions has been recently proposed in .. The distinction between deformations and vacua of an UV fixed point helps but does not solve the problem. Supersymmetric and supersymmetric-inspired solutions however are uniquely selected because the equations of motion can be reduced to first order ones . N=4 Coulomb branch solutions have been studied in . Here we focus on the flow to N=1 SYM. Despite the singularity, we obtain a good qualitative agreement with quantum field theory expectations already at the level of supergravity.
Since singularities are apparently unavoidable in interesting supergravity solutions, it is mandatory to understand their fate in the full string theory, where they must be resolved. Available options are the chance that the singularity is an artifact of the dimensional reduction to 5 dimensions, mechanisms such that proposed in and, more generally, some help from string corrections.
The supergravity solutions with an asymptotic AdS region certainly have many other applications. Relaxing the $`d`$-Poincaré invariance, we have examples of RG flow due to finite temperature. This is indeed the firstly proposed me- thod for discussing non-conformal theories from AdS and the one not suffering from unpleasant singularities. Cutting the AdS-boundary, we can describe CFTs coupled to gravity and make contact with the large extra-dimension scenario . We will not discuss this issue here, but we simply notice that singular solutions have been recently considered in this context.
## 2 RG Flow from 5$`d`$ Supergravity
In general, we interpret the $`(d+1)`$-th coordinate $`y`$ of AdS<sub>d+1</sub> as an energy scale . RG flows between CFTs then correspond to type II or M-theory supergravity solutions interpolating (along $`y`$) between AdS$`{}_{d+1}{}^{}\times _{W}^{}H`$ vacua.
The very first example of RG flow in the AdS/CFT correspondence is manifest in the multi-centre supergravity solution for D3-branes . This represents the Coulomb branch of N=4 SYM. Given two sets of $`N`$ and $`M`$ branes at different points, the near-horizon geometry is AdS<sub>5</sub> with radius $`\sqrt{N+M}`$ far from both sets of branes, and AdS<sub>5</sub> with radius $`\sqrt{N}`$ near one set. In QFT this is the RG flow between the $`U(N+M)`$ N=4 CFT in the UV, where the Higgs VEVs can be neglected, and the $`U(N)`$ N=4 CFT in the IR. A more sophisticated example was found in . A supergravity solution interpolating between AdS$`{}_{5}{}^{}\times S^5/Z_2`$ and AdS$`{}_{5}{}^{}\times T^{1,1}`$ was also interpreted on the QFT side as a RG flow between CFTs. It is a supersymmetric massive deformation of the N=2 $`SU(N)\times SU(N)`$ theory corresponding to a $`Z_2`$ orbifold of N=4 SYM which flows to an N=1 IR fixed point. Many successful checks of this interpretation have been performed .
However, interpolating 10$`d`$ backgrounds are difficult to find. Sometimes dimensional reduction to 5 dimensions helps.
The RG flow has a natural description in 5$`d`$.
Consider a certain UV CFT and suppose we have the corresponding 5$`d`$ Lagrangian and that it contains all the fields/modes we are interested in. The effective 5$`d`$ Lagrangian we need is just the most general Lagrangian for scalars coupled to gravity
$$L=\sqrt{g}\left[\frac{R}{4}+\frac{1}{2}g^{IJ}_I\lambda _a_J\lambda _bG^{ab}+V(\lambda )\right].$$
(1)
The scalars $`\lambda _a`$ can either be the massless modes or Kaluza-Klein modes of the compactification to 5 dimensions. The form of the potential depends on the particular case we are considering. We may have, for example, N=8 gauged supergravity, which describes N=4 SYM and most of its bilinear relevant operators (almost all of the masses for scalars and fermions). Or we may have an N=4 theory describing the orbifold $`R^4/Z_2`$ and the supersymmetric mass term that drives the theory to an N=1 IR fixed point. Or else we may have the Lagrangian for some of the KK modes. The interactions among the modes in the graviton multiplet in 5$`d`$ can be found using supersymmetry. In particular, for the N=4 SYM case, the 5$`d`$ Lagrangian for the massless modes is uniquely fixed by supersymmetry in the form of the N=8 gauged supergravity . All mass terms for the scalars and the fermions contained in the KK spectrum are associated to modes in the gauged supergravity. 5-dimensional supersymmetric Lagrangians have been discussed also for less supersymmetry, but the uniqueness of N=8 supergravity is lost and interesting modes are split into various vector, tensor and hyper-multiplets. One needs some help from QFT intuition in identifying the right potential. In principle, $`V(\lambda )`$ can be obtained for all modes (often with non-trivial effort) by dimensional reduction from 10 dimensions.
If the UV CFT perturbed by a particular operator $`O_\lambda `$ flows in the IR to another CFT, the potential $`V`$ must have a critical point for non-zero value of the scalar field $`\lambda `$. Analogously, the dual of the flow to a non-conformal field theory is given by the flow from one minimum of the potential to infinity.
The 5$`d`$ description of the RG flow between conformal theories is a kink solution, which interpolates between the two critical points. A 4$`d`$ Poincaré invariant metric is
$$ds^2=dy^2+e^{2\varphi (y)}dx^\mu dx_\mu ,\mu =0,1,2,3.$$
(2)
AdS corresponds to $`\varphi =y/R`$. We then look for solutions with asymptotics: $`\varphi (y)y/R_{UV,IR}`$ for $`y\pm \mathrm{}`$; $`\lambda (y)0`$ for $`y\mathrm{}`$, while $`\lambda (y)\lambda _{IR}`$ for $`y\mathrm{}`$. We associate larger energies with increasing $`y`$.
The equations of motion for the scalars and the metric read
$`\ddot{\lambda }_a+4\dot{\varphi }\dot{\lambda }={\displaystyle \frac{V}{\lambda _a}},`$
$`6(\dot{\varphi })^2={\displaystyle \underset{a}{}}(\dot{\lambda }_a)^22V.`$ (3)
With the above boundary conditions and a reasonable shape of the potential, a kink interpolating between critical points always exists .
As an example of flows between conformal field theories, we can discuss the mass deformations of N=4 SYM. These can be studied in the context of N=8 gauged supergravity, where the form of the potential $`V`$ is known. N=8 gauged supergravity is the low energy effective action for the “massless” modes of the compactification of type IIB on $`AdS_5\times S_5`$. It is believed to be a consistent truncation of type IIB on $`S^5`$ in the sense that every solution of the 5$`d`$ theory can be lifted to a consistent 10$`d`$ type IIB solution. Five-dimensional gauged supergravity has 42 scalars, which transform under the N=4 YM R-symmetry $`SU(4)`$ as $`\underset{¯}{1},\underset{¯}{20},\underset{¯}{10}`$. The singlet is associated with the marginal deformation corresponding to a shift in the coupling constant of the N=4 theory. The mode in the $`\underset{¯}{20}`$ has mass square $`M^2=4`$ and is associated with a symmetric traceless mass term for the scalars $`\text{Tr }\varphi _i\varphi _j`$, ($`i,j=1,\mathrm{},6`$) with $`\mathrm{\Delta }=2`$. The $`\underset{¯}{10}`$ has mass square $`M^2=3`$ and corresponds to the fermion mass term $`\text{Tr }\lambda _A\lambda _B`$, ($`A,B=1,\mathrm{},4`$) of dimension 3. Thus the scalar sector of N=8 gauged supergravity is enough to discuss at least all mass deformations that have a supergravity description<sup>4</sup><sup>4</sup>4 The only missing state is $`\text{Tr }_i^6\varphi _i^2`$, the prototype of a stringy states in the correspondence. Even without this state, we can study almost all massive deformations of the N=4 theory and all these deformations can be described by just the Lagrangian for the massless multiplet..
The scalar potential $`V`$ in eq.(1) is known and it turns out to have only isolated minima (apart from one flat direction, corresponding to the dilaton). Up to now, all critical points with at least $`SU(2)`$ symmetry have been classified . There is a central critical point with $`SO(6)`$ symmetry and with all the scalars $`\lambda _a`$ vanishing: it corresponds to the unperturbed N=4 YM theory. There are three N=0 theories with residual symmetry $`SU(3)\times U(1)`$, $`SO(5)`$ and $`SU(2)\times U(1)^2`$. They correspond to non-zero VEV for some of the scalars in the $`\underset{¯}{10}`$, $`\underset{¯}{20}`$, and $`\underset{¯}{10}+\underset{¯}{20}`$, respectively. Then there is an N=2 point with symmetry $`SU(2)\times U(1)`$, obtained giving VEV to scalars in the $`\underset{¯}{10}+\underset{¯}{20}`$ . According to the AdS/CFT correspondence, these other minima should correspond to IR conformal field theories<sup>5</sup><sup>5</sup>5The symmetries of field theories can be read from those of the supergravity minima according to the correspondence : gauge symmetry in supergravity $``$ global symmetry in field theory, supersymmetry in supergravity $``$ superconformal symmetry in field theory.. The following IR CFT theories can be obtained as mass deformations of N=4 SYM:
* Three N=0 theories with symmetry $`SU(3)\times U(1)`$, $`SO(5)`$ and $`SU(2)\times U(1)^2`$. All these theories are unstable and correspond to non-unitary CFTs. A natural question arises: are all the N=0 critical points unstable?
* A stable N=1 theory with symmetry $`SU(2)`$ $`\times U(1)`$. It corresponds to the N=4 theory deformed with a mass for one of the three N=1 chiral superfields. Results and supergravity description are almost identi- cal to the $`T^{1,1}`$ case, which is just a $`Z_2`$ projection of this example.
### 2.1 Central charges
In a supersymmetric gauge field theory in 4$`d`$, the trace and R-symmetry anomaly are given by
$`T_\mu ^\mu `$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\beta }}{2g^2}}F_{\mu \nu }^2+{\displaystyle \frac{c}{16\pi ^2}}W_{\mu \nu \rho \sigma }^2{\displaystyle \frac{a}{16\pi ^2}}\stackrel{~}{R}_{\mu \nu \rho \sigma }^2`$ (4)
$`+`$ $`{\displaystyle \frac{c}{6\pi ^2}}V_{\mu \nu }^2+{\displaystyle \frac{b}{32\pi ^2}}B_{\mu \nu }^2,`$
$`_\mu \sqrt{g}R^\mu `$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{\beta }}{3g^2}}F_{\mu \nu }\stackrel{~}{F^{\mu \nu }}{\displaystyle \frac{ac}{24\pi ^2}}R_{\mu \nu \rho \sigma }\stackrel{~}{R}^{\mu \nu \rho \sigma }`$ (5)
$`+`$ $`{\displaystyle \frac{5a3c}{9\pi ^2}}V_{\mu \nu }\stackrel{~}{V}^{\mu \nu }{\displaystyle \frac{b}{48\pi ^2}}B_{\mu \nu }\stackrel{~}{B^{\mu \nu }}.`$
Here $`W_{\mu \nu \rho \sigma }`$ and $`R_{\mu \nu \rho \sigma }`$ are the Weyl and curvature tensors for an external metric $`g_{\mu \nu }`$ that couples to the energy-momentum tensor $`T_{\mu \nu }`$. Similarly $`V_{\mu \nu }`$ and $`B^{\mu \nu }`$ are the field strengths of the external sources $`V_\mu ,B_\mu `$ that couple to the R-symmetry and flavour currents, respectively. $`F_{\mu \nu }`$ is the gauge field strength and $`\stackrel{~}{\beta }`$ is the numerator of the exact beta-function .
The external anomaly coefficients $`a`$ and $`c`$ have a straightforward interpretation in the dual supergravity theory.
$`c`$ is the central charge of the CFT, and it is associated with the cosmological constant at the critical points. From eq. (1), we can see by a simple scaling that, at least at the fixed points, where $`ds^2=R^2[dy^2+\mathrm{exp}(2y)_idx_i^2]`$,
$$T(x)T(0)=\frac{c}{|x|^8}cR^3(\mathrm{\Lambda })^{3/2}.$$
(6)
This scaling reproduces the known results for $`c`$ . More interestingly, one can prove that for the class of field theories that have a supergravity dual a $`c`$-theorem exists. Indeed we can exhibit a $`c`$-function that is monotonically decreasing along the flow . The $`c`$-function
$$c(y)(T_{yy})^{3/2},$$
(7)
is constructed with the $`y`$ component of the stress-energy tensor
$$T_{yy}=6(\dot{\varphi })^2=\underset{a}{}(\dot{\lambda }_a)^22V.$$
(8)
At the critical points, where $`\dot{\lambda }_a=0`$,
$$c(y)=c_{UV,IR}(V)_{UV,IR}^{3/2}\mathrm{\Lambda }_{UV,IR}^{3/2},$$
(9)
and using the equations of motion ($`\ddot{\varphi }<0`$) and the boundary conditions one can easily check that $`c(y)`$ is monotonic .
Let us consider $`a`$. AdS computations showed that $`a=c`$ for all CFTs that have an AdS dual.
It is then natural to ask what can AdS/CFT correspondence say about the coefficient $`b`$ <sup>6</sup><sup>6</sup>6These results have been obtained in collaboration with D. Anselmi and L. Girardello.. The coefficient $`b`$ is related to the two-point function of the flavour (global) symmetry currents . According to AdS/CFT correspondence the R-symmetry and flavour currents are associated to the gauge fields of the SUGRA Lagrangian
$$J_\mu ,R_\mu A_\mu .$$
(10)
One should then be able to read the $`b`$ (and $`a`$) coefficient from the kinetic terms of the corresponding SUGRA modes. The generic 5$`d`$-Lagran- gian we are interested in has the following structure
$$L=\sqrt{g}\left[\frac{R}{4}+\mathrm{\Lambda }+fF_{\mu \nu }^2+f_RF_{\mu \nu R}^2\right].$$
(11)
Here $`F_{\mu \nu R}`$ and $`F_{\mu \nu }`$ represent the kinetic terms for the fields corresponding to the R-symmetry and flavour symmetry currents, respectively. At the critical points (or generically for a metric of the form (2)), one obtains by scaling
$$J(x)J(0)=\frac{b}{|x|^6}bfRfc^{1/3}.$$
(12)
A similar behaviour is obtained for the R-symme- try currents. In this case, supersymmetry <sup>7</sup><sup>7</sup>7The R-symmetry currents are in the same multiplet as the energy-momentum tensor. implies $`b=c`$, and the previous equation can be used as a check of the consistency of the procedure.
The values of the coefficients $`f`$ and $`f_R`$ depend on the particular model under consideration. Consider for example the massive deformations of N=4 SYM, for which we have the dual supergravity Lagrangian: that of N=8 gauged supergravity. In this case, the kinetic term for the gauge fields is expressed in terms of the vielbein parametrising the scalar manifold . To determine $`f`$ and $`f_R`$ we have then to evaluate the contractions of the vielbein and therefore these coefficients depend on the critical point and on the way the UV $`SU(4)`$ group is broken (for instance, $`SU(4)`$ $``$ $`SU(3)`$ $`\times `$ $`U(1)_R`$, or $`SU(4)`$ $``$ $`SU(2)`$ $`\times `$ $`U(1)_R`$, …). We now want to compute the charge $`b`$ for the global non-abelian symmetry group preserved a- long the flow (e.g. $`SU(3)`$, $`SU(2)`$, …). The computation of the coefficients $`f`$ can be performed using the results of for most of the critical points. Alternatively, using the parametrisation in appendix A of , it is easy to convince themselves that
$$f=e^{4\alpha }.$$
(13)
Here $`\alpha `$ is the scalar in the $`\underset{¯}{20}`$ of $`SU(4)`$ corresponding to a mass term for the scalars in N=4 SYM . The value of the scalar $`\alpha `$ and $`c`$ for the various fixed points can be found in . One then gets the following results for the coefficient $`b`$ :
* N=1 point with symmetry $`SU(2)\times U(1)`$. $`\frac{b_{IR}}{b_{UV}}=\frac{3}{2}`$. This is the only case where comparison with field theory is possible. Consider a set of N=1 chiral superfields $`X_i`$ in the representation $`R_i`$ of the gauge group and in the representation $`T_i`$ of the flavour symmetry group. Then, because of supersymmetry, the following formula holds
$$b_{UV}b_{IR}=3\underset{ij}{}\left(dimR_i\right)\left[\left(r_i\frac{2}{3}\right)T_i^jT_j^i\right],$$
(14)
where $`r_i`$ is IR R-symmetry charge of the field $`X_i`$ and $`T_i^j`$ are the generators of the flavour group in the representation $`T_i`$. It is straightforward to check that the supergravity and the field theory computations agree.
* N=0 theories. For the $`SU(3)\times U(1)`$, $`SO(5)`$ and $`SU(2)\times U(1)^2`$ symmetric points, we have $`\frac{b_{IR}}{b_{UV}}=\frac{2\sqrt{2}}{3}`$, $`\frac{b_{IR}}{b_{UV}}=\sqrt{2}`$ and $`\frac{b_{IR}}{b_{UV}}=2`$, respectively.
In it was observed that for several examples of supersymmetric gauge theory $`b`$ increa- ses going from the UV to the IR. This was suggestive of possible anti-$`b`$-theorem. The same authors however pointed out that for non-supersym- metric gauge theories $`b`$ has no universal behaviour, and that also a large class of supersymmetric theories violates the relation $`b_{IR}/b_{UV}>1`$. Then it is not possible to state any anti-$`b`$-theorem in field theory. It is interesting to see what are the supergravity results. Consider first the non-supersymmetric cases. For the point $`SU(3)\times U(1)`$ we have $`b_{IR}/b_{UV}<1`$, which violates the anti-$`b`$-theorem. The situation is different for the supersymmetric point $`SU(2)\times U(1)`$. In this case the coefficient $`b`$ increases along the flow. The same analysis carried on for the massive flow to N=1 super Yang-Mills (see section 4) or for the Coulomb branch of N=4 SYM seems to indicate a similar behaviour.
Notice that the theories that have a supergravity dual represent a very restricted class of gauge theories. First of all these theories always have $`a=c`$, which is in general not the case in field theory. It has been argued that the requirement $`a=c`$ simplifies the structure and OPEs of a CFT, making it most similar to a two dimensional conformal field theory . Secondly it has been suggested (see and next section) that all these theories could be characterised by having a pre-potential. It could then be possible, and interesting to check, whether an anti-$`b`$-theorem could hold for this particular class of gauge theories.
The previous results on $`b`$ could have been obtained from the analysis of the Chern-Simons terms of the N=8 Lagrangian, which contain all information about global anomalies . In particular, $`b`$ can be read from the $`SU(2)^2\times U(1)_R`$ anomaly coefficient, which can be extracted from the Chern-Simon terms. It is easy to check, using the results in , that the result for $`b`$ coincides with the previously obtained one<sup>8</sup><sup>8</sup>8It is crucial to pay attention to normalisations and the definition of $`U(1)_R`$, which varies from UV to IR.. Notice that the Chern-Simon terms uniquely determine the form of a supersymmetric gauge supergravity. From the knowledge of the global anomaly, we should be able to reconstruct the entire AdS Lagrangian for massless modes for a given supersymmetric CFT fixed point .
### 2.2 Vacua and deformations
We end this section with a brief discussion of a point that will play an important role in our analysis, namely the fact that supergravity solutions can represent both deformations of a CFT and different vacua of the same theory . The running of coupling constants and parameters along the RG flow can be induced in the UV theory in two different ways: by deforming the CFT with a relevant operator, or by giving a nonzero VEV to some operators. The asymptotic UV behaviour discriminates between the two options. In the asymptotic AdS-region, we just need a linearised analysis. A scalar fluctuation $`\lambda (y)`$ in the asymptotically AdS background must satisfy
$$\ddot{\lambda }+4\dot{\lambda }=M^2\lambda ,$$
(15)
where the dot means the derivative with respect to $`y`$. The previous equation has a solution depending on two arbitrary parameters
$$\lambda (y)=Ae^{(4\mathrm{\Delta })y}+Be^{\mathrm{\Delta }y},$$
(16)
where $`\mathrm{\Delta }`$ is the dimension of the operator, $`M^2=\mathrm{\Delta }(\mathrm{\Delta }4)`$ . We are interested in the case of relevant operators, where $`\mathrm{\Delta }4`$. From the basic prescription of the AdS/CFT, we associate solutions behaving as $`e^{(4\mathrm{\Delta })y}`$ with deformations of the N=4 theory with the operator $`O_\lambda `$. On the other hand, solutions asymptotic to $`e^{\mathrm{\Delta }y}`$ (the subset with $`A=0`$) are associated with a different vacuum of the UV theory, where the operator $`O_\lambda `$ has a non-zero VEV <sup>9</sup><sup>9</sup>9We are not careful about subtleties for particular values of $`\mathrm{\Delta }`$ ..
Since in general the UV-IR interpolating solution is not known, it is not even obvious whether a particular solution corresponds to a deformation or to a different vacuum. For many problems, we may invoke supersymmetry. It helps in finding the solution all along the flow and in unambiguously identifying the UV behaviour. In ref. the conditions for a supersymmetric flow were found. As usual, a solution for which the fermionic shifts vanish, automatically satisfies the equations of motion. Moreover, this shortcut reduces the second order equations to first order ones. For a supersymmetric solution, the potential $`V`$ can be written in terms of a superpotential $`W`$ as
$$V=\frac{1}{8}\underset{a=1}{\overset{n}{}}\left|\frac{W}{\lambda _a}\right|^2\frac{1}{3}\left|W\right|^2,$$
(17)
where $`W`$ is one of the eigenvalues of the tensor $`W_{ab}`$ defined in . The equations of motion reduce to
$`\dot{\lambda }_a`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{W}{\lambda _a}},`$ (18)
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{1}{3}}W.`$ (19)
It is easy to check that a solution of eq.(19) satisfies also the second order equations (3).
It is quite plausible and generally assumed that all the supergravity flows connecting fixed points correspond to deformations of the UV fixed point.
## 3 Confining Solutions
Solutions flowing to infinity represent RG flows to non-conformal theories, which may exist in various phases in the IR. These kinds of solution are difficult to classify. In many cases the asymptotic IR behaviour is known, but the entire solution along the flow can not be found. Typically, we encounter a singularity somewhere along the flow. Many solutions exhibit a logarithmic divergence at finite $`y_0`$ for the scalar fields, $`\lambda _aB_a\mathrm{log}|yy_0|`$, and the metric, $`\varphi A\mathrm{log}|yy_0|`$. There are many criteria for studying the IR properties and the phase of these solutions. One of them, the Wilson loop, will be discussed later. The spectrum can be determined also from two-point functions, where physical bound states appear as poles. Poles in the two-point function corresponding to a minimally coupled scalar, for example, correspond to $`F^2`$ glueball masses in the field theory. The analysis of the spectrum can be reduced, as usual in the AdS/CFT correspondence, to the solution of a Schroedinger problem . After a change of variable $`yz`$ to the conformally flat metric $`ds^2=e^{2\varphi (z)}((dz)^2+(dx)^2)`$ and a field redefinition $`\mathrm{\Phi }_k(z)=e^{3\varphi (z)/2}\psi (z)`$, the 5$`d`$ equation for a minimally coupled scalar $`\mathrm{\Phi }(x,y)=e^{ikx}\mathrm{\Phi }_k(y)`$ takes the Schroedinger form
$$(_z^2+V(z))\psi =E\psi $$
(20)
where $`V=\frac{3}{2}\varphi ^{\prime \prime }+\frac{9}{4}(\varphi ^{})^2`$. The eigenvalues $`E`$ give the poles in the two-point function and the spectrum.
The form of $`V`$ immediately tells us whether the theory has a mass gap and a discrete spectrum or a continuous one, whether it confines or not. Unfortunately, in very few examples $`V`$ is known along the entire flow. We can nevertheless extract some information from the IR behaviour. For the logarithmically divergent flows discussed above, if $`A<1`$, the singularity is mapped to a finite $`z_0`$ and we have
$$V\frac{3A(5A2)}{4(1A)^2(zz_0)^2}.$$
(21)
This behaviour looks potentially dangerous, but, as discussed in all quantum mechanics textbooks, $`Vk/z^2`$ has a discrete spectrum bounded from below, provided $`k1/4`$. It is easy to check that, for the logarithmically divergent flows, this condition is always satisfied. The value $`k=1/4`$ is obtained for $`A=1/4`$. This is the value that appears in many solutions where the supergravity potential is irrelevant in the IR , but also in one of the examples of N=4 coulomb branch in . If $`A>1`$, the singularity is mapped to $`z=\mathrm{}`$, the potential goes to zero and we may expect portions of continuous spectrum. Clearly, any sensible prediction about the spectrum requires the full knowledge of $`V`$. The same Schroedinger equation is to be considered when looking at generalisations of the RS scenario.
### 3.1 Supersymmetric and non-supersymmet- ric examples
We now briefly discuss few examples in the literature.
In , the class of non-supersymmetric solutions where the potential can be neglected in the IR have been discussed. They all have $`A=1/4`$. It was argued that they may exhibit a variety of IR behaviours, from confinement to screening, depending on the values of the constants $`B_a`$. Since we can not follow the solution from UV to IR, it is difficult to make more meaningful claims. We do not even know whether these solutions correspond to deformations or to different vacua of the UV fixed point.
In the N=4 Coulomb branch solutions discussed in , $`A`$ assumes various values. There is one solution with $`A=1/5`$, one with $`A=1/4`$ and all the other have $`A>1/4`$. The UV behaviour can be unambiguously determined using the first-order equations (19). All these solutions correspond to different vacua (Coulomb branch) of the UV fixed point.
The supersymmetric massive flow from N=4 to N=1 SYM was discussed in . It has $`A=1/2`$. The qualitative properties of the solution agree with QFT expectations. They are discussed in the next section.
Due to the IR singularity, not all the previous solutions are expected to be physical. A possible criterion for selecting the physical solutions has been proposed in . According to this criterion, the supergravity potential must be bounded above along the flow. This seems to eliminate all solutions with $`A<1/4`$. The case $`A=1/5`$ in the examples of N=4 Coulomb branch is indeed known to correspond to a singular 10$`d`$ solution with negative tension branes. The criterion can be also understood as follows. It selects solutions for which the IR ambiguities noticed in are absent. The action for a (canonically normalized) scalar $`S=e^{4\varphi }(\lambda )^2`$ predicts an IR contribution to the condensate
$$<O_\lambda >=\frac{\delta S}{\delta \lambda }e^{4\varphi }\lambda |yy_0|^{4A1}$$
(22)
for all logarithmic flows. This IR ambiguities diverges when $`A<1/4`$. The case $`A=1/4`$ is borderline. It is possible that, as noticed in , only the $`A=1/4`$ solutions representing vacua have a physical interpretation.
## 4 The Flow to N=1 SYM
We now present a holographic RG flow from N=4 SYM to pure N=1 SYM in the IR. We find agreement with field theory expectations: quarks confine, monopoles are screened, and there is a gau- gino condensate.
Consider a deformation of N=4 Super Yang-Mills theory with a supersymmetric mass term for the three fermions in the chiral N=1 multiplets. In N=1 notations, this is a mass term for the three chiral superfields $`X_i`$
$$d^2\theta m_{ij}\text{Tr }X_iX_j+\mathrm{c}.\mathrm{c}.,$$
(23)
where $`m_{ij}`$ is a complex, symmetric matrix.
The theory flows in the IR to pure N=1 Yang-Mills, which confines. To obtain the standard N=1 pure Yang-Mills with fixed scale $`\mathrm{\Lambda }`$, we need a fine tuning of the UV parameters, in which the mass $`m`$ diverges while the ’t Hooft coupling constant, $`x`$, goes to zero as an (inverse) logarithm of $`m`$. This is outside the regime of validity of supergravity, which requires a large $`x`$. We can think of $`m`$ as a regulator for N=1 SYM. When embedded in N=4 SYM, the theory is finite. To get a well defined N=1 SYM, we remove the cut-off ($`m\mathrm{}`$) with a fine tuning of the coupling ($`x(m)0`$). However, if we use supergravity, we are in the large $`x`$ regime. The massive modes have a mass comparable with the scale of N=1 SYM and they do not decouple. We can think of this as a theory with an ultraviolet cut-off. A good analogy is with lattice gauge theory. $`1/m`$ corresponds to the lattice spacing. The continuum limit is obtained with a fine tuning $`a0,g(a)0`$. However we can study the lattice theory at strong coupling, far from the continuum limit. A standard computation at strong coupling (by Wilson) gives the area law. We are just doing analogous computations with supergravity. Qualitative features of the theory should hold also at strong coupling.
The 5-dimensional action for the scalars
$$L=\sqrt{g}\left[\frac{R}{4}\frac{1}{24}\text{Tr }(U^1U)^2+V(U)\right],$$
(24)
is written in terms of a $`27\times 27`$ matrix $`U`$, transforming in the fundamental representation of $`E_6`$ and parametrising the coset $`E_6/USp(8)`$. In a unitary gauge, $`U`$ can be written as $`U=e^X,X=_a\lambda _aT_a`$, where $`T_a`$ are the generators of $`E_6`$ that do not belong to $`USp(8)`$. This matrix has exactly 42 real independent parameters, which are the scalars of the supergravity theory. They transform in the following $`SO(6)`$ representations: $`\underset{¯}{10_c}`$, $`\underset{¯}{20}`$, and $`\underset{¯}{1_c}`$. The supersymmetric mass term for the chiral multiplets, $`m_{ij}`$, transforms as the $`\underset{¯}{6}`$ of $`SU(3)SO(6)`$, and the corresponding supergravity mode appears in the decomposition of the $`\underset{¯}{10}\underset{¯}{1}+\underset{¯}{6}+\underset{¯}{3}`$ of $`SU(4)`$ under $`SU(3)\times U(1)`$. The term $`\underset{¯}{1}`$ in this decomposition corresponds instead to the scalar $`\sigma `$ dual to the gaugino condensate in N=1 SYM. In principle, a generic non-zero VEV for $`m_{ij}`$ will induce non-zero VEVs for other scalars as well, due to the existence of linear couplings of $`m`$ to other fields in the potential. However, if we further impose $`SO(3)`$ symmetry by taking $`m_{ij}`$ proportional to the identity matrix, a simple group theory exercise shows that all the remaining fields can be consistently set to zero. This is true also if we consider a two-parameter Lagrangian depending on both $`m`$ and $`\sigma `$. This felicitous circumstance makes an apparently intractable problem very simple and exactly solvable.
The actual computation is reported in . The result for the action for $`m`$ and $`\sigma `$ (the reason why we are considering both modes will be clear very soon) is
$`L`$ $`=`$ $`\sqrt{g}\{{\displaystyle \frac{R}{4}}+{\displaystyle \frac{1}{2}}(m)^2+{\displaystyle \frac{1}{2}}(\sigma )^2+`$ (25)
$`{\displaystyle \frac{3}{8}}[(\mathrm{cosh}{\displaystyle \frac{2m}{\sqrt{3}}})^2+4\mathrm{cosh}{\displaystyle \frac{2m}{\sqrt{3}}}\mathrm{cosh}2\sigma `$
$``$ $`(\mathrm{cosh}2\sigma )^2+4]\}.`$
The action has the supersymmetric form (17) with $`W=\frac{3}{4}\left(\mathrm{cosh}\frac{2m}{\sqrt{3}}+\mathrm{cosh}2\sigma \right)`$. The first order equations (19) read
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(1+\mathrm{cosh}{\displaystyle \frac{2m}{\sqrt{3}}}\right)`$ (26)
$`\dot{m}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2}}\mathrm{sinh}{\displaystyle \frac{2m}{\sqrt{3}}},`$ (27)
$`\dot{\sigma }`$ $`=`$ $`{\displaystyle \frac{3}{2}}\mathrm{sinh}2\sigma .`$ (28)
One interesting feature of the solution is that the equations can be analytically solved. To the best of our knowledge, there is only another example of analytically solvable flow, describing the Coulomb branch of N=4 SYM . The solution in our case is:
$`\varphi (y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}[2\mathrm{sinh}(yC_1)]+`$ (29)
$`+`$ $`{\displaystyle \frac{1}{6}}\mathrm{log}[2\mathrm{sinh}(3yC_2)],`$
$`m(y)`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{2}}\mathrm{log}\left[{\displaystyle \frac{1+e^{(yC_1)}}{1e^{(yC_1)}}}\right],`$ (30)
$`\sigma (y)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left[{\displaystyle \frac{1+e^{(3yC_2)}}{1e^{(3yC_2)}}}\right].`$ (31)
The metric has a singularity at $`y=C_1`$ with $`A=1/2`$
$$ds^2=dy^2+|yC_1|dx^\mu dx_\mu .$$
(32)
Around this point $`m`$ behaves as
$$m\frac{\sqrt{3}}{2}\mathrm{log}(yC_1)+\text{const}.$$
(33)
Here we assumed that $`C_23C_1`$, so that at the point where $`m`$ is singular, $`\sigma `$ is still finite.
Let us notice that this intuitive criterion for selecting physical solutions is in agreement with the one proposed in , which exactly selects the solutions with $`C_23C_1`$. For $`C_2>3C_1`$, $`\sigma `$ diverges first with a value $`A=1/6`$. For these and other reasons, we regard these solutions as unphysical.
### 4.1 Properties of the solution
Let us discuss the qualitative properties of the N=1 SYM solution.
It is easy to see that the solution corresponds to a true deformation of the gauge theory. Indeed, $`m`$ approaches the boundary in the UV ($`y\mathrm{}`$) as $`me^y`$, which is the required behaviour of a deformation (see eq.(16)). On the other hand, $`\sigma `$ has the UV behaviour appropriate for a condensate $`\sigma e^{3y}`$. Let us stress that this behaviour is enforced by the requirement of N=1 supersymmetry along the flow. The interpretation of the solution is therefore the following: upon perturbation with a mass term for the three chiral fields, the N=4 SYM theory flows in the IR to pure N=1 SYM in a vacuum with a non-zero gaugino condensate. The existence of a gaugino condensate is one of the QFT expectations for N=1 SYM.
We also expect the gauge theory to exhibit confinement in the IR. We can easily compute a two-point function for a minimally-coupled scalar in the background with $`\sigma =0`$. In our example, the Schroedinger potential is
$$V(z)=\frac{6\mathrm{cos}(2z)+9}{\mathrm{sin}^2(2z)}.$$
(34)
It is obvious from the figure below that there is mass gap and a discrete spectrum. The AdS boundary is at $`z=0`$ and the singularity at $`z=\pi /2`$.
The two-point function for the massless scalar corresponding to $`F^2`$ can be explicitly computed :
$$F^2(k)F^2(0)k^2(k^2+4)\text{Re }\psi (2+ik).$$
(35)
It approaches the conformal expression $`k^4\mathrm{log}k`$ in the UV and it is analytic for small $`k`$, as appropriate for a confining theory. It has poles for $`M^2=k^2=n^2,n=2,3,\mathrm{}`$, corresponding to the $`F^2`$ glueball states in the spectrum.
Despite the presence of a singularity that invalidates the supergravity approximation in the IR, the qualitative properties of the solution agree with the QFT expectations. There is however a disturbing point: our solution depends on two independent parameters $`C_1`$ and $`C_2`$. The first one fixes the position of the singularity and it is related to the magnitude of the mass deformation. The second one is instead related to the magnitude of the gaugino condensate. We have a chirally-symmetric vacuum and, more disturbing, a continuous degeneracy of vacua with arbitrary small condensate. We certainly expect that the correct treatment of the singularity and its resolution in string theory fixes the relation between $`C_1`$ and $`C_2`$ in agreement with field theory expectations. We do not still known how to resolve or deal with the singularity, therefore we limit ourself to a brief discussion of the QFT expectations and possible interpretations of the singularity.
### 4.2 QFT and string expectations
Strong coupling QFT results for N=1 SYM have been recently obtained and differ considerably from the weak coupling ones . At weak coupling, spontaneous breaking of the chiral symmetry $`Z_N`$ gives $`N`$ vacua that only differ for the phase of the gaugino condensate $`<\lambda \lambda >e^{2\pi ik/N}\mathrm{\Lambda }_{\mathrm{N}=1}^3`$. In the large $`N`$ limit, we obtain a circle of vacua. The magnitude of the gaugino condensate is fixed in terms of the SYM scale $`\mathrm{\Lambda }_{\mathrm{N}=1}me^{1/3Ng^2}`$. At strong coupling instead, it was shown in that there is, at least for $`\theta =0`$, a distribution of vacua with condensate $`<\lambda \lambda >m^3x^3/j^2,j=1,2,\mathrm{}`$ with zero phase. The weakly coupled circle is lost, the condensate magnitude is not fixed and the vacua have an accumulation point at the origin (zero condensate). However, we notice that the structure of vacua found in has many similarities with our supergravity result. As independently noticed in , it is tempting to identify the solution with $`C_2=3C_1`$ with the $`j=1`$ vacuum in . The other solutions with $`C_2<3C_1`$ should correspond to the $`j1`$ vacua. To see how the continuum of vacua in supergravity is reduced to a discrete numerable set, we should understand how to include string corrections in our computation. Notice that the solution with $`\sigma =0`$, which is not appealing on the ground of weak coupling intuition, could be nevertheless used as a (reasonable?) approximation for the many vacua with small condensate at strong coupling.
It was also proposed in to fix the relation between $`C_1`$ and $`C_2`$ by considering the finite temperature version of our solution, where conditions to be imposed at the horizon fix the parameters. One finds $`C_2=3C_1`$. This is the only special value for our parameters, since, exactly for $`C_2=3C_1`$, the two scalars $`m`$ and $`\sigma `$ diverge at the same point in $`y`$. In SYM the breaking of supersymmetry will select the vacuum with minimal energy. At weak coupling, where all the vacua have a condensate with the same magnitude, this procedure should give us also the value of the N=1 condensate. At strong coupling, with condensates of almost arbitrary magnitude, this would give information at most about one particular vacuum ($`j=1`$?).
The knowledge of the full 10 dimensional solution would greatly help in understanding the properties of the RG flow and in studying possible resolutions of the singularity. It may even happen that the singularity is an artifact of the dimensional reduction, that disappears in 10$`d`$. This happens, for example, in the case of the Coulomb branch of N=4 SYM , where the 10 dimensional background is just a regular continuous distribution of D3-branes. However, even in this context, some other equally nice<sup>10</sup><sup>10</sup>10But not satisfying the criterion in . 5$`d`$ solutions have a lift to still singular 10$`d`$ solutions, representing D3-branes with negative tension. The complete ansatz for the 10$`d`$ lifting of 5$`d`$ solutions is known only for a subset of scalars, the $`\underset{¯}{20}`$, coming from the KK modes of the internal metric. This is sufficient to lift all solutions representing the Coulomb branch, but it is not of help with our solution, where the modes $`\underset{¯}{10}`$ from the anti-symmetric tensors are excited.
A ten dimensional interpretation of the N=1 solution in terms of a background with also D5-branes has been proposed in . We only notice that the ingredients in this interpretation (D5 and NS-branes) have been independently suggested in on the basis of the strong coupling QFT analysis.
Finally, we mention that a mechanism for resolving singularities in distributions of branes which may help, after the 10$`d`$ lifting, has been proposed in .
### 4.3 The Wilson loop
A complementary approach for checking confinement is the computation of a Wilson loop, which should manifest an area law behaviour. We need to minimise the action for a string whose endpoints are constrained on a contour $`C`$ on the boundary. The detailed computation is reported in . In the coordinates used in those papers, the quark-antiquark energy reads
$$E=S/T=𝑑x\sqrt{(_xu)^2+f(u)}.$$
(36)
where $`f(u)=T^2(u)e^{4\varphi (u)}`$. The phase of the theory can be inferred by the IR behaviour of this function (see for a review of the various cases). $`T(u)`$ is the tension of the fundamental (in the case of a quark loop) or of the D1 string (monopole) in five dimensions. They are in general non-trivial functions of the scalar fields. The 5$`d`$ N=8 gauged supergravity has an $`SL(2,Z)`$ symmetry that allows to discriminate electric and magnetic strings. They should couple to the 5$`d`$ antisymmetric tensors $`B_{\mu \nu }^{I\alpha }`$, transforming in the $`(\underset{¯}{6},2)`$ of $`SO(6)\times SL(2,Z)`$. The $`SO(6)`$ index should account for the orientation of the strings on the five-sphere, while the $`SL(2,Z)`$ index should iden- tify electric and magnetic quantities. On the basis of naive dimensional reduction from ten dimensions, the tensions can be read from the coefficients of the kinetic term for the antisymmetric tensors. In 10 dimensions, the tension of the fundamental string (or the D1-string) can be read from the NS-NS (or R-R) antisymmetric tensor Lagrangian evaluated in the Einstein frame,
$$\frac{1}{T_{F1}^2}H_{NS\text{-}NS}^2+\frac{1}{T_{D1}^2}H_{R\text{-}R}^2.$$
(37)
A simple Weyl rescaling shows that this property is valid also in the five-dimensional theory in the Einstein frame.
The kinetic terms for the anti-symmetric ten- sors can be computed for the N=1 SYM solution and behave asymmetrically in the $`SL(2,Z)`$ indices . The final result for the tensions $`T(u)`$ of the fundamental strings and of the D1-strings are, respectively,
$`T_{F1}^2`$ $`=`$ $`4\left(\mathrm{cosh}{\displaystyle \frac{4m}{\sqrt{3}}}+\mathrm{cosh}{\displaystyle \frac{2m}{\sqrt{3}}}\right),`$ (38)
$`T_{D1}^2`$ $`=`$ $`8\left(\mathrm{cosh}{\displaystyle \frac{m}{\sqrt{3}}}\right)^2,`$ (39)
so that the asymptotic behaviour of the corresponding functions $`f(u)`$ is
$$f_{(q\overline{q})}(u)1,f_{(m\overline{m})}(u)\left|uC_1\right|.$$
(40)
It is easy to check that $`f_{(q\overline{q})}(u)`$ is bounded from below. It follows that the energy $`EcL`$, where $`L`$ is the quark distance. It can be easily proven that it is in fact $`E=cL`$, implying an area law behaviour for the Wilson loop, as expected for a confining theory. The IR behaviour of $`f_{(m\overline{m})}(u)`$ implies, on the other hand, that monopoles are screened (see for a review).
There is an apparent contradiction in the previous reasoning. The 5$`d`$ dilaton is not running in our solution. If the 10$`d`$ dilaton were also constant, the tension for a fundamental string would be proportional to the tension of a D1-string and the same would be true also after dimensional reduction to 5 dimensions. The 5$`d`$ tensions would be then complicated functions of the scalars, but invariant under $`SL(2,Z)`$. We instead find an $`SL(2,Z)`$ asymmetric result from the N=8 gauged supergravity evaluated along our solution. A possible way out is to assume that, against naive expectations, the 10$`d`$ dilaton is not constant. Clearly, it also exists the option that the 10$`d`$ dilaton is constant and that the argument which determines the 5$`d`$ tensions via dimensional reduction is too naive. However, we are not aware of any argument that rules out the possibility of a running 10$`d`$ dilaton. Since we are not expert in reconstructing 10$`d`$ solutions from 5$`d`$ ones, we just limit ourselves to consider this option and perform some very preliminary check on the equations of motion.
The 10$`d`$ dilaton equation of motion is
$$^2\varphi G_{MNP}G^{MNP}.$$
(41)
Therefore, a non-vanishing anti-symmetric tensor is a source for the dilaton. We can perform a check on our solution at the linearised level. Consider a generic fluctuation of the anti-symmetric tensor $`B_{ab}=f_I(y)Y_{[ab]}^{I\pm }`$. We refer to for notations and useful equations. Here $`Y_{[ab]}^{I\pm },a,b=1,\mathrm{},5`$ are harmonic functions on the five-sphere, transforming in the representation $`I`$ of $`SO(6)`$. They satisfy $`ϵ_{abcde}_cY_{[de]}=\pm 2i(k+2)Y_{[ab]}`$, where $`k`$ is an integer labelling the harmonic degree. It is then easy to check that
$$^2\varphi \frac{1}{3}((_yf)^2(k+2)^2f^2)Y_{[ab]}Y_{[ab]}.$$
(42)
In our case ($`I=\underset{¯}{10}`$) $`k=1`$. Since we are considering a deformation of the UV fixed point, $`fe^x`$, we see that the dilaton must run. Notice that instead, considering a different vacuum of the UV theory, one has $`fe^{3x}`$, and the dilaton remains constant (at least at the first perturbative order).
We still need to check that $`Y_{[ab]}Y_{[ab]}0`$. There is at least one example where $`Y_{[ab]}Y_{[ab]}=0`$: the $`SU(3)\times U(1)`$ critical point of the N=8 supergravity, whose 10$`d`$ solution is explicitly known . In the product $`\underset{¯}{10}\times \underset{¯}{10}=\underset{¯}{20}+\mathrm{}`$, only the indicated term contains scalar terms ($`SO(5)SO(6)`$ singlets). It is easy to check that, decomposing $`\underset{¯}{10}=\underset{¯}{1}+\underset{¯}{3}+\underset{¯}{6}`$ under $`SU(3)\times U(1)`$, the $`\underset{¯}{1}`$ term (related to the $`SU(3)\times U(1)`$ critical point) has vanishing square. The N=1 mass term $`\underset{¯}{6}`$, however, has non vanishing square.
This argument is certainly not a proof that the 10$`d`$ dilaton runs. However, we find this option appealing. A running of the 10$`d`$ dilaton would agree with an interpretation of our solution that includes branes others than the D3s. In many respects, the knowledge of the explicit 10$`d`$ solution would help us in understanding the system, from the constituent branes to the fate of the singularity. Using a D3-brane probe in the 10$`d`$ background we could also explicitly compute the running of the gauge coupling along the flow.
###### Acknowledgments.
The content of this paper has to appear divided in two parts in the proceedings for the TMR conference in Paris, September 99. For sake of economy and not to bother the potential readers, we just unified the two parts in the version for the archives. We would like to thank our collaborators L. Girardello and M. Porrati, with whom most of the results reported here were obtained. We also thank D. Anselmi for useful discussions and collaboration at various stages. We also thank N. Dorey, S. S. Gubser, S. P. Kumar, N. Warner, C. Pilch and E. Witten for useful discussions and criticisms. A. Z. is partially supported by INFN, and by the European Commission TMR program ERBFMRX-CT96-0045, wherein he is associated to the University of Torino. M. P is partially supported by INFN, MURST, the European Commission TMR program ERBFMRX-CT96-0045, wherein she is associated to Imperial College, London, and the PPARC SPG grant PPA/G/S/1998/00613. |
warning/0002/cond-mat0002230.html | ar5iv | text | # Charge-Reversal Instability in Mixed Bilayer Vesicles
## I Introduction
The self-assembly of colloidal particles offers an attractive route to the synthesis of highly ordered, nanostructured materials. Typically these materials have been extremely soft, being stabilized by entropic effects.
For example, classical colloidal crystals are three-dimensional arrays of mutually repelling spheres . Entropic effects maintain their crystalline order in spite of a density well below that of close-packing. As a result, these arrays are easily disrupted by small mechanical shear, dilution, etc. More recently, depletion forces have been harnessed to assemble spheres into crystalline arrays on the walls of their container . Again the physical forces between the spheres are repulsive, and again the resulting arrays are extremely soft.
Attempts to create strong ordered materials from physically attracting components have generally produced instead highly disordered aggregates. Recently, however, Ramos et al. reported the observation of robust two-dimensional crystallites formed from negatively-charged latex spheres introduced into a suspension of bilayer vesicles . The membranes forming the vesicles consist of a mixture of positively-charged and neutral surfactants. The immense electrostatic attraction between the negative spheres and positive membranes led to the crystallites’ great strength; their ordered 2d character arose via the intermediary role of the vesicles as templates for the initial self-assembly of the spheres.
In this paper we develop some of the physics of the crucial intermediate step just mentioned, elaborating and extending the discussion in . This stage begins when the latex spheres are first introduced to the vesicle suspension, and lasts for hours to days. Initially the spheres adsorb avidly onto the vesicles, and indeed many vesicles become completely covered with spheres. However, a significant subpopulation of vesicles content themselves with only partial coverage: on these vesicles the adsorbed spheres form a self-limiting ‘raft’. Once the raft forms, no further spheres attach to the vesicle anywhere, though they are present in excess. Instead, particles in suspension are seen to approach, then wander away from, the vesicle.
The theory of colloidal surface interactions is vast (for introductions see ). Our goal is to introduce a very simple mechanism for adhesion saturation, summarized graphically in Fig. 1 below, then present some calculations to show how it works in the parameter regime relevant to experiments. We will argue that our effect should be qualitatively unchanged after many other surface-interaction effects are included in the analysis, but much work remains to be done to show this in detail. Sect. II sketches the physics of our mechanism. Sect. III begins the analysis using linearized Poisson–Boltzmann theory, considering in turn a series of more complicated situations. The linearized theory is familiar and helps to connect the analysis to the physical picture, but it proves to be inadequate for the interesting range of parameter values. Thus in Sect. IV we upgrade to the full nonlinear theory, which proves to be quite easy in this context. Finally we consider the effects of ion correlations, neglected in Poisson–Boltzmann theory, in Sect. V. A glossary of symbols appears in the Appendix.
## II Physical Picture
We first briefly review the physical picture developed in and summarized in Fig. 1.
Consider first two dielectric surfaces bearing fixed charge densities $`\sigma _\pm `$ of the same magnitude but opposite sign in an electrolyte solution. When they are separated by several screening lengths they feel little mutual attraction, since each maintains a neutralizing cloud of counterions. As the surfaces approach closer, eventually their screening clouds begin to interpenetrate. Then negative counterions from the positive surface, and positive counterions from the negative surface, can escape to infinity without violating overall charge neutrality. The corresponding gain in entropy reduces the system’s free energy: counterion release drives the surfaces into contact.
Next consider the case of two surfaces of opposite sign and unequal magnitude; for instance, suppose that $`\sigma _+<|\sigma _{}|`$. In this case counterion release will be incomplete; after exhausting all the negative counterions, some positive ones will remain, trapped by the requirement of charge neutrality. The osmotic pressure of the trapped ions will prevent the surfaces from coming into perfect contact. If one surface has variable charge density, say $`\sigma _+`$, then additional surface charges will be pulled into the contact region in order to improve the contact with the approaching negative surface . A surface charge density can for instance vary because the composition of the surface is variable: for instance, the surface may be a mixture of charged and neutral surfactants, as in the experiments of . In this case the recruitment of charge to the contact region will deplete the other regions, in turn rendering them less attractive to additional negative dielectric objects. Fig. 1c depicts this situation: surfactant rearrangement in the outer monolayer of the membrane has permitted the release of two more ion pairs than would otherwise (panel b) have been possible.
The rearrangement of membrane charges is limited: the relative concentration of charged surfactants cannot exceed unity. The maximum charge density on the outer monolayer may still be less than that of the approaching dielectric, and so the final contact may still be imperfect, as shown in Fig. 1c. We will assume this to be the case in the rest of this paper.
Nevertheless, a further reduction in free energy density from Fig. 1c is still possible, once we remember that the inner membrane monolayer and its counterions need not play a passive role. Fig. 1d shows how the remaining trapped counterions in panel c can leave the gap, even if the membrane is impermeable, by following the dashed horizontal arrows in panel c. After this rearrangement some of the charge on the negative dielectric is neutralized by surfactants on the inner monolayer, whose own interior counterions migrate to the nonadhesion region.<sup>*</sup><sup>*</sup>*Even if the dielectric’s charge exceeds twice the monolayer charge density, as assumed in the text below, additional $`\pm `$ ion pairs can be brought from the membrane interior, with the positive ions remaining in the adhesion region to help neutralize the dielectric and the negative ones migrating to the nonadhesion region, driving its net charge still more negative. Panel c also shows a rearrangement of the surfactants on the inner monolayer, further depleting the charge of the noncontact zone.
Fig. 1d raises an intriguing question: will the migration of interior counterions ever overwhelm and effectively reverse the charge of the membrane as seen from outside, as shown in the figure? Of course, cartoons alone will not settle this question, but we can argue physically that such an effect may well happen as follows. First we note that the positive charge density of the noncontact zone “n” is already very small in Fig. 1c, since the interior monolayer and its counterion cloud cancel, and as we will see below the electrostatic interaction driving the depletion of charged surfactants from the outer monolayer is very strong. Thus only a small migration of interior counterions will suffice to get charge reversal. Second, the entropic cost of creating a nonuniform charge density in the interior counterion cloud is quadratic in the amount of charge which migrates, since the uniform distribution is an equilibrium state. But the free energy gain from this redistribution is linear in the amount of charge migration, being dominated by the derivative $`\mathrm{d}/\mathrm{d}\sigma _+`$ of the attractive self-energy (see formula (10) below). Thus a finite amount of counterion migration will occur, and this amount may well exceed the small net charge on the nonadhesion region, effectively reversing it.
The rest of this paper is devoted to a quantitative justification of the intuitive argument just given. Before passing on to the analysis, we should remark on another feature of Fig. 1d. Charge reversal requires that electric fields (represented schematically by the vertical lines in the figure) penetrate the interior of the membrane. Since the membrane interior is a low dielectric constant medium, the energetic cost of these fields can be significant, another term quadratic in the amount of charge migration from panel c to d. If the membrane is sufficiently thick, this cost will reduce the charge migration below the point of charge reversal, a point we will need to examine quantitatively in Sect. IV C below.
## III Linearized Mean-Field Theory
In this section we begin the mathematical implementation of the ideas in Sect. II. We begin with the linearized (Debye-Hückel) limit of low charge density, even though ultimately we will argue that the experiments studied here require a full nonlinear treatment. We do this partly because of the simplicity of the formulæ, and partly to make contact with earlier work.
To fix notation and keep the article self-contained we begin by rederiving some key results from . The Appendix summarizes our units and all symbols used throughout the paper.
### A Basic formulæ
The electrostatic potential energy of a distribution of free charges of density $`\rho (𝐫)`$ is $`\frac{1}{2}d𝐫\rho (𝐫)\psi (𝐫)`$, where $`\psi `$ is the electric potential . The potential created by a single point charge $`q`$ in an infinite, uniform, dielectric medium is $`\psi (𝐫)=q/4\pi ϵ|𝐫|`$. In a more complicated situation, $`\psi (𝐫)`$ is related to $`\rho (𝐫^{})`$ by some Green function $`G(𝐫,𝐫^{})`$ and obeys Poisson’s equation, $`^2\psi =\rho /ϵ`$.
We first imagine a uniform charge distribution of density $`\sigma `$ on the surface $`\{z=0\}`$.The assumption of fixed charge is appropriate for surfaces with fully-ionized groups at the pH used, such as those in the experiments of . We also implicitly assume that the surfactants used are insoluble in water, so that their numbers in the membrane are fixed. This assumption may need further scrutiny, since in the experiments one surfactant species forms micelles. The halfspace $`z<0`$ is filled with a dielectric with no free charges, and so the electric field must everywhere vanish here. The other halfspace is a univalent salt solution in equilibrium with a reservoir at concentration $`\widehat{n}`$. The reservoir must remain neutral, but it can supply ion pairs at a cost in free energy given by a chemical potential $`\mu k_\mathrm{B}T`$. The total free energy of the mobile ions near the surface is then $`F=F_{\mathrm{ent}}+F_{\mathrm{es}}`$, where the entropic and electrostatic energies in mean-field approximation are
$$F_{\mathrm{ent}}=k_\mathrm{B}Td𝐫\left[n_+(\mathrm{ln}n_+v_01)+n_{}(\mathrm{ln}n_{}v_01)\mu (n_++n_{})+\xi (n_+n_{}+n_\mathrm{f})\right]$$
(1)
$$F_{\mathrm{es}}=\frac{1}{2}d𝐫d𝐫^{}en_{\mathrm{tot}}(𝐫)G(𝐫,𝐫^{})en_{\mathrm{tot}}(𝐫^{}).$$
(2)
In the above formulæ, $`n_\pm `$ are the number densities of ions, while $`n_{\mathrm{tot}}(𝐫)=\rho /e=n_+n_{}+n_\mathrm{f}`$ is the total signed density, including fixed surface charges with signed density $`n_\mathrm{f}`$. We introduced a Lagrange multiplier $`\xi `$ to enforce overall neutrality. The symbol $`v_0`$ is a microscopic volume factor which will drop out of all physical results. We have fixed the arbitrary constant in $`F_{\mathrm{es}}`$ by setting the electrostatic energy to zero when the mobile counterions form a sheet coinciding with the fixed surface charge. Thus $`F_{\mathrm{es}}`$ is the work needed to pull this sheet away from the surface, and so is a positive quantity.
In equilibrium we have $`\frac{\delta F}{\delta n_\pm (𝐫)}=0`$. Away from the plane this fixes
$$n_\pm (𝐫)v_0=\mathrm{e}^{\mu (\overline{\psi }(𝐫)+\xi )},z>0.$$
(3)
Here $`\overline{\psi }=e\psi /k_\mathrm{B}T`$ and we have fixed the additive constant in $`\psi `$ by choosing $`\psi (\mathrm{})=0`$. Since $`n_+=n_{}=\widehat{n}`$ at infinity, we get $`\xi =0`$ and $`\mu =\mathrm{ln}\widehat{n}v_0`$, or
$$n_\pm (𝐫)=\widehat{n}\mathrm{e}^{\overline{\psi }(𝐫)}.$$
(4)
Substituting then gives the free energy
$$F=k_\mathrm{B}T\widehat{n}d𝐫\left[\overline{\psi }\mathrm{sinh}\overline{\psi }2\mathrm{cosh}\overline{\psi }+\frac{1}{2}(n_\mathrm{f}/\widehat{n})\overline{\psi }+2\right].$$
(5)
The last term of (5) is a constant which we have added by hand to cancel a term proportional to the volume of the world.
Eqn. (5) simplifies if the dimensionless potential $`\overline{\psi }`$ is everywhere $`1`$; in this case we simply get $`F=k_\mathrm{B}Td𝐫\frac{1}{2}n_\mathrm{f}\overline{\psi }`$. Since the fixed charge $`n_f`$ is confined to a plane, the free energy is a purely surface term once $`\overline{\psi }`$ has been found.
To find $`\overline{\psi }`$, we note that it satisfies the Poisson equation, a property of the Green function used to define it. Using the charge density $`en_\pm (𝐫)`$ found above in (4) gives the Poisson-Boltzmann equation,
$$^2\overline{\psi }=\frac{2e^2\widehat{n}}{ϵk_\mathrm{B}T}\mathrm{sinh}\overline{\psi }.$$
(6)
Linearizing then gives the familiar Debye-Hückel equation: $`^2\overline{\psi }=\kappa ^2\overline{\psi }`$, where $`\kappa =\sqrt{2e^2\widehat{n}/ϵk_\mathrm{B}T}`$.
The objects we want to consider are much bigger than the screening length $`\lambda _D=1/\kappa `$ (see Fig. 1a). Thus our geometry is essentially planar, and we need the planar solutions $`\overline{\psi }(z)=B\mathrm{e}^{\pm \kappa z}`$ to the Debye-Huckel equation. The electric field is then $`𝐄=\mathbf{}\psi `$, which indeed decays exponentially on the length scale $`\lambda _D`$.
For a single wall we must choose the decaying solution to (6). We fix the constant $`B`$ by imposing Gauss’s law at the surface: $`𝐄=\frac{\psi }{z}\widehat{𝐳}=\frac{\sigma }{ϵ}\widehat{𝐳}`$. Then $`B=\sigma e/\kappa ϵk_\mathrm{B}T`$, the solution is
$$\overline{\psi }(z)=\frac{\sigma e}{\kappa ϵk_\mathrm{B}T}\mathrm{e}^{\kappa z}.(\mathrm{linearized}\mathrm{approximation})$$
(7)
and the free energy per unit area of the isolated, charged surface is
$$f_{\mathrm{self}}F/(\mathrm{area})=k_\mathrm{B}T\sigma B/2e=\sigma ^2/2\kappa ϵ.(\mathrm{linearized}\mathrm{approximation})$$
(8)
Another well-known solution to (6) arises in the opposite case of very high charge density, where $`\overline{\psi }1`$ at the surface. In this case the Poisson-Boltzmann equation has a solution of “Gouy-Chapman” form: $`\overline{\psi }(z)=\mathrm{ln}\left[\frac{2ϵk_\mathrm{B}T}{e^2\widehat{n}}\frac{1}{(z+\lambda _{\mathrm{GC}})^2}\right]`$. Here the free parameter is the offset $`\lambda _{\mathrm{GC}}`$, chosen to enforce Gauss’s law: $`\lambda _{\mathrm{GC}}=2ϵk_\mathrm{B}T/e\sigma `$. More highly-charged surfaces thus have smaller $`\lambda _{\mathrm{GC}}`$ and so a more nearly singular potential. The pathological behavior of $`\overline{\psi }`$ at large $`z`$ simply reflects the end of the regime $`\overline{\psi }1`$ at large enough $`z`$. Note that the electric field $`E_z=2k_\mathrm{B}T/e(z+\lambda _{\mathrm{GC}})`$ of the Gouy-Chapman solution is independent of the ambient salt concentration $`\widehat{n}`$, as it should be: the electric forces near a highly charged surface depend only on the surface charge. The salt concentration determines only the extent of the region in which the strong-field approximation is valid.
### B Two dielectrics
We minimized the free energy of an isolated surface, obtaining (8). To extract any useful work from this stored free energy, we would have to remove some constraint. One way to do this is to bring in another semiinfinite, planar dielectricNothing is really infinite. The phrase “semiinfinite planar dielectric” will mean a finite dielectric object whose surface curvature is much smaller than $`\kappa `$, whose interior contains no free charges, and whose volume is large enough that any interior electric field would be prohibitively expensive in energy. As two such objects approach, the gap $`\mathrm{}`$ between them decreases but the total volume occupied by solution doesn’t change; this is why the constant we subtracted from (5) really is a constant. bearing opposite surface charge, thus changing the solution region from a half-space to a planar slab of thickness $`\mathrm{}`$. Let us suppose that a surface with $`\sigma _+>0`$ approaches another surface with $`\sigma _{}<0`$.
Parsegian and Gingell studied this situation in the linearized approximation , arguing as in Sect. II that the surfaces attract via counterion release until all of one species of counterions in the gap (the “minority” species) has been exhausted. If $`\sigma _+|\sigma _{}|`$, a residual cloud of the other (“majority”) species remains in the gap and the system equilibrates at a finite gap spacing $`\mathrm{}_{}`$. Nevertheless, the final state has less free energy per unit area than it did originally; the difference is the adhesion strength $`W`$.
We could compute $`W`$ by again solving a boundary-value problem as in Sect. III A, but there is a shortcut. Suppose that $`\sigma _+<|\sigma _{}|`$, so that the $`+`$ counterions are the “majority” species. In mechanical equilibrium the hydrostatic pressure pushing the walls together vanishes. The planar Poisson-Boltzmann equation is a second-order ordinary differential equation, and so its solutions form a two-parameter family. One integration constant is fixed by the Gauss-law boundary condition on the negative wall, while in equilibrium the other is fixed by the condition of vanishing pressure. Hence the solution $`\overline{\psi }(z)`$ is exactly the same for two walls as it is for the isolated negative wall; the only difference is that in the former case we truncate the solution at $`z=\mathrm{}_{}`$, while in the latter case $`z`$ extends to infinity. The equilibrium gap spacing $`\mathrm{}_{}`$ is then just the value of $`z`$ at which Gauss’s law for the positive wall is satisfied: $`(\frac{\psi }{z})=\frac{\sigma _+}{ϵ}`$. Then (7) gives the equilibrium spacing $`\mathrm{}_{}`$ by $`\mathrm{e}^\kappa \mathrm{}_{}=|\sigma _{}/\sigma _+|`$ in the linearized approximation. Note that indeed the right side of this formula is positive and greater than unity, as it must be since $`\mathrm{}_{}0`$.
We now recall that the linearized approximation retains only the boundary term of (5), so
$`f_{\mathrm{gap}}(\sigma _+,\sigma _{})`$ $`=`$ $`{\displaystyle \frac{k_\mathrm{B}T}{2e}}\left[\sigma _{}\overline{\psi }(0)+\sigma _+\overline{\psi }(\mathrm{}_{})\right]`$ (9)
$`=`$ $`{\displaystyle \frac{1}{2\kappa ϵ}}\left((\sigma _{})^2(\sigma _+)^2\right).`$ (10)
Repeating these steps for the opposite case where $`\sigma _+>|\sigma _{}|`$, we find that in general (Fig. 2a)
$$f_{\mathrm{gap}}(\sigma _+,\sigma _{})=|f_{\mathrm{self}}(\sigma _+)f_{\mathrm{self}}(\sigma _{})|.$$
(11)
Remarkably, the simple combination formula (11) will continue to hold in the full nonlinear Poisson-Boltzmann treatment of Sect. IV A below.<sup>§</sup><sup>§</sup>§Behrens and Borkovec have independently used this fact to simplify the study of nonlinear PB solutions . Formula (11) is certainly reasonable: when $`\sigma _+=|\sigma _{}|`$ all counterions get released, the two surfaces coincide, and this was our reference state of zero energy. Also, when we reverse the signs of all the charges the free energy should not change; (11) has this property.
Finally we find the adhesion energy $`W`$ as
$$W=f_{\mathrm{self}}(\sigma _+)+f_{\mathrm{self}}(\sigma _{})f_{\mathrm{gap}}(\sigma _+,\sigma _{})=\mathrm{min}\{(\sigma _+)^2,(\sigma _{})^2\}/ϵ\kappa .\text{(linearized approximation)}$$
(12)
Note that $`W`$ is completely independent of the majority charge density, a property noted by Nardi et al. In light of the physical picture in Sect. II, we can readily interpret that fact: The total counterion release is limited by the smaller of the two counterion populations.
Since $`W`$ is always positive we find, as expected, that oppositely-charged dielectrics always attract via the counterion-release mechanism . Of course this is not the behavior we were seeking to explain (see Sect. I). We must now proceed to generalize the above arguments, incorporating the relevant differences between the above system and the one studied in the experiments of .
### C Thick membrane
We just found that two oppositely-charged dielectrics attract, as expected. But the experiments we are studying involve dielectric (latex) spheres interacting not with other dielectrics, but with a bilayer membrane. In this subsection we begin to incorporate the new physics associated with this situation. We first study the interaction of a dielectric of fixed charge density $`\sigma _{}<0`$ with a positively-charged, very thick, membrane, recapitulating some results of Nardi et al. .
The new physical feature of this situation is that the bilayer membranes in the experiments are fluid mixtures of positively-charged and neutral surfactants. This means that the charge density $`\sigma _+`$ on the membrane is not a fixed number, but may vary subject to $`\sigma _+>0`$ and the overall constraint that the total membrane charge $`dA\sigma _+`$ is fixed. Let $`\sigma _{+,\mathrm{av}}`$ denote the average charge density, so that the total membrane charge is $`A\sigma _{+,\mathrm{av}}`$. In addition we will suppose that the charge density cannot exceed a maximum of $`\sigma _{\mathrm{max}}=2e/a_0`$ determined by the area per headgroup $`a_0`$ of the charged surfactants in each of the two monolayers constituting the membrane, and that $`|\sigma _{}|>\sigma _{\mathrm{max}}`$.
Throughout this paper we will adopt a highly simplified, generic picture of membrane compositional changes, retaining only the entropy of mixing of the two surfactant types. Thus we neglect other entropic or enthalpic packing effects in the assumed membrane free energy $`f_\mathrm{m}`$. Moreover, at first we will for simplicity neglect the bilayer structure of the membrane; later on we will use formulæ appropriate to a bilayer. With these simplifications $`f_\mathrm{m}`$ takes the form
$$f_\mathrm{m}=\frac{2}{a_0}k_\mathrm{B}T\left[\frac{\sigma _+}{\sigma _{\mathrm{max}}}\mathrm{ln}\frac{\sigma _+}{\sigma _{\mathrm{max}}}+\left(1\frac{\sigma _+}{\sigma _{\mathrm{max}}}\right)\mathrm{ln}\left(1\frac{\sigma _+}{\sigma _{\mathrm{max}}}\right)\right].$$
(13)
As discussed in Sect. II, we wish to explore the possibility of a spontaneous partition of the membrane into two uniform regions, which we will call zones “a” and “n”. (Ultimately we hope to find that “a” is adhering while “n” is nonadhering, but for the moment these are arbitrary names.) The areas of the two zones are not known in advance, but they must add up to the total area $`A`$, so we take them to be $`\gamma A`$ and $`(1\gamma )A`$ respectively.
The two zones exchange one conserved quantity, namely membrane charge.A second conserved quantity, the total charge of the counterions, is not independent but instead fixed by charge neutrality. Thus the system can divide into zones of charge density $`\sigma _+^{(a)}`$ and $`\sigma _+^{(n)}`$, subject to
$$\gamma \sigma _+^{(a)}+(1\gamma )\sigma _+^{(n)}=\sigma _{+,\mathrm{av}}.$$
(14)
This separation will be energetically advantageous if the corresponding total free energy $`A\left(\gamma f(\sigma _+^{(a)})+(1\gamma )f(\sigma _+^{(n)})\right)`$ is less than $`Af(\sigma _{+,\mathrm{av}})`$. Here the free energy density $`f(\sigma _+)`$ is computed for a uniform zone with fixed membrane charge density $`\sigma _+`$, minimizing over all other variables.
The instability just described will not occur if $`f(\sigma _+)`$ is a convex function, i.e. $`\mathrm{d}^2f/\mathrm{d}\sigma _{+}^{}{}_{}{}^{2}>0`$. If $`\mathrm{d}^2f/\mathrm{d}\sigma _{+}^{}{}_{}{}^{2}`$ is negative anywhere within the allowed region $`0<\sigma _+<\sigma _{\mathrm{max}}`$, we apply the Maxwell construction from thermodynamics to the graph of $`f`$. This involves drawing a straight line tangent to the graph and spanning the region of concavity. Let the two points of tangency be located at $`\sigma _+^{(a)}`$ and $`\sigma _+^{(n)}`$. If the average membrane composition $`\sigma _{+,\mathrm{av}}`$ lies between these two values, then the uniform system will be unstable to partitioning into two zones with compositions $`\sigma _+^{(a)}`$ and $`\sigma _+^{(n)}`$.
In the case of a thick membrane, we have $`f=f_{\mathrm{gap}}(\sigma _+)+f_\mathrm{m}(\sigma _+)`$. Consulting (11), (8), and (13), we find that the first (electrostatic) term is destabilizing, while the second (entropic) term is stabilizing. For future use we introduce two convenient abbreviations, one parameterizing the relative strengths of the two terms, the other a dimensionless measure of charge density:
$$\beta 2\widehat{n}a_0/\kappa =\kappa a_0/4\pi \mathrm{}_B,\overline{\sigma }\sigma /\sigma _{\mathrm{max}}.$$
(15)
With these abbreviations we obtain
$`f={\displaystyle \frac{\sigma _{\mathrm{max}}^{}{}_{}{}^{2}}{2ϵ\kappa }}\left[|(\overline{\sigma }_{})^2(\overline{\sigma }_+)^2|+\beta \left(\overline{\sigma }_+\mathrm{ln}\overline{\sigma }_++(1\overline{\sigma }_+)\mathrm{ln}(1\overline{\sigma }_+)\right)\right].(\mathrm{linearized}\mathrm{approximation})`$
Nardi et al. pointed out that this function has an inflection point, giving a region of instability (Fig. 2b), when $`\beta <1/2`$. According to (15), this means that either the maximum charge density $`e/a_0`$ must be large, or else the salt concentration $`\widehat{n}`$ very small. Substituting some typical values for the charge per headgroup $`a_0=0.5`$nm and salt concentration $`\widehat{n}=1`$mM$`=0.0006`$nm<sup>-3</sup> gives $`\beta =0.006`$, well into the regime of instability.
Though we have found an instability, two remarks limit its interest. First, we have insisted that $`0<\sigma _+<\sigma _{\mathrm{max}}`$, so of course the charge densities $`\sigma _+^{(a)}`$ and $`\sigma _+^{(n)}`$ on our two zones are both positive: both zones are adhesive, unlike the experimental phenomenon we are trying to explain.
Moreover, we found no instability at all unless the charge density $`\sigma _{\mathrm{max}}`$ is quite large (recall also that $`|\sigma _{}|`$ is assumed to be even greater than this). But at such large charges our linearized approximation breaks down! Our calculation becomes inconsistent just as it gets interesting. Much of this paper is dedicated to correcting this deficiency. We ask the reader to suspend disbelief momentarily while we implement the physical picture sketched in Sect. II in the linearized theory, where the formulæ are simple. Our claim is that the physical picture is robust and holds beyond this inadequate mathematical framework; we will support this claim by improving the calculation in Sect. IVV.
### D Thin, permeable membrane
The previous subsection found that a highly-charged, thick, membrane can partition into a zone of strong adhesion and a second zone of weaker adhesion. In this subsection we will introduce another element of realism by accounting for the interior charges in the vesicle. To highlight the key role of the membrane as a barrier to counterions, we will first study the simpler case of a thin membrane permeable to ions, finding uniform attraction. This sets the stage for the more interesting case of an impermeable membrane in Sect. III E below.
Thus the new feature introduced in this subsection is that a membrane separates the world into two compartments, with electrolyte solution on each side (Fig. 3a). We continue to neglect the internal structure of the membrane, treating it as a single thin sheet of charge; in Sect. IV C below we will improve the analysis to include the bilayer structure and finite internal capacitance of real membranes.
To organize the calculation we first note that once again there is only one independent conserved quantity exchanged laterally between zones on the membrane, namely $`\sigma _+`$. Let
$$\sigma _{\mathrm{in}}_{\mathrm{}}^0dze(n_+(z)n_{}(z)).$$
(16)
be the areal density of mobile interior counterions. Note that unlike $`\sigma _+`$, which must be positive and less than $`\sigma _{\mathrm{max}}`$, the interior density $`\sigma _{\mathrm{in}}`$ can in principle have any sign and magnitude. We will hold $`\sigma _{\mathrm{in}}`$ fixed while optimizing over the gap spacing $`\mathrm{}`$ as in Sect. III B above. Since in this subsection we are assuming a permeable membrane, we then minimize over $`\sigma _{\mathrm{in}}`$ as well to obtain $`f(\sigma _+)`$. We will suppress explicit mention of the dependence on the dielectric charge $`\sigma _{}`$, because $`\sigma _{}`$ is fixed.
The free energy density $`f`$ can be regarded as an interior term from (8), plus a gap term, $`f_{\mathrm{gap}}`$ from (11), plus the membrane free energy $`f_\mathrm{m}`$ from (13). The interior term is $`f_{\mathrm{self}}(\sigma _{\mathrm{in}})=(\sigma _{\mathrm{in}})^2/2ϵ\kappa `$. The gap sees opposing charge densities of $`\sigma _{}`$ and $`(\sigma _++\sigma _{\mathrm{in}})`$, so (11) gives $`f_{\mathrm{gap}}=|\sigma _{}^{}{}_{}{}^{2}(\sigma _++\sigma _{\mathrm{in}})^2|/2ϵ\kappa `$ and the equilibrium spacing is $`\mathrm{}_{}(\sigma _{\mathrm{in}},\sigma _+)=\frac{1}{\kappa }\mathrm{ln}\frac{|\sigma _{}|}{\sigma _++\sigma _{\mathrm{in}}}`$.
The membrane free energy $`f_\mathrm{m}`$ is independent of $`\sigma _{\mathrm{in}}`$, so minimizing over $`\sigma _{\mathrm{in}}`$ gives (Fig. 3b) $`\sigma _{\mathrm{in},}=|\sigma _{}|\sigma _+`$, a positive value corresponding to $`\mathrm{}_{}=0`$: the membrane comes into tight contact with the dielectric. Evaluating the free energy at this point gives
$$f(\sigma _+)=\frac{\sigma _{\mathrm{max}}^{}{}_{}{}^{2}}{2ϵ\kappa }\left[(|\overline{\sigma }_{}|\overline{\sigma }_+)^2+\beta \left(\overline{\sigma }_+\mathrm{ln}\overline{\sigma }_++\left(1\overline{\sigma }_+\right)\mathrm{ln}\left(1\overline{\sigma }_+\right)\right)\right].$$
(17)
Computing the second derivative we see that this time every term of $`f`$ is separately convex: there is no instability. Our result is physically reasonable. As the positive membrane approaches the dielectric, the latter’s negative (minority) counterions and some of the positive (majority) counterions get released to the exterior. Since we have assumed the membrane is permeable, the remaining positive counterions pass through it, where many more pair up with interior negative ions from the membrane and get released to the interior. Since no counterions need to remain in the gap, we get tight contact between membrane and dielectric, which in effect become a single object of reduced charge density $`\sigma _+|\sigma _{}|`$. The electrostatic self-energy of this composite object is a convex function, entropy never favors phase separation, and so there is no instability.
### E Thin, impermeable membrane
Previous subsections have shown that charge mobility alone can lead to an instability, but not to charge reversal (Sect. III C), and that introducing a coupled interior compartment alone does not even lead to instability (Sect. III D). Surprisingly, in this section and Sect. IV B below we will find that combining these two unpromising ingredients with the hypothesis of a membrane impermeable to ions can lead to a charge-reversing instability.We need not assume the membrane to be impermeable to water; because the bulk salt concentration is assumed the same on both sides, there will be no net osmotic flow. The key observation is that an impermeable membrane has two conserved quantities independently exchanged between zones: the membrane charge $`\overline{\sigma }_+`$ and the interior counterion charge density $`\overline{\sigma }_{\mathrm{in}}`$.<sup>\**</sup><sup>\**</sup>\**A third exchanged quantity, the net counterion charge density outside the membrane, is then fixed by charge neutrality: $`\sigma _0=(\sigma _1+\sigma _++\sigma _{})`$. Similarly the density of neutral surfactants in the membrane is not independent, being given by $`(2/a_0)(\sigma _+/e)`$. The numbers of individual counterions of each species are not, however, conserved, since neutral $`\pm `$ pairs can be exchanged with large reservoirs (the bulk solution inside and outside the vesicle) without macroscopic charge separation. In the previous subsection $`\overline{\sigma }_{\mathrm{in}}`$ could relax within a zone by passing through the membrane, and so we simply optimized $`f`$ over it before applying the Maxwell construction. For an impermeable membrane we must instead apply the Maxwell construction to both $`\overline{\sigma }_+`$ and $`\overline{\sigma }_{\mathrm{in}}`$ jointly.
The geometry is the same as Fig. 3a. It will shorten some formulæ to define the total charge density
$$\overline{\sigma }_\mathrm{t}\overline{\sigma }_++\overline{\sigma }_{\mathrm{in}}.$$
(18)
The free energy density is the same as in Sect. III D, but this time we need a more general formulation than (17), since we are not simply evaluating at the optimal value of $`\overline{\sigma }_{\mathrm{in}}`$. In fact, there are three physically separate cases we must distinguish: i) The membrane plus its trapped interior counterions may have greater charge density than the dielectric: $`\overline{\sigma }_\mathrm{t}>|\overline{\sigma }_{}|`$. ii) The membrane plus its trapped counterions may have lower charge density than the dielectric, but still be positive: $`0<\overline{\sigma }_\mathrm{t}<|\overline{\sigma }_{}|`$. iii) The trapped counterions may overwhelm and effectively reverse the charge of the membrane: $`\overline{\sigma }_\mathrm{t}<0`$. This is the charge-reversal we seek. In this case the equilibrium distance between the membrane and a negative dielectric is infinity; the membrane actually repels incoming negative objects.
The total free energy density in each of these cases (and still in the linearized approximation) now reads
$`f(\overline{\sigma }_+,\overline{\sigma }_\mathrm{t})`$ $`={\displaystyle \frac{\sigma _{\mathrm{max}}^{}{}_{}{}^{2}}{2\kappa ϵ}}[(\overline{\sigma }_\mathrm{t}\overline{\sigma }_+)^2+\left\{\begin{array}{cc}|(\overline{\sigma }_\mathrm{t})^2(\overline{\sigma }_{})^2|\mathrm{if}\hfill & \overline{\sigma }_\mathrm{t}>0\hfill \\ (\overline{\sigma }_\mathrm{t})^2+(\overline{\sigma }_{})^2\mathrm{if}\hfill & \overline{\sigma }_\mathrm{t}<0\hfill \end{array}\right\}`$ (21)
$`+`$ $`\beta (\overline{\sigma }_+\mathrm{ln}\overline{\sigma }_++(1\overline{\sigma }_+)\mathrm{ln}(1\overline{\sigma }_+))].(\mathrm{linearized}\mathrm{approximation})`$ (22)
The function $`f`$ defines a surface over the $`(\overline{\sigma }_+,\overline{\sigma }_\mathrm{t})`$-plane. If this surface is everywhere convex-down then there is no instability. If not, then it may be possible to bring a straight line up to the surface from below, touching it at two points of tangency but lower than the surface at some point $`(\overline{\sigma }_{+,\mathrm{av}},\overline{\sigma }_{\mathrm{t},\mathrm{av}})`$ lying between those tangency points. In this case a homogeneous system with average composition $`(\overline{\sigma }_{+,\mathrm{av}},\overline{\sigma }_{\mathrm{t},\mathrm{av}})`$ will be able to reduce its free energy by partitioning into zones whose compositions are given by the two points of tangency.
We will assume that initially, when the membrane vesicle was formed and no dielectric spheres were present, the ions were in equilibrium across the membrane, so that half of them got trapped inside: $`\overline{\sigma }_{\mathrm{t},\mathrm{av}}=\frac{1}{2}\overline{\sigma }_{+,\mathrm{av}}`$. We will also for illustration generally take the membrane composition to be half charged and half neutral surfactants, so that the mole fraction $`\overline{\sigma }_{+,\mathrm{av}}=\frac{1}{2}`$. Finally we assume the approaching dielectric to have greater charge density than the maximum possible value for the membrane. For illustration we take $`\overline{\sigma }_{}=3/2`$. Summarizing, we will consider the illustrative case
$`\overline{\sigma }_{+,\mathrm{av}}=\frac{1}{2},\overline{\sigma }_{\mathrm{t},\mathrm{av}}=\frac{1}{2}\overline{\sigma }_{+,\mathrm{av}},\overline{\sigma }_{}=3/2.(\mathrm{illustrative}\mathrm{case})`$
In general there may be many lines in the $`(\overline{\sigma }_+,\overline{\sigma }_\mathrm{t})`$-plane, all passing through the point $`(\overline{\sigma }_{+,\mathrm{av}},\frac{1}{2}\overline{\sigma }_{+,\mathrm{av}})`$ and all exhibiting the instability. In this case we must examine all the lines and choose the one which gives the absolute minimum in free energy. To do this systematically, we label the lines by the real number $`p`$ and write each parametrically as
$$(\overline{\sigma }_+,\overline{\sigma }_\mathrm{t})=(s,\frac{1p}{2}\overline{\sigma }_{+,\mathrm{av}}+\frac{p}{2}s),0<s<1.$$
(23)
In principle one could now plot $`f`$ from (22) along the family of lines defined by (23), find the points of tangency, optimize over $`p`$, and finally obtain the sought instability and the charge densities $`\overline{\sigma }_\mathrm{t}^{(a)}`$ and $`\overline{\sigma }_\mathrm{t}^{(n)}`$ in the two zones as points of tangency, as described in Sect. III C above. If one of these (conventionally $`\overline{\sigma }_\mathrm{t}^{(n)}`$) proves to be negative, then we conclude that the membrane exhibits a charge-reversal instability, as was to be shown. In fact these steps are now rather easy to complete. But we have already remarked that the linearized theory is not accurate in the regime of high charge densities of interest to us. Accordingly we will now improve our theory by solving the full nonlinear Poisson-Boltzmann equation, then carry out the steps just described.
## IV Nonlinear Poisson–Boltzmann Theory
### A Basic formulæ
We introduce the useful new variable
$$\zeta \mathrm{e}^{\kappa z}.$$
(24)
It will also be convenient to define another nondimensional form of the charge density by
$$\stackrel{~}{\sigma }\sigma \kappa /2\widehat{n}e=2\overline{\sigma }/\beta .$$
(25)
and a nondimensional form of the free energy density by
$$\overline{f}=\frac{\kappa }{\widehat{n}k_\mathrm{B}T}f.$$
(26)
The general solution to the Poisson-Boltzmann equation can be written in terms of elliptic functions (see for example ). Fortunately, however, we need only the zero-pressure solutions, corresponding to two walls which are free to adopt their equilibrium spacing $`\mathrm{}_{}`$, and these solutions consist of elementary functions :
$$\overline{\psi }_\pm =\pm 2\mathrm{ln}\frac{\zeta +1}{\zeta 1}.$$
(27)
At $`\zeta \mathrm{}`$ we then have $`\overline{\psi }_\pm \pm 4/\zeta =\pm 4\mathrm{e}^{\kappa z}`$, precisely the weak-field solution we found in Sect. III A.
The Poisson-Boltzmann equation is second-order, and so its general solution has two integration constants. We have fixed one of these by restricting to the zero-pressure case. The other one enters (27) rather trivially, due to the translation invariance of the PB equation: in (27) we are free to shift $`z`$, or equivalently multiply $`\zeta `$ by an arbitrary constant, thus obtaining a one-parameter family of zero-pressure solutions. In practice we will use (27) in the unshifted forms given above, but select the region $`\zeta _1<\zeta <\zeta _2`$ to enforce Gauss’s law at each of the two charged surfaces.
For example, for an isolated surface of charge density $`\sigma _+>0`$ we choose the solution $`\overline{\psi }_+`$ with one limit at infinity and the other at $`\zeta _+`$, which we choose by requiring
$`{\displaystyle \frac{\sigma _+}{ϵ}}={\displaystyle \frac{k_\mathrm{B}T}{e}}{\displaystyle \frac{\mathrm{d}}{\mathrm{d}z}}\overline{\psi }_+|_{z_+}.`$
or
$$\zeta _+=\frac{2}{\stackrel{~}{\sigma }_+}\left(1+\sqrt{1+\stackrel{~}{\sigma }_+^2/4}\right).$$
(28)
The free energy formula analogous to (8) is then obtained by substituting (27) into (5), to get
$$\overline{f}_{\mathrm{self}}=\frac{8\left(2\zeta _+\mathrm{ln}(\frac{1+\zeta _+}{1+\zeta _+})\right)}{1+\zeta _{+}^{}{}_{}{}^{2}}+2\stackrel{~}{\sigma }_+\mathrm{ln}\frac{\zeta _++1}{\zeta _+1}(\mathrm{nonlinear}\mathrm{theory}).$$
(29)
Two surfaces of the same sign charge will repel to infinite separation, so we use this formula for each one separately. For a negatively-charged surface we simply replace $`\stackrel{~}{\sigma }_+`$ by $`|\stackrel{~}{\sigma }_{}|`$ in (28).
It is instructive to compare (29) to the corresponding formula in the linearized approximation, formula (8) (Fig. 4). While the two formulæ agree at low charge density, the linearized formula overestimates the free energy by almost an order of magnitude at the high charge densities of interest to us. The nonlinear PB equation also predicts a narrower cloud of counterions than the linearized approximation at any given charge density.
For two oppositely charged surfaces at their equilibrium spacing we can generalize the argument given in Sect. III B above, again obtaining (11). Again suppose first that $`\sigma _+<|\sigma _{}|`$, and so $`\zeta _+>\zeta _{}`$. By the same logic as in Sect. III B, the potential $`\overline{\psi }`$ in the gap is just the same as that of an isolated surface of charge density $`\sigma _{}`$, but truncated at some finite $`\zeta _+`$.
Let us abbreviate the local integrand in (5) by $`\mathrm{\Phi }=\overline{\psi }\mathrm{sinh}\overline{\psi }2\mathrm{cosh}\overline{\psi }+2`$. This is the same for either of the two solutions $`\overline{\psi }_\pm `$. We thus get
$`\overline{f}_{\mathrm{self}}(\sigma _\pm )=\pm {\displaystyle \frac{\kappa \sigma _\pm }{2e\widehat{n}}}\overline{\psi }_+(\zeta _\pm )+{\displaystyle _{\zeta _\pm }^{\mathrm{}}}d\zeta \mathrm{\Phi }`$
and
$`\overline{f}_{\mathrm{gap}}(\sigma _+,\sigma _{})={\displaystyle \frac{\kappa }{2e\widehat{n}}}\left(\sigma _{}\overline{\psi }(\zeta _{})+\sigma _+\overline{\psi }(\zeta _+)\right)+{\displaystyle _\zeta _{}^{\zeta _+}}d\zeta \mathrm{\Phi }.`$
But the last expression just equals $`\overline{f}_{\mathrm{self}}(\sigma _{})\overline{f}_{\mathrm{self}}(\sigma _+)`$. Repeating for the opposite case $`\sigma _+>|\sigma _{}|`$, we get the desired combination formula (11).
### B Thin, impermeable membrane
Proceeding now as in Sect. III E, we combine (25), (26), (28), (29), (11), (13), and (15) to obtain the analog of the linearized formula (22): the nondimensional free energy density of the membrane+dielectric system at its equilibrium spacing $`\mathrm{}_{}`$, as a function of the local membrane charge density $`\sigma _+`$ and the net charge density $`\sigma _\mathrm{t}`$ of counterions trapped inside the membrane vesicle, is now
$`\overline{f}(\overline{\sigma }_+,\overline{\sigma }_\mathrm{t})`$ $`=`$ $`\overline{f}_{\mathrm{self}}(\overline{\sigma }_\mathrm{t}\overline{\sigma }_+)+\left\{\begin{array}{cc}|\overline{f}_{\mathrm{self}}(\overline{\sigma }_\mathrm{t})\overline{f}_{\mathrm{self}}(\overline{\sigma }_{})|\mathrm{if}\hfill & \overline{\sigma }_\mathrm{t}>0\hfill \\ \overline{f}_{\mathrm{self}}(\overline{\sigma }_\mathrm{t})+\overline{f}_{\mathrm{self}}(\overline{\sigma }_{})\mathrm{if}\hfill & \overline{\sigma }_\mathrm{t}<0\hfill \end{array}\right\}`$ (32)
$`+`$ $`{\displaystyle \frac{4}{\beta }}\left(\overline{\sigma }_+\mathrm{ln}\overline{\sigma }_++(1\overline{\sigma }_+)\mathrm{ln}(1\overline{\sigma }_+)\right).`$ (33)
Here $`\overline{\sigma }_\mathrm{t}=\overline{\sigma }_++\overline{\sigma }_{\mathrm{in}}`$ as before and $`\overline{f}_{\mathrm{self}}(\overline{\sigma })`$ is the function defined by (28), (29) evaluated at $`\stackrel{~}{\sigma }=2\overline{\sigma }/\beta `$. Thus at low charge densities the first two terms of (33) contain factors of $`1/\beta ^2`$, and so dominate the last (mixing-entropy) term, just as in (22).
We can now carry out the program outlined in Sect. III E for the illustrative parameter values $`\widehat{n}=1`$mM, $`a_0=0.5`$nm<sup>2</sup>, $`\beta =0.006`$, $`\overline{\sigma }_{}=3/2`$, and $`\overline{\sigma }_{+,\mathrm{av}}=1/2`$ discussed earlier. Fig. 5 shows the surface defined by (33). For clarity we have shown $`\overline{f}`$ instead of $`\overline{f}`$, so that thermodynamic stability would correspond to an inverted bowl shape. We have also tilted the graph by adding a convenient linear function to $`\overline{f}`$, to highlight the saddle shape. The linear function was selected by trial and error. Adding it does not change the points of tangency between the surface and a straight line.
The graph clearly displays the instability we were seeking. Moreover, one of the two hills on the surface clearly lies to the left of the line of charge reversal, $`\{\overline{\sigma }_\mathrm{t}=0`$}. To make this qualitative observation precise, we must now evaluate (33) along the family of lines specified by (23), perform the Maxwell construction on each line, and choose the value of $`p`$ whose tangent line has the lowest value of $`\overline{f}`$ at the point $`\overline{\sigma }_{+,\mathrm{av}}`$. Fig. 6 shows the result of this analysis for the illustrative values $`p=2.4`$ and 3.8, and the optimal value $`p=2.9`$.
The figure shows coexistence between a zone with $`\overline{\sigma }_+^{(a)}=0.95`$, and another zone with $`\overline{\sigma }_+^{(n)}=0.25`$. The latter zone thus presents total charge density $`\overline{\sigma }_\mathrm{t}=0.11`$ to the outside of the vesicle. Since this is negative, this zone is charge-reversed and deserves its name as a “nonadhesive” zone. Indeed the effect is large: $`\overline{\sigma }_\mathrm{t}`$ is $`45`$% as great as the charge $`\overline{\sigma }_{+,\mathrm{av}}/2=1/4`$ presented to the outside world when there are no adhering dielectrics spheres. Recalling that $`\gamma \overline{\sigma }_+^{(a)}+(1\gamma )\overline{\sigma }_+^{(n)}=\overline{\sigma }_{+,\mathrm{av}}`$, we find that the adhesion zone covers 36% of the vesicle. These results were announced in .
### C Finite thickness, bilayer membrane
While the above results are encouraging, and show the mathematical possibility of a charge-reversal instability, our model needs considerable refinement before we can take its results seriously. In this subsection we begin this task by acknowledging the bilayer character of the membrane and its finite capacitance, both neglected up to this point. The results in this section were also announced in .
Instead of idealizing the membrane as a thin sheet of charge density $`\sigma _+`$, we now regard it as two sheets of charge density $`u\sigma _+`$ and $`(1u)\sigma _+`$ representing the charged headgroups of the inner and outer surfactant layers respectively (see Fig. 1d). These two layers of charge are separated by a dielectric layer of thickness $`t`$ and dielectric constant $`ϵ_\mathrm{m}`$, creating a capacitor of capacitance $`c=ϵ_\mathrm{m}/t`$ per area. We will estimate $`c`$ using the value 0.01 pF/$`\mu \mathrm{m}^2`$ typical for artificial bilayer membranes and make the useful abbreviation
$$\tau t\kappa ϵ/ϵ_\mathrm{m}=\kappa ϵ/c.$$
(34)
Then $`\tau 7`$ at salt concentration $`\widehat{n}=1`$mM, or more generally $`\tau 7\sqrt{\widehat{n}/\mathrm{mM}}`$.
The free energy formula (33) used in Sect. IV B needs only two simple modifications: 1) Since the membrane still presents charge density $`\sigma _+\sigma _\mathrm{t}`$ to the interior solution and $`\sigma _\mathrm{t}`$ to the exterior (Fig. 1), the terms involving $`\overline{f}_{\mathrm{self}}`$ are unchanged. Now, however, when $`\sigma _{\mathrm{in}}+u\sigma _+`$ is nonzero (or equivalently $`\sigma _\mathrm{t}+(u1)\sigma _+0`$), there will be a nonzero electric field in the membrane’s interior, with a capacitive energy cost per area of $`\left(\sigma _\mathrm{t}+(u1)\sigma _+\right)^2/2c`$. 2) The membrane now consists of two fluid monolayers of mixed charged and neutral surfactants. Each monolayer has a maximal density $`\frac{1}{2}\sigma _{\mathrm{max}}=e/a_0`$, attained when the density of neutrals is zero. Accordingly we replace the mixing entropy $`f_\mathrm{m}`$ from (13) by
$`{\displaystyle \frac{k_\mathrm{B}T}{a_0}}\left[{\displaystyle \frac{u\sigma _+}{\sigma _{\mathrm{max}}/2}}\mathrm{ln}{\displaystyle \frac{u\sigma _+}{\sigma _{\mathrm{max}}/2}}+\mathrm{}\right].`$
Casting everything into the nondimensional forms defined above, we see that we must add to the formula (33) for $`\overline{f}`$ the expression
$`{\displaystyle \frac{4\tau }{\beta ^2}}\left(\overline{\sigma }_\mathrm{t}+(u1)\overline{\sigma }_+\right)^2`$ $`+{\displaystyle \frac{2}{\beta }}[2u\overline{\sigma }_+\mathrm{ln}(2u\overline{\sigma }_+)+(12u\overline{\sigma }_+)\mathrm{ln}(12u\overline{\sigma }_+)`$ (35)
$`+2(1u)`$ $`\overline{\sigma }_+\mathrm{ln}2(1u)\overline{\sigma }_++(12(1u)\overline{\sigma }_+)\mathrm{ln}(12(1u)\overline{\sigma }_+)].`$ (36)
To use our formulæ we add (36) to (33) and again hold fixed the two conserved quantities $`\overline{\sigma }_+`$ and $`\overline{\sigma }_\mathrm{t}`$. As before we must optimize over all other variables, in this case just $`u`$, before performing the Maxwell construction as in Sect. IV B. We can simply optimize (36) over $`u`$, since $`u`$ does not enter (33). However, this optimization is subject to the four inequalities which $`u`$ must obey:
$`0<u\overline{\sigma }_+<1/2,0<(1u)\overline{\sigma }_+<1/2.`$
A graphical analysis similar to the one shown in Fig. 6 now gives (Fig. 7) that the best line through $`(\overline{\sigma }_+,\overline{\sigma }_\mathrm{t})=(1/2,1/4)`$ has $`p2`$ (see (23)), with points of tangency at $`\overline{\sigma }_+^{(a)}=0.247`$, and $`\overline{\sigma }_+^{(n)}=0.65`$. Proceeding as in Sect. IV B gives coverage $`\gamma _{}=63`$% at equilibrium and $`\overline{\sigma }_\mathrm{t}=0.003`$, or about $`1.2`$% of the charge density presented to the outside world when no dielectric spheres are present.
We can readily understand the qualitative features of these results. Since $`\tau `$ is large, electric fields inside the membrane are energetically costly and the two sides of the membrane are nearly independent. Thus the interior ion charge density $`\sigma _{\mathrm{in}}`$ remains nearly uniform, and hence nearly equal to $`\sigma _{+,\mathrm{av}}/2`$, and similarly the inner monolayer charge density $`u\sigma _+\sigma _{+,\mathrm{av}}/2`$. Then the total charge density $`\overline{\sigma }_\mathrm{t}(1u)\overline{\sigma }_+\frac{1u}{2u}\overline{\sigma }_{+,\mathrm{av}}`$ (see Fig. 1d), and (23) requires that either $`u=1/2`$ or $`p2`$. The solution $`u=1/2`$ is unphysical; the solution $`p2`$ is just what we found numerically. To reverse its charge, the membrane must allow electric fields in its interior (see Fig. 1d); the high cost of doing this accounts for the very sharp left-hand dip in Fig. 7 compared to Fig. 6, and for the greatly reduced degree of charge reversal in the finite-thickness case. The charge-reversal effect diminishes for larger values of $`\tau `$, as we predicted in Sect. II. According to (34) this means the effect will disappear either for thick membranes or at large enough ion strength $`\widehat{n}`$. Numerically we find the critical value to be about $`20`$mM, roughly as seen in the experiments of .
Though the charge-reversal effect seems small, it is enough to cause the rejection of additional negative dielectric spheres. To estimate the magnitude of this effect, consider what is needed to increase $`\gamma `$ from its equilibrium value $`\gamma _{}`$ to $`\gamma _{}+\delta `$. To do this we must choose new values of $`\overline{\sigma }_+^{(a)}+ϵ^{(a)}`$ and $`\overline{\sigma }_+^{(n)}+ϵ^{(n)}`$, subject to the condition (14), which now reads
$`(\gamma _{}+\delta )(\overline{\sigma }_+^{(a)}+ϵ^{(a)})+(1\gamma _{}\delta )(\overline{\sigma }_+^{(n)}+ϵ^{(n)})=\overline{\sigma }_{+,\mathrm{av}}.`$
We then minimize the total free energy $`F`$ over $`ϵ^{(a)}`$ and $`ϵ^{(n)}`$ subject to this constraint, finding that the increase in $`F`$ when we force a nonequilibrium value of $`\gamma `$ is
$`\mathrm{\Delta }F={\displaystyle \frac{k_\mathrm{B}T\widehat{n}}{\kappa }}{\displaystyle \frac{1}{2}}(\overline{\sigma }_+^{(a)}\overline{\sigma }_+^{(n)})^2\left[{\displaystyle \frac{1\gamma _{}}{\overline{f}_{(n)}^{\prime \prime }}}+{\displaystyle \frac{\gamma _{}}{\overline{f}_{(a)}^{\prime \prime }}}\right]^1A\delta ^2.`$
In this expression $`\overline{f}_{(a)}^{\prime \prime }`$ denotes $`\frac{\mathrm{d}^2\overline{f}}{\mathrm{d}\overline{\sigma }_{+}^{}{}_{}{}^{2}}|_{\overline{\sigma }_+^{(a)}}`$, etc., and $`A`$ is the total membrane area. Bringing an additional 1$`\mu `$m<sup>2</sup> of negative dielectric into contact with a vesicle of area $`4\pi (10\mu \mathrm{m})^2`$ gives $`\delta =0.00083`$. Evaluating numerically then gives $`\mathrm{\Delta }F3000k_\mathrm{B}T`$, a huge barrier to adhesion, and similarly if we pull away 1$`\mu `$m<sup>2</sup> of adhering dielectric.
## V Effects of Ion Correlations
The charge density near a highly-charged surface can become so great as to invalidate the mean-field theory we have used so far. The resulting changes in the force between two plates have been the object of intense study since the discovery that the total force can become attractive for two like-charged plates, in the presence of multivalent counterions . The situation we will need to study will be much simpler than that one, for two reasons. First, we are only interested in the free energy density at zero force (i.e., equilibrium). Secondly, our effect has arisen already at the level of mean-field theory. Since we consider the case of monovalent counterions only, i.e. in the regimes of small to moderate ion interactions, correlation effects will turn out to be a modest correction to the main, mean-field, contribution. Thus we are in the regime opposite to that recently studied in refs. .
The criterion for mean-field theory to be valid is roughly that the electrostatic potential energy of two ions at the mean ionic separation be smaller than the thermal energy: $`e^2n_\pm ^{1/3}/4\pi ϵ<k_\mathrm{B}T`$, or equivalently that $`n_\pm \mathrm{}_B^3<1`$. Any isolated surface, no matter how highly charged, will have a distance beyond which this criterion is satisfied, and so our universal Poisson-Boltzmann solution (27) will be valid there. We will call the boundary of this region $`z=z_s`$. We find $`z_s`$ using (27) and (4), obtaining $`z_s=0.28`$nm for our illustrative case of ambient salt concentration $`\widehat{n}=1`$mM.
Thus any isolated planar surface has exactly the same potential (equation (27)) as any other, for $`z>z_s`$. The charge density $`\sigma _+`$ enters only via the location $`z_+`$ of the surface in the coordinate $`z`$. Given a surface of charge density $`\sigma _+`$, we compute $`z_+=\kappa ^1\mathrm{ln}\zeta _+`$, where $`\zeta _+`$ is given by (28). If $`z_+>z_s=0.28`$nm, then mean-field theory is everywhere accurate and there is no correlated-ion cloud near the surface. In the opposite case, that part of the ion cloud lying within the layer $`z_+<z<z_s`$ will have nonnegligible correlations. Since $`z_+`$ is always positive, this layer is never any thicker than a typical ion radius, and so may be treated as a two-dimensional classical charged gas . This approach may be regarded as an approximation to other, more refined, calculations (e.g. refs. ).
The effect of correlations will be to reduce the free energy density, as ions can avoid each other, reducing their electrostatic self-energy. To apply the results of Totsuji, originally derived for use in the study of electrons adsorbed onto liquid helium , we need to know the two-dimensional density $`m`$ of counterions in the correlated layer. $`m`$ simply equals $`\sigma _+/e`$ minus the total density in the uncorrelated region $`z>z_s`$. Again using (27) and (4) in the latter region, we find $`m=(\sigma _+/e)0.81\mathrm{nm}^2`$. If this quantity is negative then there simply is no correlated layer and $`m=0`$. Defining the plasma parameter as $`\mathrm{\Gamma }=\mathrm{}_B\sqrt{\pi m}`$, the correlation energy density can then be represented by the interpolation formula $`E_c=mk_\mathrm{B}T\mathrm{\Gamma }(1.07+\frac{1}{2.2\mathrm{\Gamma }+1.3})`$, which is approximately valid over the range $`0<\mathrm{\Gamma }<5000`$ . Fig. 8 shows the resulting change in the free energy density, obtained from the thermodynamic formula $`f_c(m)=mk_\mathrm{B}T_0^\mathrm{\Gamma }\frac{\mathrm{d}\mathrm{\Gamma }}{\mathrm{\Gamma }}\frac{E_c}{mk_\mathrm{B}T}`$ for the correlation contribution $`f_c`$ to the free energy density.
When two oppositely-charged surfaces face each other, we have seen how the minority surface will be stripped of its counterions; the majority surface may, however, have a correlated ion cloud, so we add $`f_c(m)`$ to its free energy density. Since in the situations of interest to us the majority surface is always the dielectric sphere, and since $`m`$ depends only on the surface’s charge (not on the presence of the other surface), we find that the correlation correction to the free energy of the negatively-charged dielectric surface is a constant, and may simply be dropped. We do need to include the correlation free energy of the membrane’s interior surface, but for a finite-thickness membrane this too is nearly a constant, since as we have seen the inner monolayer’s charge density $`u\sigma _+`$ deviates only slightly from $`\sigma _{+,\mathrm{av}}/2`$. Finally, in conditions of charge reversal the outer monolayer becomes isolated and can have an ion cloud of its own. Since as we have seen the degree of charge reversal is very small, the density $`m`$ of this last ion cloud is very small and the correlation contribution is negligible.
We have just outlined qualitatively why counterion correlations may be expected to have little effect on the results given in Sect. IV C. Indeed, the numerically calculated graph analogous to Fig. 7 is not appreciably different from that graph, and we do not display it here.
## VI Discussion
We have proposed a theoretical explanation for the phenomenon of electrostatic adhesion saturation observed experimentally in . While the experimental system has not been systematically explored yet, our model reproduces qualitatively the surprising the phenomenon of charge reversal and several salient experimental facts : 1) Adhesion saturation occurs only with mixed bilayer vesicles, that is, at mole ratios $`\overline{\sigma }_{+,\mathrm{av}}`$ not too close to zero or unity. 2) It occurs only under conditions of sufficiently low salt. 3) The saturated state has a very definite number of adhering objects ($`\gamma _{}`$ is fixed for each vesicle).
Our analysis has omitted many familiar colloidal-force effects. Many of these are short-ranged (e.g. solvation forces), weak compared to electrostatic forces (e.g. undulation repulsion), or rapidly decreasing with distance (e.g. van der Waals forces). In addition we have neglected all finite ion-size effects. We believe that our conclusions will be robust when such effects are introduced, in part because the crucial physics of charge reversal involves the immediate neighborhood of the left-hand dip in Fig. 7, namely, the separation between the charge-reversal point and the tangency point. But the distance $`\mathrm{}_{}`$ between the membrane and dielectric diverges as we approach the charge-reversal point from the right, so this physics is controlled by the long-distance behavior of the forces. Certainly the exact location of the tangent point depends on the right-hand part of Fig. 7 as well, where our theory is not reliable. But this dependence is small due to the sharpness of the left-hand dip in the free energy density. Even if the right-hand side of the graph differs from what we computed, there should be a range of membrane compositions $`\overline{\sigma }_{+,\mathrm{av}}`$ greater than $`\overline{\sigma }_+^{(n)}`$ but low enough to be in the left part of the graph, and hence yielding the sort of zone separation we have studied.
We have examined only equilibrium states. It is quite possible that the experimental system of is not in equilibrium, i.e. that the observed coverage $`\gamma `$ is less than the equilibrium value $`\gamma _{}`$ because the last one or two balls is initially repelled by a finite free-energy barrier. But our goal was to understand the surprising existence of any barrier, not to predict a specific value for $`\gamma _{}`$, which in any case depends on the membrane composition.<sup>††</sup><sup>††</sup>††Moreover, the observed ball coverage does not directly give $`\gamma `$, since the degree of each ball’s coverage is not optically observable; see .
The analysis suggests a number of experimental tests of our mechanism. A mixed vesicle adhering to a charged dielectric surface may provide a more controlled geometry than that of ; in this case adhesion saturation suggests the possibility of observing an adhering, yet flaccid, vesicle. A more ambitious test could be arranged by washing out the exterior solution, replacing it by another of different ion strength but the same osmolarity, while pinning a single vesicle for observation with a micropipette. Our formulæ generalize readily to the case where the ionic strength $`\widehat{n}`$ is different inside and outside the vesicle. In this way may be possible reversibly to turn adhesion saturation on and off.
###### Acknowledgements.
We thank R. Bruinsma, I. Rouzina, and B. Shklovskii for discussions, and especially H. Aranda-Espinoza, N. Dan, T. C. Lubensky, L. Ramos, and D. A. Weitz for an earlier collaboration leading to the ideas presented here. This work was supported in part by NSF grant DMR98-07156.
## Appendix: Notation
### Constants
We work in SI units. Thus the potential around a point charge $`q`$ in vacuum is $`\psi (r)=q/4\pi ϵ_0r`$, where $`ϵ_0=910^{12}`$Farad/meter. We treat water as a continuum dielectric with $`ϵ=80ϵ_0`$; inside the membrane $`ϵ_m2ϵ_0`$. The Bjerrum length in water is $`\mathrm{}_B=e^2/4\pi ϵk_\mathrm{B}T`$; thus $`4\pi \mathrm{}_B=8.7`$nm.
### Parameters
We take for illustration a typical ambient salt concentration of $`\widehat{n}=1`$mM=$`610^4`$nm<sup>-3</sup>. Then the inverse Debye length is $`\kappa =\sqrt{2\widehat{n}e^2/ϵk_\mathrm{B}T}=\sqrt{\widehat{n}/\mathrm{mM}}/(9.8`$nm). The salt concentration inside the vesicle is the same, due to osmotic clamping.
We suppose a mixture of surfactants, which for simplicity have equal area per headgroup $`a_0=0.5`$nm<sup>2</sup>. Then $`\sigma _{\mathrm{max}}=2e/a_0`$ is the maximum bilayer charge density and the parameter $`\beta =2\widehat{n}a_0/\kappa =0.006`$ measures the relative importance of mixing-entropy and electrostatic effects.
We use a typical artificial bilayer capacitance of $`c=0.01`$pF/$`\mu `$m<sup>2</sup>, which enters only in combination with the membrane thickness $`t`$ via $`\tau =t\kappa ϵ/ϵ_\mathrm{m}7\sqrt{\widehat{n}/\mathrm{mM}}`$.
For illustration we take the experimentally-controllable mole fraction of charged surfactants to be $`\overline{\sigma }_{+,\mathrm{av}}=1/2`$ and one half of the corresponding counterions to be trapped on the vesicle interior, so that $`\overline{\sigma }_{\mathrm{t},\mathrm{av}}=\overline{\sigma }_{+,\mathrm{av}}/2=1/4`$. We also take the approaching charged dielectric objects to have charge density 50% greater than the membrane, or $`\overline{\sigma }_{}=3/2`$.
### Variables
We generally denote nondimensionalized quantities with a bar or tilde: thus $`\overline{\sigma }\sigma /\sigma _{\mathrm{max}}`$, while $`\stackrel{~}{\sigma }=\sigma \kappa /2\widehat{n}e=2\overline{\sigma }/\beta `$. Also the free energy density $`f`$ gives rise to $`\overline{f}=\kappa f/\widehat{n}k_\mathrm{B}T=f/(610^3k_\mathrm{B}T/\mu \mathrm{m}^2)\sqrt{\widehat{n}/\mathrm{mM}}`$, while the electrostatic potential $`\psi `$ gives $`\overline{\psi }=e\psi /k_\mathrm{B}T`$. Various contributions to $`f`$ include the mixing entropy of membrane surfactants $`f_\mathrm{m}`$ and the correlation contribution $`f_c`$.
The charge density $`\sigma `$ of a surface determines its Gouy-Chapman length $`\lambda _{\mathrm{GC}}=2ϵk_\mathrm{B}T/e\sigma `$. Various charge densities in the text are defined in Fig. 1, for example $`\sigma _t=\sigma _{\mathrm{in}}+\sigma _+`$. $`m`$ denotes the 2d number density of ions in the dense correlated cloud near a surface. $`m`$ in turn determines the plasma parameter $`\mathrm{\Gamma }\mathrm{}_B\sqrt{\pi m}`$.
Geometrical quantities include the gap width $`\mathrm{}`$, the total membrane area $`A`$, the fraction $`\gamma `$ of $`A`$ in the adhesion zone and its equilibrium value $`\gamma _{}`$. The distance $`z`$ from a surface is sometimes expressed using $`\zeta =\mathrm{e}^{\kappa z}`$. |
warning/0002/math0002065.html | ar5iv | text | # A generalization of Cayley submanifolds
## 1 Introduction
Cayley submanifolds were defined by Harvey-Lawson \[HL1\] and by McLean \[McL\], as calibrated submanifolds of $`\mathrm{Spin}(7)`$-manifolds. Each such manifold $`M`$ admits a parallel calibration $`\mathrm{\Phi }^4(M)`$ whose stabilizer is $`\mathrm{Spin}(7)`$. It is called the *Cayley calibration*, due to the link with the octonians, and the corresponding minimal varieties are called *Cayley submanifolds*.
If the ambient manifold is a Calabi-Yau manifold, the Cayley calibration is not unique. Indeed, given any parallel normalized<sup>1</sup><sup>1</sup>1I.e. such that
$$\frac{\omega ^4}{4!}=\mathrm{vol}_g=\left(\frac{\sqrt{1}}{2}\right)^4\mathrm{\Omega }\overline{\mathrm{\Omega }}.$$
This amounts to requiring that $`\mathrm{\Omega }^{}=1`$, where $`||||^{}`$ denotes the comass norm. complex volume form $`\mathrm{\Omega }^{4,0}(M)`$, the form
$$\mathrm{\Phi }_\mathrm{\Omega }:=\mathrm{Re}\mathrm{\Omega }+\frac{\omega ^2}{2}$$
is a calibration whose stabilizer is isomorphic to $`\mathrm{Spin}(7)`$. Therefore there is an $`\mathrm{S}^1`$-family of Cayley calibrations. A submanifold calibrated by any of these forms, has the property that its Kähler angles coincide. It is therefore natural to consider the submanifolds defined by the latter condition, without any assumption relating to the calibrating forms. This makes sense whether or not the ambient manifold is Ricci-flat, and gives rise to an interesting family of not necessarily minimal submanifolds, including both the Lagrangian and the complex submanifolds as extreme cases.
In this paper we start collecting some facts from linear algebra, making precise the relation between the Kähler angles on one side, and the Cayley calibrations on the other. Then we define the submanifolds with equal Kähler angles, which we call *Cayley*, and prove a formula (38) relating the angle, the mean curvature and the Ricci form of the ambient manifold. Finally we apply this formula to the case where the ambient manifold is Kähler-Einstein, obtaining the following two results:
###### Theorem 1
Let $`(M,J,g)`$ be a Calabi-Yau manifold. Then a Cayley submanifold of $`M`$, is minimal iff it is calibrated by some parallel Cayley calibration.
###### Theorem 2
Let $`(M,J,g)`$ be a Kähler-Einstein manifold of non-zero scalar curvature. Then any (connected) minimal Cayley submanifold of $`M`$ is either complex or (minimal) Lagrangian.
The first result shows the relation with the theory of Harvey and Lawson.
The last result has been obtained, independently and very recently, also by Isabel Salavessa and Giorgio Valli \[SV\], by quite different methods.
Acknowledgments: The author wants to thank Gang Tian for proposing him the subject of this work, and for the constant encouragement. He is also grateful to his advisor, Paolo de Bartolomeis, and to Claudio Arezzo, for interesting discussions.
## 2 Linear algebra of real 4-planes in $`^4`$
Let $`(V,J,g)`$ be a Hermitian vector space of real dimension 8. Denote by $`\omega (X,Y)=g(JX,Y)`$ the associated Kähler form. Given a subspace $`WV`$ we denote by $`\pi _W`$ the orthogonal projection onto $`W`$, and we put
$$B_W:=\pi _WJ_{|W}.$$
$`B_W`$ is a skew-hermitian operator on $`W`$ with respect to $`g`$. We let $`\mathrm{G}(p,V)`$ denote the Grassmannian of *oriented* $`p`$-planes in $`V`$.
Let us recall an important lemma proved by Harvey and Lawson \[HL2\], applied to our situation.
###### Lemma 1 (Canonical form of a 4-plane over $`U(4)`$)
Let $`(V,J,g)`$ be a Hermitian vector space of real dimension $`8`$. Then, given $`\xi \mathrm{G}(4,V)`$, there is a unitary basis $`u_1,u_2,u_3,u_4`$ of $`V`$ and angles $`\theta _1,\theta _2`$, with
$$\begin{array}{c}0\theta _1\frac{\pi }{2}\\ \theta _1\theta _2\pi \end{array}$$
(1)
such that
$$\begin{array}{cc}\hfill \xi =& u_1\left(\mathrm{cos}\theta _1Ju_1+\mathrm{sin}\theta _1u_2\right)\hfill \\ \hfill & u_3\left(\mathrm{cos}\theta _2Ju_3+\mathrm{sin}\theta _2u_4\right).\hfill \end{array}$$
(2)
Therefore
$$B_\xi =\left(\begin{array}{cc}\begin{array}{cc}0& \mathrm{cos}\theta _1\\ \mathrm{cos}\theta _1& 0\end{array}& 0\\ 0& \begin{array}{cc}0& \mathrm{cos}\theta _2\\ \mathrm{cos}\theta _2& 0\end{array}\end{array}\right)$$
(9)
and
$$\omega _{|\xi }=\mathrm{cos}\theta _1e^{12}+\mathrm{cos}\theta _2e^{34}.$$
(10)
The numbers $`\theta _1`$ and $`\theta _2`$ are called the *Kähler angles* of the 4-plane $`\xi `$.
###### Definition 1
$`\xi \mathrm{G}(4,V)`$ is called a Cayley 4-plane if
$$\omega _{|\xi }=_\xi =\omega _{|\xi }.$$
(11)
Here $`_\xi `$ is the Hodge operator of the metric $`g_{|\xi }`$.
###### Lemma 2
An oriented 4-plane $`\xi \mathrm{G}(4,V)`$ is a Cayley subspace if and only if its Kähler angles coincide. In this case, putting $`\mathrm{cos}\theta _1=\mathrm{cos}\theta _2=\lambda [0,1]`$, we have
$$B_\xi ^2=\lambda ^2\mathrm{Id}$$
(12)
$$(\omega ^2)_{|\xi }=2\lambda ^2\mathrm{vol}_{g_{|\xi }}$$
(13)
and there is a a positive orthonormal basis $`e_1,e_2,e_3,e_4`$ of $`\xi `$ such that
$$\omega _{|\xi }=\lambda (e^{12}+e^{34})$$
(14)
$$B_\xi =\left(\begin{array}{cc}\begin{array}{cc}0& \lambda \\ \lambda & 0\end{array}& 0\\ 0& \begin{array}{cc}0& \lambda \\ \lambda & 0\end{array}\end{array}\right)$$
(21)
P r o o f.
Just apply $`_\xi `$ to (10).
Q.D.E.
A positive orthonormal basis $`\{e_1,\mathrm{},e_4\}`$ in which (14) hold is called a *Cayley basis*. If we put
$$𝔛:=\{\xi \mathrm{G}(4,V):\xi \text{ is Cayley}\},$$
then we have a well defined function
$$\lambda :𝔛[0,1].$$
###### Lemma 3
* $`𝔛`$ is a closed subset of $`\mathrm{G}(4,V)`$ and $`\lambda `$ is a continuous function.
* $`\lambda ^1(1)`$ is the Grassmannian of complex planes in $`(V,J)`$, while
$$𝔛_r:=\{\xi 𝔛:\lambda (\xi )<1\}$$
consist of totally real subspaces.
* $`\lambda ^1(0)`$ is the (oriented) Lagrangian grassmannian, while every $`\xi 𝔛`$ with $`\lambda (\xi )>0`$ is a symplectic subspace of $`(V,\omega )`$.
P r o o f.
Let us consider the following subset of the Stiefel manifold of orthonormal quadruples of vectors in $`V`$:
$$\begin{array}{c}𝒴=\{(e_1,e_2,e_3,e_4):\omega (e_1,e_2)=\omega (e_3,e_4)\hfill \\ \hfill \omega (e_1,e_3)=\omega (e_1,e_4)=\omega (e_2,e_3)=\omega (e_2,e_4)=0\}\end{array}$$
(22)
$`𝒴`$ is a closed subset, and the projection $`\pi :𝒴𝔛`$ is onto, therefore it is an identification, i.e. $`𝔛`$ has the quotient topology. As $`\lambda \pi (e_1,e_2,e_3,e_4)=\omega (e_1,e_2)`$, $`\lambda \pi `$ is a continuos function on $`𝒴`$, hence the same is true of $`\lambda `$. The remaining statements are trivial.
Q.D.E.
###### Lemma 4
Let $`\xi `$ be a non-complex, hence totally real Cayley 4-plane. Given any Cayley basis $`\{e_i\}`$ of $`\xi `$, we put
$$\begin{array}{cc}\hfill u_1& =e_1\hfill \\ \hfill u_3& =e_3\hfill \end{array}\begin{array}{cc}\hfill u_2& =\frac{1}{\sqrt{1\lambda ^2}}(e_2\lambda Je_1)\hfill \\ \hfill u_4& =\frac{1}{\sqrt{1\lambda ^2}}(e_4\lambda Je_3).\hfill \end{array}$$
(23)
Then $`\{u_j\}`$ is a unitary basis of $`V`$ and
$$\begin{array}{cc}\hfill \xi =& u_1\left(\lambda (\xi )Ju_1+\sqrt{1\lambda ^2(\xi )}u_2\right)\hfill \\ \hfill & u_3\left(\lambda (\xi )Ju_3+\sqrt{1\lambda ^2(\xi )}u_4\right).\hfill \end{array}$$
(24)
P r o o f.
A straightforward computation shows that
$$g(u_i,u_j)=\delta _{ij}\omega (u_i,u_j)=0.$$
(24) follows immediately from (23).
Q.D.E.
###### Lemma 5
* If $`\xi \mathrm{G}(4,V)`$ is totally real (i.e. if $`\mathrm{cos}\theta _10\mathrm{cos}\theta _2`$) there exists a unique normalized (4,0)-form $`\mathrm{\Omega }_\xi `$ such that
$$\mathrm{\Omega }_\xi (\xi )>0.$$
If we write $`\xi `$ in the form (2), then $`\stackrel{}{u}=u_1u_2u_3u_4`$ satisfies $`\mathrm{\Omega }_\xi (\stackrel{}{u})=1`$. In particular, two basis $`\{u_i\}`$ such that (2) hold differ by an element of $`\mathrm{SU}(4)`$.
* If $`\xi `$ is Cayley and totally real, it is calibrated by the Cayley calibration associated to $`\mathrm{\Omega }_\xi `$:
$$\mathrm{\Phi }_\xi =\mathrm{Re}\mathrm{\Omega }_\xi +\frac{\omega }{2}\mathrm{\Phi }_\xi (\xi )=1.$$
* If $`\xi `$ is calibrated by some Cayley calibration, $`\mathrm{\Phi }_\mathrm{\Omega }(\xi )=1`$, then it is a Cayley subspace, and $`\mathrm{\Omega }_\xi =\mathrm{\Omega }`$.
P r o o f.
From the constraints (1) descends that
$$\begin{array}{c}\mathrm{sin}\theta _1\mathrm{sin}\theta _20\\ \mathrm{cos}\theta _10.\end{array}$$
If $`\xi `$ is totally real, then $`\mathrm{sin}\theta _i0`$, and $`\mathrm{sin}\theta _1\mathrm{sin}\theta _2>0.`$ If we let $`\mathrm{\Omega }_\xi `$ be the unique (4,0)-form such that $`\mathrm{\Omega }_\xi (\stackrel{}{u})=1`$, then $`\mathrm{\Omega }_\xi (\xi )=\mathrm{sin}\theta _1\mathrm{sin}\theta _2>0`$. This shows $`\mathrm{\Omega }_\xi `$ only depends on $`\xi `$ and proves (a). Using the representation (24) we see that
$$\mathrm{\Omega }_\xi (\xi )=1\lambda ^2(\xi )\omega ^2(\xi )=2\lambda ^2(\xi ),$$
thus proving (b). On the other hand, using the representation (2) we see that
$$\mathrm{\Omega }_\xi (\xi )=\mathrm{sin}\theta _1\mathrm{sin}\theta _2\omega ^2(\xi )=2\mathrm{cos}\theta _1\mathrm{cos}\theta _2.$$
Therefore, if $`\mathrm{\Omega }=e^{\sqrt{1}\alpha }\mathrm{\Omega }_\xi `$,
$$\begin{array}{cc}\hfill \mathrm{\Phi }_\mathrm{\Omega }(\xi )& =\mathrm{Re}\left(e^{\sqrt{1}\alpha }\mathrm{\Omega }_\xi (\xi )\right)+\mathrm{cos}\theta _1\mathrm{cos}\theta _2=\hfill \\ & =\mathrm{cos}\alpha \mathrm{sin}\theta _1\mathrm{sin}\theta _2+\mathrm{cos}\theta _1\mathrm{cos}=\theta _2\hfill \end{array}$$
and this can be 1, only if $`\mathrm{cos}\alpha =1`$ and $`\theta _1=\theta _2`$.
Q.D.E.
## 3 Cayley submanifolds of Kähler manifolds
Let $`(M,J,g)`$ be a Kähler manifold of complex dimension 4. We consider an oriented submanifold $`NM`$ of real dimension 4.
We let $`_N`$ denote the Hodge operator of the metric $`g_{|N}`$.
###### Definition 2
We call $`N`$ a Cayley submanifold if the equation
$$\omega _{|N}=_N\omega _{|N}$$
(25)
is satisfied on $`N`$.
This just means that for any point $`x`$ of $`N`$, the oriented tangent space $`T_xN`$ is a Cayley subspace of $`T_xM`$.
We stress that this definition does NOT agree with the one given by Harvey and Lawson, which makes senses on any $`\mathrm{Spin}(7)`$-manifolds and implies that the submanifold is volume-minimizing. The above definition on the contrary makes sense on any Kähler manifold, and does not imply minimality. Just consider that any Lagrangian submanifold has equal (and zero) Kähler angles, and is therefore Cayley, according to the above definition.
The relation between this definition and the one of Harvey and Lawson, in the case where the ambient manifold is Calabi-Yau, is the subject of theorem 1.
As the tangent spaces to $`N`$ are Cayley subspaces, if we denote by $`B_x`$ the endomorphism $`\pi J_x_{|T_xN}`$, then $`B_x^2`$ is a multiple of the identity at each point $`x`$ of $`N`$. We can define a function $`\lambda =\lambda (x)0`$, such that
$$B_x^2=\lambda ^2(x)\mathrm{Id}.$$
As $`\omega _{|N}^2=2\lambda ^2\mathrm{vol}`$, we deduce that $`\lambda ^2`$ is a smooth function on $`N`$, with values in $`[0,1]`$.
Given any 4-dimensional submanifold of $`M`$, not necessarily Cayley, we denote by $`N_r`$ the totally real part of $`N`$, and by $`N_c`$ the set of complex points. If $`N`$ is Cayley, then $`N_r=\{xN:\lambda (x)<1\}`$ and $`N_c=\lambda ^1(1)`$. In particular $`N=N_rN_c`$. Taking the square root of $`\lambda ^2`$ we deduce that $`\lambda `$ is a continuos function on $`N`$, smooth on $`N_r`$, i.e. away from complex points.
On $`N_r`$ is defined a section $`\mathrm{\Omega }_N`$ of $`^{4,0}(M)_{|N}=K_M_{|N}`$, determined by the condition that $`\mathrm{\Omega }_N`$ be normalized and satisfy $`\mathrm{\Omega }_N(T_xN)>0`$ at each point. This is seen applying lemma 5.
###### Lemma 6
Given a Cayley submanifold $`N`$, near each non-complex point $`x`$ of $`N`$, one can find a smooth Cayley frame $`\{e_1,e_2,e_3,e_4\}`$, and a smooth unitary frame $`\{u_1,u_2,u_3,u_4\}`$ in $`TM_{|N}`$ such that
$$\begin{array}{cc}\hfill T_xN=& u_1\left(\lambda (x)Ju_1+\sqrt{1\lambda ^2(x)}u_2\right)\hfill \\ \hfill & u_3\left(\lambda (x)Ju_3+\sqrt{1\lambda ^2(x)}u_4\right).\hfill \end{array}$$
(26)
In particular
$$\mathrm{\Omega }_N(u_1,u_2,u_3,u_4)=1,$$
and therefore $`\mathrm{\Omega }_N:NK_M_{|N}`$ is a smooth section.
P r o o f.
For $`xN_r`$, let us consider the endomorphism
$$j_x=\frac{B_x}{\lambda (x)}$$
of $`T_xN`$. It is a $`g_{|N}`$-orthogonal almost complex structure, compatible with the orientation of $`N`$, smooth on all of $`N_r`$. Therefore we know that near any $`xN_r`$ we can find a smooth $`j`$-unitary frame $`\{e_1,e_3\}`$ in $`TN`$, i.e. a positive orthonormal frame in $`TN`$ of the form $`\{e_1,je_1,e_3,je_3\}`$. Putting $`e_2=je_1`$, $`e_4=je_3`$ we obtain the Cayley frame. Using the formulas (23) to define $`\{u_i\}`$ we find a smooth unitary frame of $`TM_{|N}`$ with the desidered properties.
Q.D.E.
We let $``$ denote the Levi-Civita connection of the Kähler metric $`g`$ on $`M`$, and $`D`$ the induced connection on the submanifold $`N`$. The metric $`g`$ being Kähler, $``$ gives a connection on $`K_M`$, and this in turn can be pulled back to a connection on $`K_M_{|N}`$. We denote both connections by $``$, too. Let $`\nu _N`$ denote the normal bundle to $`NM`$, $`h:TNTN\nu _N`$ the II fundamental form, and $`\stackrel{}{H}`$ the mean curvature vector.
###### Proposition 1
If $`N`$ is a Cayley submanifold
$$g(h(X,Y),JZ)g(h(X,Z),JY)=\left(D_X\omega \right)(Z,Y)$$
(27)
$$\omega (X,\stackrel{}{H})=\underset{i=1}{\overset{4}{}}g(h(X,e_i),Je_i)$$
(28)
where $`X,Y,Z`$ are arbitrary vectors tangent to $`N`$, and $`\{e_i\}`$ is any orthonormal basis of $`TN`$.
P r o o f.
$$\begin{array}{cc}\hfill g(h(X,Y),JZ)& =g((_XY)^{},JZ)=\hfill \\ & =g(_XY,JZ)g(D_XY,JZ)=\hfill \\ & =g(_XY,JZ)\omega (Z,D_XY),\hfill \end{array}$$
(29)
therefore
$$\begin{array}{c}g(h(X,Y),JZ)g(h(X,Z)JY)=\hfill \\ \hfill =g(_XY,JZ)g(_XZ,JY)+\omega (Y,D_XZ)\omega (Z,D_XY)\end{array}$$
(30)
now
$$\begin{array}{cc}\hfill g(_XY,JZ)g& (_XZ,JY)=\hfill \\ & =Xg(Y,JZ)g(Y,J_XZ)g(_XZ,JY)=\hfill \\ & =X\omega (Z,Y).\hfill \end{array}$$
(31)
Therefore
$$\begin{array}{cc}\hfill g(h(X,Y),JZ)g& (h(X,Z)JY)=\hfill \\ & =X\omega (Z,Y)\omega (D_XZ,Y)\omega (Z,D_XY)=\hfill \\ & =\left(D_X\omega \right)(Z,Y).\hfill \end{array}$$
(32)
This proves (27).
The second formula follows by taking the trace,
$$\underset{i=1}{\overset{4}{}}\left\{g(h(e_i,X)Je_i)g(h(e_i,e_i),JX)\right\}=\underset{i=1}{\overset{4}{}}\left(D_{e_i}\omega \right)(e_i,X),$$
and
$$\underset{i=1}{\overset{4}{}}\left(D_{e_i}\omega \right)(e_i,X)=d^{}\left(\omega _{|N}\right)(X).$$
$`N`$ being Cayley, the restriction of $`\omega `$ to $`N`$ is selfdual (and closed), hence coclosed (with respect to the metric $`g_{|N}`$). Therefore $`d^{}\omega _{|M}=0`$, and
$$\underset{i=1}{\overset{4}{}}g(h(e_i,X)Je_i)=\underset{i=1}{\overset{4}{}}g(h(e_i,e_i),JX)=\omega (X,\stackrel{}{H}).$$
(33)
Q.D.E.
We will now use a Cayley frame and a unitary basis as in 6, defined in some open subset of the totally real part of $`N`$, to prove some formulas relating $`\lambda `$ and $`\stackrel{}{H}`$.
###### Lemma 7
$$\mathrm{\Omega }_N(u_1,\mathrm{},_Xu_k,\mathrm{},u_4)=\sqrt{1}g(_Xu_k,Ju_k).$$
(34)
P r o o f.
$$_Xu_k=\underset{j=1}{\overset{4}{}}\left\{g(_Xu_k,u_j)u_j+g(_Xu_k,Ju_j)Ju_j\right\}$$
therefore, using the fact that $`\mathrm{\Omega }_N`$ is a complex form, i.e. of type (4,0), we compute
$$\begin{array}{cc}\hfill \mathrm{\Omega }_N(u_1,\mathrm{},_Xu_k,& \mathrm{},u_4)=\hfill \\ & =\underset{j}{}g(_Xu_k,u_j)\mathrm{\Omega }_N(u_1,\mathrm{},u_j,\mathrm{},u_4)+\hfill \\ & +\underset{j}{}g(_Xu_k,Ju_j)\mathrm{\Omega }_N(u_1,\mathrm{},Ju_j,\mathrm{},u_4)=\hfill \\ \hfill =& g(_Xu_k,u_k)\mathrm{\Omega }_N(u_1,\mathrm{},u_k,\mathrm{},u_4)+\hfill \\ & +g(_Xu_k,Ju_k)\sqrt{1}\mathrm{\Omega }_N(u_1,\mathrm{},u_k,\mathrm{},u_4)=\hfill \\ & =\sqrt{1}g(_Xu_k,Ju_k)\mathrm{\Omega }_N(u_1,\mathrm{},u_k,\mathrm{},u_4)\hfill \end{array}$$
because
$$g(_Xu_k,u_k)=\frac{1}{2}Xu_k^2=0.$$
Q.D.E.
###### Lemma 8
$$(1\lambda ^2)\underset{k=1}{\overset{4}{}}g(_Xu_k,Ju_k)=\omega (X,\stackrel{}{H}).$$
(35)
P r o o f.
Let us use the definition (23) of $`u_j`$:
$`\begin{array}{cc}\hfill _X& u_2=\hfill \\ & =\left(X{\displaystyle \frac{1}{\sqrt{1\lambda ^2}}}\right)(e_2\lambda Je_1)+{\displaystyle \frac{1}{\sqrt{1\lambda ^2}}}\left(_Xe_2(X\lambda )Je_1\lambda J_Xe_1\right)\hfill \end{array}`$
$`\begin{array}{cc}\hfill g(_X& u_2,Ju_2)=\hfill \\ & =\left(X{\displaystyle \frac{1}{\sqrt{1\lambda ^2}}}\right){\displaystyle \frac{1}{\sqrt{1\lambda ^2}}}g(e_2\lambda Je_1,J\left(e_2\lambda Je_1\right))+\hfill \\ & +{\displaystyle \frac{1}{1\lambda ^2}}g(_Xe_2(X\lambda )Je_1\lambda J_Xe_1,Je_2+\lambda e_1)=\hfill \\ & =\mathrm{\hspace{0.17em}0}+{\displaystyle \frac{1}{1\lambda ^2}}g(_Xe_2,Je_2)+{\displaystyle \frac{\lambda }{1\lambda ^2}}g(_Xe_2,e_1)+\hfill \\ & {\displaystyle \frac{X\lambda }{1\lambda ^2}}g(Je_1,Je_2){\displaystyle \frac{X\lambda \lambda }{1\lambda ^2}}g(Je_1,e_1)+\hfill \\ & {\displaystyle \frac{\lambda }{1\lambda ^2}}g(_Xe_1,e_2)+{\displaystyle \frac{\lambda ^2}{1\lambda ^2}}g(_Xe_1,Je_1)\hfill \end{array}`$
therefore
$`\begin{array}{cc}\hfill (1\lambda ^2)g(_Xu_2,Ju_2)& =g(_Xe_2,Je_2)+\lambda ^2g(_Xe_1,Je_1)+\hfill \\ & \lambda g(_Xe_1,e_2)+\lambda g(_Xe_2,e_1)=\hfill \\ & =g(_Xe_2,Je_2)+\lambda ^2g(_Xe_1,Je_1)+\hfill \\ & g(_Xe_1,Be_1)g(_Xe_2,Be_2)\hfill \end{array}`$
$`\begin{array}{cc}\hfill (1\lambda ^2)[g(_Xu_1,Ju_1& )+g(_Xu_2,Ju_2)]=\hfill \\ & =g(_Xe_1,Je_1)+g(_Xe_2,Je_2)+\hfill \\ & g(_Xe_1,Be_1)g(_Xe_2,Be_2)=\hfill \\ & =g(_Xe_1,(Je_1)^{})+g(_Xe_2,(Je_2)^{})=\hfill \\ & =g(h(X,e_1),Je_1)+g(h(X,e_2)Je_2).\hfill \end{array}`$
The same computation works for the last two indices, 3 and 4. Summing the two terms and using(28) one gets
$$(1\lambda ^2)\underset{k=1}{\overset{4}{}}g(_Xu_k,Ju_k)=\underset{i}{}g(h(X,e_i),Je_i)=\omega (X,\stackrel{}{H}).$$
Q.D.E.
###### Proposition 2
$$\sqrt{1}\left(_X\mathrm{\Omega }_N\right)(\stackrel{}{u})=\frac{i_\stackrel{}{H}\omega }{\lambda ^21}(X).$$
(36)
P r o o f.
By construction $`\mathrm{\Omega }_N(\stackrel{}{u})1`$.
$$0=X.\mathrm{\Omega }_N(\stackrel{}{u})=\left(_X\mathrm{\Omega }\right)(\stackrel{}{u})+\underset{k=1}{\overset{4}{}}\mathrm{\Omega }_N(u_1,\mathrm{},_Xu_k,\mathrm{},u_4).$$
Using the last two lemmas
$$\sqrt{1}\left(_X\mathrm{\Omega }\right)(\stackrel{}{u})=\underset{k=1}{\overset{4}{}}g(_Xu_k,Ju_k)=\frac{\omega (X,\stackrel{}{H})}{1\lambda ^2}.$$
Q.D.E.
Let $`\rho `$ denote the Ricci form of $`\omega `$ and let us define $`\gamma ^1(N_r)`$ by
$$\begin{array}{cc}\hfill \gamma (X)& :=\sqrt{1}\left(_X\mathrm{\Omega }_N\right)(\stackrel{}{u})=\hfill \\ & =\frac{i_\stackrel{}{H}\omega }{\lambda ^21}(X)=\underset{k=1}{\overset{4}{}}g(_Xu_k,Ju_k).\hfill \end{array}$$
(37)
###### Theorem 3
$$d\gamma =\rho _{|N}.$$
(38)
P r o o f.
$$X\left(\left(_Y\mathrm{\Omega }_N\right)(\stackrel{}{u})\right)=\left(_X_Y\mathrm{\Omega }_N\right)(\stackrel{}{u})+\left(_Y\mathrm{\Omega }_N\right)(_X\stackrel{}{u}).$$
$`(M,g)`$ being Kähler, $`\mathrm{\Omega }_N^{4,0}`$ implies that $`_Y\mathrm{\Omega }_N^{4,0}`$ too, therefore the same computations as above apply:
$$\begin{array}{cc}\hfill _Y\mathrm{\Omega }_N(u_1,\mathrm{},& _Xu_k,\mathrm{},u_4)=\hfill \\ & =\underset{j=1}{\overset{4}{}}\{g(_Xu_k,u_j)\left(_Y\mathrm{\Omega }_N\right)(u_1,\mathrm{},u_j,\mathrm{},u_4)+\hfill \\ & +g(_Xu_k,Ju_j)\left(_Y\mathrm{\Omega }_N\right)(u_1,\mathrm{},Ju_j,\mathrm{},u_4)\}=\hfill \\ & =g(_Xu_k,u_k))(_Y\mathrm{\Omega }_N\left)\right(u_1,\mathrm{},u_k,\mathrm{},u_4)+\hfill \\ & +\sqrt{1}\underset{j=1}{\overset{4}{}}\left\{g(_Xu_k,Ju_j)\left(_Y\mathrm{\Omega }_N\right)(u_1,\mathrm{},u_j,\mathrm{},u_4)\right\}=\hfill \\ & =\sqrt{1}g(_Xu_k,Ju_k)\left(_Y\mathrm{\Omega }_N\right)(\stackrel{}{u}),\hfill \end{array}$$
$$\left(_Y\mathrm{\Omega }_N\right)(_X\stackrel{}{u})=\sqrt{1}\left[\underset{k}{}g(_Xu_k,Ju_k)\right]\left(_Y\mathrm{\Omega }_N\right)(\stackrel{}{u})=\gamma (X)\gamma (Y).$$
Then applying the usual formula for the differential of a 1-form we find
$$d\gamma (X,Y)=\sqrt{1}\left(\mathrm{R}_{XY}\mathrm{\Omega }_N\right)(\stackrel{}{u})$$
but
$$\mathrm{R}_{XY}\mathrm{\Omega }_N=\sqrt{1}\rho (X,Y)\mathrm{\Omega }_N.$$
Q.D.E.
Finally we make the following remark.
###### Proposition 3
Let $`M`$ be a Kähler manifold and let $`N,N^{}`$ be two closed Cayley submanifolds in the same homology class,
$$[N]=[N^{}]H_4(M,).$$
Then, if $`N`$ is Lagrangian, the same is true of $`N^{}`$.
P r o o f.
It is enough to observe that
$$\lambda _N_{L^2(N)}^2=_N\lambda _N^2\mathrm{dvol}=\frac{1}{2}_N\omega ^2=\frac{1}{2}<\omega ^2,[N]>$$
is a topological invariant. If $`N`$ is Lagrangian, $`\lambda _N^{}_{L^2}^2=\lambda _N_{L^2}^2=0`$, therefore $`\lambda _N^{}0`$ and $`N^{}`$ is Lagrangian.
Q.D.E.
## 4 Minimal Cayley submanifolds in Kähler-Einstein manifolds
We now apply the formula (38) to the cases where $`\rho =s\omega `$, i.e. when the ambient manifold is Kähler-Einstein. This will yield proofs of theorems 1 and 2.
P r o o f of theorem 1.
Let $`\mathrm{\Omega }\mathrm{\Gamma }(K_M)`$ be a parallel normalized (4,0)-form. On $`N`$ we can write
$$\mathrm{\Omega }_N=e^{\sqrt{1}\theta }\mathrm{\Omega }$$
for some locally defined real valued function $`\theta =\theta (x)`$. Then
$`_X\mathrm{\Omega }_N`$ $`=_X\left(e^{\sqrt{1}\theta }\mathrm{\Omega }\right)=`$
$`=\sqrt{1}(X\theta )e^{\sqrt{1}\theta }\mathrm{\Omega }=`$
$`=\sqrt{1}(X\theta )\mathrm{\Omega }_N.`$
$`\gamma (X)`$ $`=\sqrt{1}\left(_X\mathrm{\Omega }_N\right)(\stackrel{}{u})=`$
$`=(X\theta )\mathrm{\Omega }_N(\stackrel{}{u})=`$
$`=X\theta `$
i.e. $`\gamma =d\theta `$. But $`\gamma =0`$ because $`\stackrel{}{H}=0`$. Therefore $`\theta \theta _0`$ is a constant, and $`N`$ is calibrated by
$$\mathrm{\Phi }_0=\mathrm{Re}\left(e^{\sqrt{1}\theta _0}\mathrm{\Omega }\right)+\frac{\omega ^2}{2}.$$
On the other side, if $`N`$ is calibrated by some parallel Cayley calibration, then it is obviously minimal, and thanks to lemma 5 (c), it is Cayley also according to our definition.
Q.D.E.
P r o o f of theorem 2.
Let $`\rho =s\omega `$, with $`s0`$.
If $`NM`$ is not complex, then $`N_r`$ is not empty. But $`\stackrel{}{H}0`$ implies that $`\gamma `$, hence $`d\gamma `$ vanish identically. Therefore
$$\omega _{|N_r}=\frac{1}{s}\rho _{|N_r}=d\gamma =0$$
i.e. $`N_r`$ is a Lagrangian submanifold, and $`\lambda =0`$ on it. This means that $`N_r=\lambda ^1(0)`$ is a closed and open set. Then $`N=N_r`$, and $`N`$ is Lagrangian.
Q.D.E. |
warning/0002/hep-th0002020.html | ar5iv | text | # 1 Introduction and Conclusions
## 1 Introduction and Conclusions
It has been recently appreciated that the dynamics of D-branes in a constant background $`B`$-field admits two equivalent descriptions: either in terms of an ‘ordinary’ gauge theory, or in terms of a gauge theory on a non-commutative worldvolume. The two descriptions arise as a result of using two different regularizations for the open string worldsheet theory . To provide evidence that both descriptions are equivalent non-perturbatively in the non-commutativity parameters is the purpose of this paper.
If Pauli-Villars regularization is used, the effective action for the massless open string fields enjoys an ordinary $`U(N)`$ gauge symmetry, where $`N`$ is the number of overlapping branes. This gauge symmetry is therefore commutative in the case that one single D-brane is present, which will be the only case considered in the present paper. The (bosonic) massless fields are a gauge potential $`A_\mu `$ ($`\mu =0,\mathrm{},p`$) and a number of scalars $`X^i`$ ($`i=p+1,\mathrm{},9`$). With Pauli-Villars regularization, the parameters entering the effective action are the closed string metric $`g`$, the closed string coupling constant $`g_s`$ and the constant background $`B`$-field itself. Furthermore, the whole dependence of the effective action on $`B`$ is accounted for by writing it in terms of the modified field strength $`F+B^{}`$, where $`F=dA`$ and the superscript ‘’ denotes the pull-back to the brane worldvolume. This is the combination which is invariant under the two $`U(1)`$ gauge symmetries which the open string $`\sigma `$-model enjoys at the classical level: a first one under which
$$\delta B=0,\delta A=d\lambda ,$$
(1.1)
and a second one under which
$$\delta B=d\mathrm{\Lambda },\delta A=\mathrm{\Lambda }^{}.$$
(1.2)
On the contrary, if point-splitting regularization is used, the gauge symmetry group of the effective action becomes non-commutative even in the case of one single D-brane. The theory is now naturally formulated as a gauge theory on a non-commutative worldvolume , namely a worldvolume whose coordinates $`\sigma ^\mu `$ do not commute but satisfy
$$[\sigma ^\mu ,\sigma ^\nu ]\sigma ^\mu \sigma ^\nu \sigma ^\nu \sigma ^\mu =i\theta ^{\mu \nu }.$$
(1.3)
The antisymmetric matrix $`\theta `$ appearing above measures the non-commutativity of the theory and defines the so-called ‘$``$-product’ as
$$(fg)(\sigma )e^{\frac{i}{2}\theta ^{\mu \nu }\frac{}{\xi ^\mu }\frac{}{\chi ^\nu }}f(\sigma +\xi )g(\sigma +\chi )|_{\xi =\chi =0}.$$
(1.4)
The only parameters entering the effective action when point-splitting regularization is used are $`\theta `$, an open string metric $`G`$ and an open string coupling constant $`G_s`$, whose relation with the closed string parameters is <sup>1</sup><sup>1</sup>1We will set $`2\pi \alpha ^{}=1`$.
$$\theta =(g+B)_{(A)}^1,G=gBg^1B,G_s=g_s\left(\frac{det(g+B)}{detg}\right)^{1/2},$$
(1.5)
where $`(A)`$ stands for the antisymmetric part of the matrix. All products of fields in the effective action are now $``$-products, and $`\theta `$ enters the action only through the definition of the $``$-product. The non-commutative field strength is defined as
$`\widehat{F}_{\mu \nu }`$ $`=`$ $`_{[\mu }\widehat{A}_{\nu ]}i[\widehat{A}_\mu ,\widehat{A}_\nu ],`$
$`[\widehat{A}_\mu ,\widehat{A}_\nu ]`$ $`=`$ $`\widehat{A}_\mu \widehat{A}_\nu \widehat{A}_\nu \widehat{A}_\mu ,`$ (1.6)
where $`\widehat{A}_\mu `$ is the non-commutative gauge field. With this regularization, the gauge symmetry under which the effective action is invariant is
$$\delta \widehat{A}=d\widehat{\lambda }+i[\widehat{\lambda },\widehat{A}],\delta \widehat{F}=i[\widehat{\lambda },\widehat{F}].$$
(1.7)
As explained in , since both the ordinary and the non-commutative descriptions arise from different regularizations of the same worldsheet theory, there should be a field redefinition which maps gauge orbits in one description into gauge orbits in the other. This requirement, plus locality to any finite order in $`\theta `$, enabled the authors in to establish the following system of differential equations: <sup>2</sup><sup>2</sup>2The two equations (1.8) are simply the dimensional reduction of that for a ten-dimensional gauge field $`\widehat{A}_M`$ $`(M=0,\mathrm{},9)`$, which was given in .
$`\delta \widehat{A}_\mu =\delta \theta ^{\alpha \beta }{\displaystyle \frac{\widehat{A}_\mu }{\theta ^{\alpha \beta }}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta \theta ^{\alpha \beta }\{\widehat{A}_\alpha ,_\beta \widehat{A}_\mu +\widehat{F}_{\beta \mu }\},`$
$`\delta \widehat{X}^i=\delta \theta ^{\alpha \beta }{\displaystyle \frac{\widehat{X}^i}{\theta ^{\alpha \beta }}}`$ $`=`$ $`{\displaystyle \frac{1}{4}}\delta \theta ^{\alpha \beta }\{\widehat{A}_\alpha ,_\beta \widehat{X}^i+D_\beta \widehat{X}^i\},`$ (1.8)
where $`\{f,g\}fg+gf`$ and $`D_\mu X_\mu \widehat{X}i[\widehat{A}_\mu ,\widehat{X}]`$. These equations, which we will call $`\theta `$-evolution equations, determine how the fields should change when $`\theta `$ is varied, in order to describe equivalent physics. Their integration provides the desired map between two descriptions with different values of $`\theta `$, to which we will refer as the Seiberg-Witten map.
The fact that two apparently so different descriptions can be equivalent is certainly remarkable. The authors in provided a direct check that this is indeed the case by showing that the effective action in one description is mapped to the effective action in the other by the Seiberg-Witten map. However, they worked in the approximation of slowly-varying fields, which consists of neglecting all terms of order $`F`$ (or $`^2X`$). This approximation was used at two different stages. First, when the effective action for the massless fields on the brane was taken to be the Dirac-Born-Infeld (DBI) action. Indeed, this action can be derived from string theory precisely by neglecting such terms. Second, this approximation was also used to simplify the Seiberg-Witten map considerably.
Although this check provides direct evidence in favour of the equivalence of the two descriptions, it would be desirable to have an exact proof. This would consist of three steps. First, one would have to determine the effective action in each description exactly, namely to all orders in $`\alpha ^{}`$. Second, one would have to integrate (1.8) also exactly, namely to all orders in $`\theta `$. Third, one would have to substitute the change of variables in one action and see that the other one is recovered. Of course, this procedure is impossible to put into practice, but we will see that it is still possible to provide some evidence that the Seiberg-Witten map works non-perturbatively in $`\theta `$.
The idea is as follows. If one exact effective action is mapped to the other by the Seiberg-Witten map, then a classical solution of one action should also be mapped to a classical solution of the other. Of course, since we do not know the exact effective actions, we do not know any non-trivial exact solutions either, except for one case: the BIon . This is a 1/2-supersymmetric solution of the ordinary DBI theory of a D-brane. In general, supersymmetric solutions of the worldvolume theories of branes have a natural interpretation as intersections of branes. The BIon is the prototype of this fact: it is the worldvolume realization of a fundamental string ending on a D-brane. What makes the BIon solution special is that, although originally discovered as a solution of the DBI action , it has been shown to be a solution of the exact effective action to all orders in $`\alpha ^{}`$ . Perhaps this should not be surprising since, after all, the fact that a fundamental string can end on a D-brane is the defining property of D-branes .
When no background $`B`$-field is present, the string ends orthogonally on the brane (see figure 1(a)). When a constant electric <sup>3</sup><sup>3</sup>3An electric $`B`$-field leads, through (1.5), to an electric $`\theta `$ in the non-commutative theory. There is some controversy about whether or not such a theory makes sense. We will postpone the discussion of this issue until the last section. $`B`$-field is turned on, supersymmetry requires the string to tilt a certain angle $`\gamma `$ determined by $`B`$ (see figure 1(b)). This can be intuitively understood, because the background $`B`$-field induces a constant electric field on the brane <sup>4</sup><sup>4</sup>4Fundamental strings ending on D-branes with constant electric fields on their worldvolumes have also been studied in . We will clarify the relationship of the results in with ours at the end of section 2. . Since the endpoint of the string is electrically charged, the string is now forced to tilt in order for its tension to compensate the electric force on its endpoint.
Our strategy will be to identify a BIon-like solution of the effective action in the non-commutative description with the value of $`\theta `$ determined by $`B`$ as in (1.5). Since the exact effective action is not known, we will work with the lowest order approximation in $`\alpha ^{}`$ (see (3.5)). We will then integrate the $`\theta `$-evolution equations exactly to find what this non-commutative BIon is mapped to in the ordinary description. The result is that it is mapped to a tilted ordinary BIon, as described above, and that the tilting angle $`\gamma `$ agrees precisely with the value determined by $`B`$ <sup>5</sup><sup>5</sup>5The fact that tilted brane configurations are related to monopoles and dyons in non-commutative gauge theories was first pointed out in by working to first order in $`\theta `$..
Our result can be interpreted in two ways. On one hand, if one accepts that the equivalence between the ordinary and the non-commutative descriptions, as determined by the Seiberg-Witten map, is valid non-perturbatively in $`\theta `$, then the result proves that the non-commutative BIon is a solution of the exact non-commutative effective action, namely it is a solution to all orders in $`\alpha ^{}`$. On the other hand, motivated by the defining property of D-branes mentioned above, one could directly conjecture that the non-commutative BIon is a solution of the exact non-commutative effective action. In this case, one could regard our result as evidence in favour of the exact equivalence of both descriptions beyond the perturbative level in $`\theta `$. Whichever point of view one adopts, the result sheds some light on another related question raised in : the convergence of the series in $`\theta `$ generated by the equations (1.8). As far as we are aware, this series has only been shown to converge for the case of constant field strength . Our example shows that it converges in a less simple case.
In this paper we will concentrate on the case of D3-branes, but the analysis for BIons applies to D-branes of arbitrary dimension. The reason for this restriction is that we will briefly comment on solutions more general than BIons, namely on dyons. These are solutions of a D3-brane worldvolume theory carrying both electric and magnetic charges which constitute the worldvolume realization of $`(p,q)`$-strings ending on a D3-brane. They fill out orbits of the $`SL(2,\text{})`$ duality group of type IIB string theory which, on the worldvolume of a D3-brane, becomes an electromagnetic duality group. In the last section we will briefly comment on the possible role of $`SL(2,\text{})`$ in the non-commutative description of D3-branes. Section 2 is mainly a review. Sections 3 and 4 contain original results.
## 2 Ordinary BIons in a background electric $`B`$-field
In this section we will consider the ordinary description of the worldvolume theory of a D3-brane in the approximation of slowly-varying fields. The action is therefore the DBI action. We will begin by reviewing its dyonic solutions in the absence of a background $`B`$-field. Then we will concentrate on the case of purely electric solutions, the so-called ‘BIons’. Finally we will see how the BIons are modified when a constant electric background $`B`$-field is turned on. In this section we will work in the static gauge $`X^\mu \sigma ^\mu `$ ($`\mu =0,\mathrm{},3`$). $`X^\mu `$, together with the scalar fields on the brane $`X^i`$, ($`i=4,\mathrm{},9`$), are target-space Cartesian coordinates, that is, the closed string metric $`g`$ is assumed to take the form $`g_{MN}=\eta _{MN}`$ ($`M,N=0,\mathrm{}9`$) when expressed in these coordinates. Throughout this paper we will only allow for one scalar field to be excited. Therefore, we will consider the target-space to be effectively five-dimensional.
The dyonic solutions we wish to describe are the worldvolume realizations of $`(p,q)`$-strings ending orthogonally on the D3-brane. Since this is a spacetime-supersymmetric configuration, we have to look for a worldvolume-supersymmetric solution. Supersymmetric solutions of the D3-brane DBI action can be found either by imposing directly preservation of supersymmetry or by looking for BPS bounds on the energy . In either case, the BPS equations for static 1/2-supersymmetric dyons carrying electric and magnetic charges $`p`$ and $`q`$ respectively, are
$$𝑬=\mathrm{sin}\alpha \mathbf{}X,𝑩=\mathrm{cos}\alpha \mathbf{}X,$$
(2.1)
where $`𝑬`$ and $`𝑩`$ are the electric and magnetic fields on the brane, $`\mathrm{tan}\alpha =p/q`$, and $`X`$ is the only scalar field involved in the solution. The bound on the energy $``$ which the solutions of (2.1) saturate is
$$\sqrt{Z_{el}^2+Z_{mag}^2},$$
(2.2)
where the electric and magnetic charges above are
$$Z_{el}=_\mathrm{\Sigma }d^3\sigma 𝑬\mathbf{}X,Z_{mag}=_\mathrm{\Sigma }d^3\sigma 𝑩\mathbf{}X,$$
(2.3)
and $`\mathrm{\Sigma }`$ is the D3-brane worldspace. Since both $`𝑬`$ and $`𝑩`$ are divergence free,
$$\mathbf{}𝑬=0,\mathbf{}𝑩=0,$$
(2.4)
(the former as a consequence of the Gauss law and the latter because of the Bianchi identity), these charges can be rewritten as surface integrals over the boundary of the brane worldspace:
$$Z_{el}=_\mathrm{\Sigma }𝑑𝐒X𝑬,Z_{mag}=_\mathrm{\Sigma }𝑑𝐒X𝑩.$$
(2.5)
This ensures that the charges are ‘topological’, in the sense that they only depend on the boundary conditions imposed on the fields. It is this topological nature which guarantees that the saturation of the bound automatically implies the equations of motion.
We see from (2.1) and (2.4) that $`X`$ must be harmonic, that is, $`^2X=0`$. Given a harmonic function $`X`$, the electric and magnetic fields are determined by (2.1). A dyon is then associated with an isolated singularity of $`X`$.
We will concentrate for the rest of this section on the BIon. It corresponds to $`\mathrm{sin}\alpha =\pm 1`$ in (2.1), and therefore satisfies
$$𝑬=\mathbf{}X,$$
(2.6)
where we have chosen the minus sign for convenience. The most general $`SO(3)`$-symmetric solution is then given (up to a gauge transformation) by
$$X=A_0=\frac{e}{4\pi |𝝈|},𝑨=0,$$
(2.7)
where $`𝝈=(\sigma ^a)`$, $`a=1,2,3`$. It corresponds to a fundamental string ending orthogonally on the brane at $`𝝈=0`$, as depicted in figure 1(a). As mentioned above, (2.7) is a solution of (classical) string theory to all orders in $`\alpha ^{}`$ , that is, when all corrections to the DBI action involving higher derivatives of the fields are taken into account.
Since the BIon (2.7) saturates the BPS bound (2.6), its energy equals its charge. Furthermore, the latter is easily calculated with the help of (2.5). The boundary of the BIon worldspace consists of a two-sphere at $`|𝝈|\mathrm{}`$ and another at $`|𝝈|0`$. The surface integral over the former vanishes. Over the latter, it diverges. We can regularize it by integrating over a small two-sphere $`S_ϵ`$ of radius $`ϵ`$. Since $`X`$ is constant on this sphere we are left with
$$=|Z_{el}|=\underset{ϵ0}{lim}\left|X(ϵ)_{S_ϵ}𝑑𝐒𝑬\right|=e\underset{ϵ0}{lim}X(ϵ).$$
(2.8)
This shows that the energy of the BIon is precisely the energy of an infinite string of constant tension . To compare with the BIon in the presence of a $`B`$-field and with the non-commutative BIon, it will be convenient for us to rewrite (2.8) as
$$=|Z_{el}|=z_{el}\underset{ϵ0}{lim}I(ϵ),$$
(2.9)
where
$$z_{el}=e^2,I(ϵ)=\frac{1}{4\pi ϵ}.$$
(2.10)
Let us now see how the BIon solution is modified when a constant electric background $`B`$-field is turned on. We assume that the target-space ten-vector $`B_{0M}`$ no longer vanishes (but we still impose the restriction that $`B`$ has no magnetic components). In general, any non-vanishing component of $`B`$ along directions ‘transverse’ to the brane can be gauged away. In our case, however, we have to be careful, because we are looking for a configuration in which one scalar field is excited. In other words, the worldspace of the brane is not a flat three-plane extending along the directions 1, 2 and 3. Therefore, we need to consider the component of $`B`$ along the direction labelled by $`X`$ (see figure 2(a)), to which we will refer as $`B_X`$. Without loss of generality, we can take the component of $`B`$ along the 123-space to point along the 1-direction. We will denote this component by $`B_\sigma `$. The remaining components of $`B_{0M}`$ can be gauged away and we will therefore set them to zero.
In the presence of the $`B`$-field, the BIon BPS equation which guarantees the preservation of some fraction of supersymmetry must be modified: the field strength $`F=dA`$ is replaced by $`=F+B^{}`$. This can be easily understood, since the supersymmetry condition has to be gauge-invariant under the two $`U(1)`$ gauge symmetries (1.1) and (1.2). Thus, (2.6) becomes
$$_{0a}=_aX,$$
(2.11)
whose solution is now
$$A_0=\frac{e}{4\pi |𝝈|},𝑨=0,X=\frac{1}{1+B_X}\frac{e}{4\pi |𝝈|}\frac{B_\sigma }{1+B_X}\sigma ^1.$$
(2.12)
Note the appearance of a term linear in the worldspace coordinate $`\sigma ^1`$ in the expression for the scalar field. This is the term responsible for the tilt of the string (see figure 2(b)). Indeed, although the spike coming out of the brane at $`|𝝈|=0`$ still points along the $`X`$-axis and the brane worldspace is still asymptotically flat, the latter no longer asymptotically extends along the 1-direction.
The energy of the solution (2.12) is computed analogously as we did with (2.7). The surface integral at infinity still vanishes (the term in the scalar field which is linear in $`\sigma ^1`$ does not contribute because it changes sign under $`𝝈𝝈`$). Thus, we obtain again (2.9), but with $`z_{el}`$ now given by
$$z_{el}=\frac{e^2}{1+B_X}$$
(2.13)
We would like to close this section with some clarifications. The first one is that it might appear that (2.12) is not the most general solution of (2.11): we could have included terms linear in the worldspace coordinate both in $`X`$ and in $`A_0`$, with appropriate coefficients such that (2.11) be satisfied. However, this apparently more general solution is related to (2.12) by a gauge transformation of the type (1.2); they are therefore physically equivalent. (Admittedly, such a transformation would shift the value of $`B`$, but this is allowed since $`B`$ is generic in our analysis.) In particular, by choosing the target-space one-form as
$$\mathrm{\Lambda }=\left(B_\sigma \sigma ^1+B_XX\right)d\sigma ^0,$$
(2.14)
the solution (2.12) in the presence of a non-vanishing $`B`$-field is mapped to a configuration with $`A_0=X`$ (where $`X`$ is still given by (2.12)) and vanishing $`B`$-field. This configuration satisfies (2.6), and is precisely of the form considered in . These considerations therefore clarify the relation between the solution studied in and the one we have presented: they are related by a gauge symmetry of the theory; whether the constant electric field on the brane is induced from the background or whether it arises from the worldvolume gauge field itself is a matter of gauge choice. We have chosen to work in a gauge in which $`B0`$ since it is in this case that a non-commutative alternative description exists.
The second remark is that the proof in that the BIon (2.7) is an exact solution to all orders in $`\alpha ^{}`$ assumed the absence of a background $`B`$-field. Therefore, one might question whether the result also holds for our BIon (2.12). The answer is affirmative, because, as we have just explained, (2.12) is related by a gauge symmetry of the theory to the configuration studied in in which $`B=0`$. As pointed out in , this latter configuration is indeed a solution to all orders in $`\alpha ^{}`$, because it satisfies (2.6), which was the only assumption in (as opposed to any assumption concerning the specific form of the solution, such as (2.7)).
## 3 Non-commutative D3-brane Dyons
In this section we will first establish the non-commutative version of the BPS equations (2.1). Then we will construct an exact solution for the case of purely electric charge, which we will refer to as the ‘non-commutative BIon’. In the next section, we will show that this solution is mapped to the ordinary BIon (2.12) in a $`B`$-field by the Seiberg-Witten map. Therefore, the metric we must use here is the open string metric $`G`$ as determined by $`g`$ and $`B`$ in (1.5). This is important since, although both $`g`$ and $`G`$ are flat, they are not necessarily simultaneously diagonal in the same coordinate system. In this section it will be convenient to work in a system in which $`G`$ takes the simplest form $`G_{MN}=\eta _{MN}`$ (note that we are referring to the target-space coordinate system, and that therefore it includes the scalar field on the brane). The system in which $`G`$ takes this form is obtained as follows (see figure 3).
Let us take the first worldspace coordinate $`\xi ^1`$ along the direction of $`B`$, and the scalar field $`Y`$ to be orthogonal to it. In this system we have
$$B=bd\xi ^0d\xi ^1,$$
(3.1)
where $`b`$ is a positive constant, but $`G`$ does not yet take the desired form. However, the simple rescaling
$$\rho ^0=\sqrt{1b^2}\xi ^0,\rho ^1=\sqrt{1b^2}\xi ^1$$
(3.2)
brings $`G`$ into such a form:
$$ds_{(G)}^2=d\rho _0^2+d\rho _1^2+d\rho _2^2+d\rho _3^2+dY^2$$
(3.3)
In this coordinate system $`\theta `$ also takes a very simple form: its only non-vanishing components are
$$\theta ^{01}=\theta ^{10}=b.$$
(3.4)
After these preliminaries, we are ready to write down the non-commutative D3-brane worldvolume action. We will take it to be the lowest order approximation in $`\alpha ^{}`$, namely
$$S=d^{1+3}\rho \left(\frac{1}{4}\widehat{F}_{\mu \nu }\widehat{F}^{\mu \nu }\frac{1}{2}D_\mu \widehat{Y}D^\mu \widehat{Y}\right).$$
(3.5)
All indices above are contracted with $`G`$.
One might think that a Hamiltonian analysis of the non-commutative action is required in order to obtain its energy functional, which we certainly need to derive a BPS bound on the energy. This seems difficult in view of the non-local nature of the action. (Note that this non-locality is not only in space but also in time, since we will not impose the restriction $`\theta ^{0a}=0`$.) The problem can be circumvented by noting that, although the presence of $`\theta `$ breaks Lorentz invariance, the theory is still translation-invariant. Therefore we can obtain the Lagrangian energy $``$ as the conserved quantity under time translations. The result is
$$=_\mathrm{\Sigma }d^3\rho \left(\frac{1}{2}\widehat{𝑬}^2+\frac{1}{2}\widehat{𝑩}^2+\frac{1}{2}(D_0\widehat{Y})^2+\frac{1}{2}(𝑫\widehat{Y})^2\right),$$
(3.6)
where $`\widehat{𝑬}`$ and $`\widehat{𝑩}`$ are the non-commutative electric and magnetic fields, that is,
$$\widehat{E}_a=\widehat{F}_{0a},\widehat{B}_a=\frac{1}{2}ϵ_{abc}\widehat{F}_{bc}.$$
(3.7)
For covariantly static configurations, namely with $`D_0\widehat{Y}=0`$, the energy can be rewritten as
$$=_\mathrm{\Sigma }d^3\rho \left[\frac{1}{2}\left(\widehat{𝑬}\mathrm{sin}\alpha 𝑫\widehat{Y}\right)^2+\frac{1}{2}\left(\widehat{𝑩}\mathrm{cos}\alpha 𝑫\widehat{Y}\right)^2+\mathrm{sin}\alpha \widehat{𝑬}𝑫\widehat{Y}+\mathrm{cos}\alpha \widehat{𝑩}𝑫\widehat{Y}\right],$$
(3.8)
from which the desired BPS bound follows immediately:
$$\sqrt{Z_{el}^2+Z_{mag}^2}.$$
(3.9)
The electric and magnetic charges above are the natural non-commutative generalizations of (2.3):
$$Z_{el}=_\mathrm{\Sigma }d^3\sigma \widehat{𝑬}𝑫\widehat{Y},Z_{mag}=_\mathrm{\Sigma }d^3\sigma \widehat{𝑩}𝑫\widehat{Y}.$$
(3.10)
The bound is saturated precisely when the non-commutative BPS equations
$$\widehat{𝑬}=\mathrm{sin}\alpha 𝑫\widehat{Y},\widehat{𝑩}=\mathrm{cos}\alpha 𝑫\widehat{Y}$$
(3.11)
hold <sup>6</sup><sup>6</sup>6The case with $`\mathrm{cos}\alpha =1`$ in (3.11), which corresponds to a monopole, was studied in in the context of a $`U(2)`$ gauge group. The non-commutative monopole equation was solved to first order in $`\theta `$. The solution exhibits a certain non-locality corresponding to the tilt of the D-string ending on the brane.. They are also the natural generalizations of their commutative counterparts (2.1).
As well as the BPS equations, the Gauss law and the Bianchi identity (2.4) should also be promoted to their non-commutative versions
$$𝑫\widehat{𝑬}=0,𝑫\widehat{𝑩}=0.$$
(3.12)
A number of comments are in order here. The new Bianchi identity follows immediately from the definition of $`\widehat{F}`$. However, it is not clear to us whether it constitutes a locally sufficient integrability condition (as it does in an ordinary gauge theory) which ensures that $`\widehat{F}`$ can be written as the covariant derivative of a gauge potential. The Gauss law above is nothing else than one of the equations of motion. However, a rigorous proof that it is a constraint in the non-commutative theory would require a careful analysis, which is difficult again due to its non-locality. Nevertheless, there are two observations which are worth noticing. First of all, at least in the case of purely magnetic $`\theta `$, that is, when $`\theta ^{0a}=0`$, the non-commutative theory becomes local in time and the Hamiltonian analysis is straightforward . In this case one can prove that the Gauss law (which contains only first order time derivatives in its Lagrangian form) becomes a Hamiltonian constraint. Second, both the Bianchi identity and the Gauss law are required in order to rewrite the electric and magnetic charges as surface integrals:
$$Z_{el}=_\mathrm{\Sigma }𝑑𝐒\widehat{Y}\widehat{𝑬},Z_{mag}=_\mathrm{\Sigma }𝑑𝐒\widehat{Y}\widehat{𝑩}.$$
(3.13)
As in the ordinary case, this is what ensures that they have a topological nature, which in turn guarantees that the saturation of the bound automatically implies the equations of motion.
Now we are ready to finally write down the non-commutative BIon solution. It is simple to check that the configuration
$$\widehat{A}_0^{(\rho )}=\widehat{Y}=\frac{q}{4\pi |𝝆|},\widehat{𝑨}^{(\rho )}=0,$$
(3.14)
solves both the non-commutative BPS equations (3.11) (with $`\mathrm{sin}\alpha =1`$ as before) and the non-commutative Gauss law (3.12). In (3.14) we have $`𝝆=(\rho ^a)`$, and the superscript $`(\rho )`$ denotes that the above are the components of the gauge potential one-form in the $`\rho `$-coordinate system, namely that $`\widehat{A}=\widehat{A}_0^{(\rho )}d\rho ^0`$. Its components in the $`\xi `$-coordinate system are, by virtue of (3.2),
$$\widehat{A}_0^{(\xi )}=\sqrt{1b^2}\widehat{A}_0^{(\rho )},\widehat{𝑨}^{(\xi )}=0.$$
(3.15)
Hereafter we will denote $`\widehat{A}_0^{(\rho )}`$ simply by $`\widehat{A}_0`$ until the moment in which we have to compare with (2.12).
The solution (3.14) is the announced non-commutative BIon. Its charge is again easily computed as we did in the previous section. The result is again (2.9) with
$$z_{el}=q^2.$$
(3.16)
and we see from (2.13) and (3.16) that we must choose
$$q=\frac{e}{\sqrt{1+B_X}}$$
(3.17)
in order for the charges (or equivalently, for the energies) of the non-commutative and of the ordinary BIon to agree.
## 4 From Non-commutative to Commutative BIons
Our goal in this section is to show that the solution (3.14) of the non-commutative theory (with non-commutativity parameter $`\theta `$ determined by $`B`$ as in (3.4)) is precisely mapped to the ordinary BIon (2.12) in the presence of the $`B`$-field. We should therefore integrate the $`\theta `$-evolution equations (1.8) exactly, with the initial conditions (3.14) for the fields at the initial value of $`\theta `$ given in (3.4). However, there is an a priori obstacle for doing so which is worth discussing. Indeed, the $`\theta `$-evolution equations constitute a system of coupled partial differential equations, whose (local) integrability conditions are that the crossed derivatives be equal, namely that
$$\frac{^2\widehat{A}_\mu }{\theta ^{\gamma \rho }\theta ^{\alpha \beta }}=\frac{^2\widehat{A}_\mu }{\theta ^{\alpha \beta }\theta ^{\gamma \rho }},$$
(4.1)
and similarly for the scalar fields. As proved in <sup>7</sup><sup>7</sup>7I thank Joan Simón for drawing my attention to this reference., these conditions are in general not satisfied. This means that the evolution of the fields in ‘$`\theta `$-space’ determined by integrating the equations (1.8) between some initial and final values $`\theta _0`$ and $`\theta _1`$ depends on the path followed from $`\theta _0`$ to $`\theta _1`$. The choice of path should be made according to some physical input <sup>8</sup><sup>8</sup>8As explained in there is no path dependence for the case of a $`U(1)`$ gauge group if terms of $`𝒪(\widehat{F})`$ are neglected..
In our case, we are interested in the $`\theta `$-evolution of the BIon (3.14) starting from an initial $`\theta `$ which is purely electric, and ending at $`\theta =0`$. This suggests that we should restrict ourselves to the hyperplane $`\theta ^{ab}=0`$ in $`\theta `$-space, and consider paths contained within this hyperplane. This restriction can be physically further motivated as follows. Suppose that we start with a D3-brane in the absence of any $`B`$-field. In this situation only the ordinary description is available. Now we smoothly turn on the electric components of $`B`$. As soon as $`B`$ is different from zero, both the ordinary and the non-commutative descriptions are available. It follows from its definition (1.5) that, in the latter description, $`\theta `$ will also be purely electric. We can now imagine following the evolution of the fields as a function of $`\theta `$ (or, equivalently, of $`B`$) as it increases. We see that in this physical situation we are following a path which lies in the $`\theta ^{ab}=0`$ hyperplane.
The restriction $`\theta ^{ab}=0`$ simplifies the $`\theta `$-evolution equations dramatically, because it implies that the $``$-product of any two time-independent functions $`f`$ and $`g`$ reduces to their ordinary product, namely that $`fg=fg`$. With this simplification the equations (1.8) become
$`{\displaystyle \frac{\widehat{A}_0}{\theta ^b}}`$ $`=`$ $`\widehat{A}_0_b\widehat{A}_0`$ (4.2)
$`{\displaystyle \frac{\widehat{A}_a}{\theta ^b}}`$ $`=`$ $`\widehat{A}_0_b\widehat{A}_a+{\displaystyle \frac{1}{2}}\widehat{A}_0_a\widehat{A}_b{\displaystyle \frac{1}{2}}\widehat{A}_b_a\widehat{A}_0,`$ (4.3)
$`{\displaystyle \frac{\widehat{Y}}{\theta ^b}}`$ $`=`$ $`\widehat{A}_0_b\widehat{Y},`$ (4.4)
where $`\theta ^a\theta ^{0a}`$. Note that in the above equations $`_b/\rho ^b`$.
The unique solution of (4.3) with the initial condition that $`\widehat{A}_a`$ vanishes at some initial value of $`\theta `$ is that it vanishes for all values of $`\theta `$, regardless of the chosen path. Furthermore, taking the derivative of (4.2) and using (4.2) itself, one can easily check that the crossed derivatives of $`\widehat{A}_0`$ coincide. Finally, (4.2) and (4.4), together with the initial condition $`\widehat{A}_0=\widehat{Y}`$, show that $`\widehat{A}_0`$ and $`\widehat{Y}`$ remain equal for all values of $`\theta `$. We thus conclude that, in the hyperplane $`\theta ^{ab}=0`$, the $`\theta `$-evolution of the BIon (3.14) is path-independent, and that to determine it we need only solve (4.2) with the initial condition
$$\widehat{A}_0(𝝆,𝜽)|_{𝜽=𝜽_0}=\frac{q}{4\pi |𝝆|},$$
(4.5)
where $`𝜽=(\theta ^a)`$ and $`𝜽_0=(b,0,0)`$.
(4.2) constitutes a system of three quasi-linear partial differential equations which can be solved by the method of characteristics:<sup>9</sup><sup>9</sup>9I am grateful to Emili Elizalde for help at this point. instead of solving (4.2) directly, one introduces a new three-vector $`𝒕=(t^a)`$ and solves
$$\frac{\rho ^b}{t^a}=\delta _a^b\widehat{A}_0,\frac{\theta ^b}{t^a}=\delta _a^b,\frac{\widehat{A}_0}{t^a}=0$$
(4.6)
with initial conditions that depend on a further three-vector $`𝒔=(s^a)`$, namely
$$𝝆=𝒔,𝜽=𝜽_0,\widehat{A}_0=\frac{q}{4\pi |𝒔|}\text{at}𝒕=0.$$
(4.7)
The physical interpretation of both $`𝒕`$ and $`𝒔`$ will become clear shortly.
(4.6) and (4.7) determine $`𝝆`$, $`𝜽`$ and $`\widehat{A}_0`$ as functions of $`𝒕`$ and $`𝒔`$. It is easy to see that the solution is
$$𝝆(𝒕,𝒔)=\frac{q}{4\pi |𝒔|}𝒕+𝒔,𝜽(𝒕,𝒔)=𝒕+𝜽_0,\widehat{A}_0(𝒕,𝒔)=\frac{q}{4\pi |𝒔|}.$$
(4.8)
The (local) inversion of the first two relations above would determine $`𝒕(𝝆,𝜽)`$ and $`𝒔(𝝆,𝜽)`$, which could be substituted into the third one to obtain the desired solution $`\widehat{A}_0(𝝆,𝜽)`$. To see that the function $`\widehat{A}_0`$ so determined satisfies (4.2) we simply have to apply the chain rule and use (4.6) to obtain:
$$0=\frac{\widehat{A}_0}{t^a}=\frac{\widehat{A}_0}{\theta ^b}\frac{\theta ^b}{t^a}+\frac{\widehat{A}_0}{\rho ^b}\frac{\rho ^b}{t^a}=\frac{\widehat{A}_0}{\theta ^a}+\widehat{A}_0\frac{\widehat{A}_0}{\rho ^a}$$
(4.9)
It is also obvious from (4.7) that $`\widehat{A}_0`$ satisfies the required initial condition (4.5).
We see from (4.8) how to interpret $`𝒕`$: it is simply the difference between the initial and the final values of the non-commutativity parameter. Therefore, in our case we will be interested in examining the values of the fields at
$$𝒕=(b,0,0).$$
(4.10)
(4.8) also shows that $`𝒔=(s^a)`$ are nothing else than intrinsic coordinates on the D3-brane worldspace. Indeed, recall that we argued that the scalar $`\widehat{Y}`$ and the potential $`\widehat{A}_0`$ coincide for all values of $`\theta `$. Therefore, for fixed $`𝒕`$, (4.8) determines the (static) embedding of the brane in spacetime as
$$𝒔(𝝆(𝒔),\widehat{Y}(𝒔))$$
(4.11)
So we see that (4.8) provides the solution (namely the gauge potential and the embedding of the brane in target-space) in parametric form.
In order to check that the solution (4.8) at $`\theta =0`$ is, as we claim, precisely the same as (2.12), we have to express (4.8) in the same coordinate system as (2.12).
First, we have to undo the rescaling (3.2), because the closed string metric $`g`$ takes the Minkowski form in the $`\xi `$-coordinates, but not in the $`\rho `$-coordinates. We thus have, using (3.2), (3.15), (4.8) and (4.10), that, at $`\theta =0`$:
$`\xi ^1`$ $`=`$ $`{\displaystyle \frac{b}{\sqrt{1b^2}}}{\displaystyle \frac{q}{4\pi |𝒔|}}+{\displaystyle \frac{1}{\sqrt{1b^2}}}s^1,Y={\displaystyle \frac{q}{4\pi |𝒔|}},`$ (4.12)
$`A_0`$ $`=`$ $`\sqrt{1b^2}{\displaystyle \frac{q}{4\pi |𝒔|}},𝑨=0.`$ (4.13)
Note that we have dropped the ‘hats’ on the fields, because the above are already their values in the ordinary description. Note also that, although we did not write the superscript ($`\xi `$) explicitly, $`A_0`$ and $`𝑨`$ now denote the components of the gauge potential in the $`\xi `$-coordinate system.
We have drawn the solution (4.12) (which is still expressed in parametric form) in figure 3. For $`|𝒔|\mathrm{}`$ we have $`Y0`$. Therefore the brane is asymptotically flat and extends along the directions labeled by $`\xi ^1`$, $`\xi ^2`$ and $`\xi ^3`$, as shown in the figure. On the contrary, for $`|𝒔|0`$, we see that
$$Y\frac{\sqrt{1b^2}}{b}\xi ^1$$
(4.14)
This means that the spike comes out of the region $`|𝝃|0`$ in a direction at an angle $`\gamma `$ with the $`Y`$-axis (see figure 3 again), where
$$\mathrm{tan}\gamma =\frac{b}{\sqrt{1b^2}}$$
(4.15)
We have labelled this direction by $`X`$, and the orthogonal direction in the $`\xi ^1`$-$`Y`$ plane by $`\sigma ^1`$. Recall that the ordinary BIon solution (2.12) is written in the static gauge and in a coordinate system in which the spike points along the transverse scalar field. In our case, this happens precisely in the $`\sigma `$-$`X`$ coordinate system. Therefore, the last step we have to take is to rotate the solution (4.12) by an angle $`\gamma `$. Defining
$$\left(\begin{array}{c}\sigma ^1\\ X\end{array}\right)=\left(\begin{array}{cc}\mathrm{cos}\gamma & \mathrm{sin}\gamma \\ \mathrm{sin}\gamma & \mathrm{cos}\gamma \end{array}\right)\left(\begin{array}{c}\xi ^1\\ Y\end{array}\right)$$
(4.16)
we find that $`\sigma ^1=s^1`$ and that
$$X=\frac{1}{\sqrt{1b^2}}\frac{q}{4\pi |𝒔|}\frac{b}{\sqrt{1b^2}}s^1.$$
(4.17)
Finally, we see from (3.1) and (4.16) that
$$B_\sigma =b\sqrt{1b^2},B_X=b^2.$$
(4.18)
Therefore, using (3.17), we arrive at the final form of our solution, expressed in the static gauge:
$$A_0=\frac{e}{4\pi |𝝈|},X=\frac{1}{1+B_X}\frac{e}{4\pi |𝝈|}\frac{B_\sigma }{1+B_X}\sigma _1.$$
(4.19)
It coincides precisely with (2.12), as we claimed.
## 5 Discussion
In this last section we wish to address a number of issues which were not fully discussed in the text.
The first one concerns the interpretation of the scalar fields in the non-commutative theory. In the ordinary description of a single D-brane, the scalars unambigously determine the embedding of the brane in spacetime. This is the reason why many phenomena in field theories acquire a clearer geometrical interpretation when such theories are realized as worldvolume theories of branes. In the non-commutative description, this interpretation is much less clear. The reason is that the scalar fields, even in the case of one single D-brane, are no longer gauge-invariant but gauge-covariant quantities; namely they transform as $`\delta \widehat{X}=i[\widehat{\lambda },\widehat{X}]`$. This problem is related to the fact that all gauge-invariant quantities in non-commutative gauge theories seem to be non-local, obtained after integrating some gauge-covariant quantity. Presumably, one can determine global properties of the brane embedding from the non-commutative scalar fields, such as winding numbers, etc., by integrating appropriate expressions, but not the local details of the embedding. The reason why we did not have to resolve this problem in our discussion of the non-commutative BIon is that we were only interested in identifying a solution of the non-commutative theory which was exactly mapped to the ordinary BIon in the presence of a $`B`$-field. Since the Seiberg-Witten map maps gauge orbits into gauge orbits and the ordinary scalar fields are gauge-invariant, any scalar field configuration in the non-commutative theory which is gauge-equivalent to (3.14) would have also been mapped to the ordinary BIon. We simply chose the simplest representative of $`\widehat{Y}`$ in its gauge-equivalence class.
The second point we wish to discuss here is whether it makes sense to consider a non-commutativity matrix with non-vanishing electric components <sup>10</sup><sup>10</sup>10I would like to thank Michael B. Green and Paul K. Townsend for a conversation on this point.. In the ordinary description, it certainly makes sense to consider $`B_{0a}0`$, which leads through (1.5) to $`\theta ^{0a}0`$. Moreover, even if in one coordinate system $`B`$ has no electric components, we can always choose to describe the physics in a frame in which $`B`$ does have electric components. The only question is whether the equivalence between the ordinary and the non-commutative descriptions still holds in this frame. It seems that our result can be regarded as evidence in favour of the affirmative answer to this question.
We would like to close this paper with a little digression on the role of S-duality in the non-commutative theory, and more generally of the full $`SL(2,\text{})`$ duality group of type IIB string theory. Consider the BIon solution in the ordinary description of D3-branes and in the absence of any $`B`$-field. As we have explained, this is the worldvolume realization of the spacetime configuration in which a fundamental string ends on the D3-brane. The S-duality of string theory maps this configuration into one in which a D-string ends on the D3-brane. Thus, from the worldvolume point of view of the latter, S-duality is an inherited symmetry which corresponds to electromagnetic duality. It maps the BIon into the monopole. These two objects are therefore equivalent, in the sense that they are related by a symmetry of the theory. One might think that this is also the case in the non-commutative worldvolume description, perhaps with a further exchange of the electric and the magnetic components of $`\theta `$. However, this is not true. The reason is that a BIon in the non-commutative theory corresponds, as we have seen, to a fundamental string ending on the D3-brane in the presence of a $`B`$-field. S-duality maps this configuration into a D-string ending on the D3-brane in the presence now of a Ramond-Ramond $`C`$-field. Clearly, this does not correspond to a monopole in the non-commutative theory, which should instead correspond to a D-string ending on the D3 in the presence of a $`B`$-field. These considerations raise the following interesting question. There exists an $`SL(2,\text{})`$-covariant worldvolume action for D3-branes coupled to a supergravity background , which can consist, in particular, of a flat background with constant $`B`$ and $`C`$ fields. The case of vanishing $`C`$-field, for which we know that a non-commutative description is possible, is mapped to the generic one by $`SL(2,\text{})`$. Since $`SL(2,\text{})`$ is a symmetry of IIB string theory, should there not exist an $`SL(2,\text{})`$-covariant non-commutative description of D3-branes in the presence of both constant $`B`$ and $`C`$ fields?
Note added: While this paper was being written, I learned about , which has some overlap with section 3.
Acknowledgments: I am grateful to Joaquim Gomis for the suggestion which motivated this work, and to both him and Joan Simón for illuminating discussions and useful comments on the manuscript. I am also grateful to Selena Ng for her helpful comments on a previous version of this paper. Finally, I would like to thank Emili Elizalde, Josep I. Latorre, Josep M. Pons, Toni Mateos and Jordi Molins for discussions. This work has been supported by a fellowship from the Comissionat per a Universitats i Recerca de la Generalitat de Catalunya. |
warning/0002/nucl-th0002019.html | ar5iv | text | # Toward chiral theory of NN interactions in nuclear matter
## Abstract
We consider an effective field theory of NN system in nuclear medium. The shallow bound states, which complicate the effective field theory analysis in vacuum do not exist in matter. We show that the next-to-leading order terms in the chiral expansion of the effective Lagrangian can be interpreted as corrections so that the expansion is systematic. The Low Energy Effective Constants of this Lagrangian are found to satisfy the concept of naturalness. The potential energy per particle is calculated. The problems and challenges in constructing the chiral theory of nuclear matter are outlined.
Effective Field Theory (EFT) is now a standard method to study nuclear dynamics. EFT is based on the use the Lagrangian with the appropriate effective degrees of freedom instead of the fundamental ones in the low-energy region (for review of EFT see, for example ). This Lagrangian includes all possible terms allowed by the symmetries of the underlying theory. The states which can be treated as heavy, compared to the typical energy scale involved, are integrated out. They are hidden in the Low Energy Effective Constants (LEC’s) of Lagrangian. The physical amplitudes can be represented as the sum of certain graphs, each of them being of a given order in $`Q/\mathrm{\Lambda }`$, where $`Q`$ is a typical momentum scale and $`\mathrm{\Lambda }`$ is a scale of the short range physics. The relative contribution of each graph can roughly be estimated using counting rules . However, being applied to the NN system EFT encounters serious problem which is due to existence of the bound states near threshold . It results in the large nucleon-nucleon scattering length and makes the perturbative expansion divergent. Weinberg suggested to apply counting rules to the certain class of the irreducible diagrams which should then be summed up to infinite order by solving the Lippmann-Schwinger (LS) equation. The irreducible diagrams can be treated as the effective potential in this case. Different aspects of the chiral NN problem have been discussed since then . The EFT method has also been used to study nuclear matter . In the effective chiral Lagrangian was constructed and the “naturalness” of the effective coupling constants has been demonstrated. The possible counting rules for nuclear matter have been discussed in . These two lines of development of the chiral nuclear physics are in some sense similar to the tendencies existed some time ago in conventional nuclear physics with the phenomenological NN forces. On the one hand, the phenomenological NN potentials were used to describe nucleon-nucleon cross sections and phase shifts. On the other hand, nuclear mean field approaches provided a reasonable description of the bulk properties of nuclear matter. The unification of these two approaches then led to the famous Bethe-Goldstone (BG) equation for the G-matrix which is an analog of scattering T-matrix, satisfying the LS equation. So one may follow the same strategy and, being equipped with the chiral theory of NN interaction in vacuum, try to construct the chiral G-matrix. One can easily see the qualitative difference between vacuum and medium cases. In nuclear medium because of Pauli blocking the intermediate states with the momenta less than Fermi momentum $`p_F`$ are forbidden. Therefore, the nucleon propagator does not exhibit a pole. Moreover, the shallow bound NN states, being a serious problem in vacuum, simply do not exist in nuclear matter because of interaction of the NN pair with nuclear mean field. It means that the effective scattering length becomes considerably smaller compared to the vacuum one. . The moderate value of the in-medium scattering length would indicate that the typical scale of the NN interactions gets “more natural” in nuclear matter. We start from the standard nucleon-nucleon effective chiral Lagrangian
$$=N^{}i_tNN^{}\frac{^2}{2M}N\frac{1}{2}C_0(N^{}N)^2\frac{1}{2}C_2(N^{}^2N)(N^{}N)+h.c.+\mathrm{}.$$
(1)
We consider the NN scattering in the $`{}_{}{}^{1}S_{0}^{}`$ state. The G-matrix is
$$G(p^{},p)=V(p^{},p)+M\frac{dqq^2}{2\pi ^2}V(p^{},q)\frac{\theta (qp_F)}{M(ϵ_1(p)+ϵ_2(p^{}))q^2}G(q,p),$$
(2)
Here $`ϵ_1`$ and $`ϵ_2`$ are the single-particle energies of the bound nucleons. For the in-medium nucleon mass we used the value $`M=0.8M_0`$, where $`M_0`$ is the nucleon mass in vacuum. The standard strategy of treating the chiral NN problem in vacuum is the following. One computes amplitudes up to a given chiral order in the terms of the effective constants $`C_0`$ and $`C_2`$ which are then determined by comparing the calculated amplitude with some experimental data. Having these constants fixed one can calculate the other observables. We will follow the similar strategy in the nuclear matter case and proceed as follows. We take exactly solvable separable potential with parameters adjusted to the value of the potential energy per particle in nuclear matter and will consider the G-matrix obtained from this separable potential as our “observable”. Then we solve the BG equation with the effective constants $`C_0`$ and $`C_2`$. The numerical values of these constants are determined comparing the phenomenological and EFT G-matrix at some fixed kinematical points. The check of consistency we used is the difference between $`C_0`$’s determined in the leading and subleading orders. In the vacuum case the this difference was found to be large . Using a separable potential with the effective strength $`\lambda `$ and the form factors
$$\eta (p)=(p^2+\beta ^2)^{1/2}$$
(3)
One can easily find the solution of the corresponding BG equation
$$G(k,k)=\eta ^2(k)\left[\lambda ^1+\frac{M}{2\pi ^2}𝑑qq^2\frac{\theta (qp_F)\eta ^2(q)}{k^2q^2}\right]^1$$
(4)
We choose $`\lambda =2.4`$ and $`\beta =1.1\mathrm{fm}`$ to provide the potential energy per particle in a good agreement with the empirical value. The parameters $`\lambda `$ and $`\beta `$ being substituted in the G-matrix lead to $`a_mr_m0(1)`$, where $`a_m`$ and $`r_m`$ are the in-medium analogs of scattering length and effective radius. The absolute value of the in-medium scattering length is considerably reduced compared to the vacuum one. It clearly indicates that, as expected, the shallow virtual nucleon-nucleon bound state is no longer present in nuclear medium. Thus, one can avoid significant part of the difficulties typical for the chiral NN problem in vacuum. Having determined the phenomenological G-matrix one can now solve the BG equation taking into account leading and sub-leading orders of the NN effective chiral Lagrangian. The solution is similar to the vacuum case
$$\frac{1}{G(k,p_F)}=\frac{(C_2I_3(k,p_F)1)^2}{C_0+C_2^2I_5(k,p_F)+k^2C_2(2C_2I_3(k,p_F))}I(k,p_F),$$
(5)
where we defined
$$I_n\frac{M}{2\pi ^2}𝑑qq^{n1}\theta (qp_F);I(k)\frac{M}{2\pi ^2}𝑑q\frac{q^2\theta (qp_F)}{k^2q^2}.$$
(6)
These integrals are divergent so the renormilization should be carried out. The procedure used is similar to that adopted in Ref. to study the EFT approach to the NN interaction in vacuum. We subtract the divergent integrals at some kinematical point $`p^2=\mu ^2`$. After subtraction the renormalized G-matrix takes the form
$$\frac{1}{G^r(k,k)}=\frac{1}{C_0^r(\mu )+2k^2C_2^r(\mu )}+\frac{M}{4\pi }[p\mathrm{log}\frac{p_Fp}{p_F+p}i\mu \mathrm{log}\frac{p_Fi\mu }{p_F+i\mu }],$$
(7)
One notes that in the $`p_F0`$ limit the vacuum chiral NN amplitude is recovered. We choose the value $`\mu `$ = 0 as a subtraction point. The $`\mu `$ dependence of LEC’s is governed by the renormalization group (RG) equations. We demand that the entire G-matrix is independent of substruction point. After differentiating the G-matrix with respect to $`\mu `$ and setting $`G/\mu `$ = 0 one can get the following RG equations
$$\frac{C_0(\mu )}{\mu }=\frac{C_0^2M}{4\pi ^2}(\frac{\mu p_F}{p_F^2+\mu ^2}+2tan^1(\frac{\mu }{p_F}))$$
(8)
$$\frac{C_2(\mu )}{\mu }=\frac{C_0C^2M}{\pi ^2}(\frac{\mu p_F}{p_F^2+\mu ^2}+tan^1(\frac{\mu }{p_F}))$$
(9)
In the limit $`p_F`$ 0 these equations transform to the ones derived by Kaplan et al. . Now one can determine the LEC’s by equating the EFT and phenomenological G-matrices at some kinematical points. We used the values $`p=\frac{p_F}{2};\frac{p_F}{3}`$ as such points. The assumed value of the Fermi-momentum is $`p_F`$ = 1.37 fm. In the following we will omit the label “r” implying that we always deal with renormalized quantities. We found $`C_0=1.86fm^2`$ in LO. In NLO one gets $`C_0=2.7fm^2`$ and $`C_2=0.84fm^4`$ so that the inclusion of the NLO corrections give rise to the approximately 40$`\%`$ change in the value of $`C_0`$. It indicates that the chiral expansion is systematic in a sense that adding of the NLO terms in the effective Lagrangian brings in a “NLO change” of the coefficients which have already been fixed at LO. The natural size of the in-medium scattering length and moderate changes experienced by the coupling constant $`C_0`$ might, in principle, indicate the possibility of the perturbative calculations. However, in spite of this, it is still more useful to treat this problem in the nonperturbative manner since the corrections themselves are quite significant. Moreover, the overall (although distant) goal of the EFT description is to derive both nuclear matter and the vacuum NN amplitude from the same Lagrangian. However, it is hard to say at what densities the dynamics becomes intrinsically nonperturbative, so it is better to treat the problem nonperturbatively from the beginning. Let’s now calculate the potential energy per particle using the expression
$$U_{tot}=\frac{1}{2}\underset{\mu ,\nu }{}<\mu \nu |G(ϵ_\mu +ϵ_\nu )|\mu \nu \nu \mu )$$
(10)
The summation goes over the states with momenta below $`p_F`$. The LO perturbative calculations (where $`GC_0`$) give $`\frac{U(^1S_0)}{A}12MeV`$ The calculations using the lowest order $`G`$-matrix result in the value $`\frac{U(^1S_0)}{A}17MeV`$. The inclusion of the next-to-leading order corrections gives rise to the value $`\frac{U(^1S_0)}{A}13.1MeV`$. Similar calculations done in the triplet s-wave channel give rise to the value $`\frac{U(^3S_1)}{A}17.3(13.2)`$ MeV in LO (NLO). The values of the potential energy obtained with chiral approach looks quite reasonable although they are somewhat smaller than the standard values usually obtained in the calculations with the phenomenological two-body forces . Thus, one can conclude that there is still a room for both pionic effects and many body correlations. One notes, however, that adding pions will average to NLO effects. It agrees with the results of Refs. where a good fit of nuclear properties was obtained in the framework of the effective Lagrangians with the point-like interactions. Many body forces are also expected to result in rather small corrections. It follows both from the nuclear phenomenology and from Weinberg counting rules . So it is reasonable to expect that the inclusion of pion effects and many-body interactions could change the exact value of the LEC’s by some factor of order unity keeping their order of magnitude the same. It is interesting to see whether the values of LEC’s obtained above satisfies the naturalness criteria elaborated for the nuclear matter case in . According to the concept of naturalness as formulated in an individual term in the effective Lagrangians can schematically be written as
$$c\left[\frac{\psi ^+\psi }{f_\pi ^2\mathrm{\Lambda }}\right]^l\left[\frac{}{\mathrm{\Lambda }}\right]^n(f_\pi \mathrm{\Lambda })^2.$$
(11)
Applying the scaling rules developed in to extract all dimensionful factors and assuming that $`\mathrm{\Lambda }`$ 600 MeV one finds $`c_0(c_2)0.7(0.65)`$. Thus the dimensionless coefficients are indeed compatible with naturalness. As was pointed out in this is a nontrivial fact as the terms in the effective Lagrangian are supposed to absorb long distance effects from the ladder and ring diagrams.
The validity of the EFT description is restricted by some cutoff parameter $`\mathrm{\Lambda }`$, following from naturalness and reflecting the short range physics effects. Its value deserves some comments in the context of applying of the EFT methods to nuclear matter. The scale where the EFT treatment ceases to be valid should approximately correspond to the scale of the short range correlations (SRC), that is, $`500600MeV`$. The description of SRC is hardly possible in the framework of EFT so the value $`\mathrm{\Lambda }500600MeV`$ might put natural constraint on the EFT description of nuclear matter. To make the chiral expansion meaningful the chiral counting rules in nuclear matter must be established. This is still open problem. However, the above obtained results suggest that the relevant expansion parameter could be something like $`\frac{<p>}{\mathrm{\Lambda }}\frac{m_\pi }{\mathrm{\Lambda }}0.30.4`$, where $`<p>`$ is the nucleon average momentum in nuclear matter. Of course, until pion effects and many-body interactions are taken into account this estimate can only be suggestive. |
warning/0002/hep-th0002244.html | ar5iv | text | # 1 Opening Remarks
## 1 Opening Remarks
In studying brane configurations related to large $`N`$ $`SU(N)`$ pure gauge theories with eight supercharges, the authors of ref. considered the BPS supergravity solutions which ought to result in taking $`gN`$ large, where $`g`$ is the string coupling, and $`N`$ is the number of constituent branes. Such supergravity solutions are afflicted by a naked singularity known as a “repulson” which is unphysical, and incompatible with the physics of the gauge theory.
Upon closer examination (by investigating how such a geometry could have arisen by constructing it out of a large number of its constituent BPS parts) it was argued that the repulson is not present. The supergravity solution may only be taken as physical down to a radius of closest approach. At that locus of points (in the case of a $`(p+1)`$–dimensional gauge theory, it is a $`(4p)`$–sphere, $`S^{4p}`$), there is an enhanced gauge symmetry in the parent string theory and new physics, consistent with the related $`SU(N)`$ gauge theory, takes over.
That locus of points —called the “enhançon”— is new type of hypersurface essentially made of D–branes. The entire curved geometry is produced by a large number of identical BPS objects. An individual unit, when separated from all the others, is an object which has a simple description in terms of (roughly) a pair of D–branes, one of which is partially wrapped on a $`K3`$ surface, and the other which is induced by the wrapping (see later). It it therefore a sharply localised and heavy object. Upon approaching the geometry produced by a large number of its counterparts, the unit becomes lighter and less sharply defined, ultimately going to zero mass while spreading out completely at the enhançon locus.
Ref. also went on to display a variety of familiar dual situations in string theory (with related gauge theory physics) in which the enhançon phenomena described above play a crucial role. For this reason, and also because it is a genuinely new mechanism by which string theory avoids an important class of spacetime singularity, the enhançon deserves to be better understood and characterised.
We show in this paper that enhançons also arise naturally in similar situations pertaining to large $`N`$ $`SO(2N+1)`$, $`USp(2N)`$ and $`SO(2N)`$ $`(p+1)`$–dimensional gauge theories, and construct these new classes using orientifolds. It is therefore clear that enhançons may be broadly classified into types $`A`$, $`B`$, $`C`$, and $`D`$. (There is no natural $`E`$–type which has a smooth geometrical interpretation, since the rank of those groups cannot be made arbitrarily large in order to make contact with a supergravity discussion.) The latter three types differ globally from type $`A`$ by having an extra $`𝐙_2`$ identification, making them into $`\mathrm{I}\mathrm{RP}^{4p}S^{4p}/𝐙_2`$ instead of $`S^{4p}`$. All types can be distinguished locally by examining their subleading behaviour in $`N`$.
As a concrete example, we shall focus in particular on $`2+1`$ dimensional gauge theory for $`N`$ large. The reason that we focus on this case is that we can explicitly write the relevant part of the supergravity solution, using the fact that one of the relevant orientifold 6–planes has a smooth M–theory realization as the Atiyah–Hitchin manifold, while D6–branes are related to Taub–NUT. While much of the structure of the final result can be deduced on general grounds (the overall global $`𝐙_2`$ is the main feature), we show that a whole family of $`1/N`$ corrections can be completely characterised using the construction that we present.
An overview is as follows: In the next section, we orient the reader and set up our notation by reviewing the salient features of ref.. In section 3, we discuss generally how orientifolds yield the other types of enhançon. Crucially, we use a familiar gauge theory fact to help us make a general statement about the result of wrapping branes and orientifolds on $`K3`$.
In section 4, in preparation for the device of using the Atiyah–Hitchin manifold to construct the eleven dimensional solution for the $`2+1`$ dimensional $`B,C,D`$ gauge theory cases, we lift the $`A`$–type case to M–theory and reconsider it in M–theory terms. In particular, we observe that while in ten dimensions there is a discussion of the geometry in terms of its constituents being dynamical objects (D–branes), there is no analogous discussion involving dynamical objects in M–theory. The geometry can be discussed only in terms of the non–dynamical $`K3`$ and multi–Taub–NUT. Taking a probe to be an M2–brane fails to show the enhançon, since the geometry cannot be constructed out of them. Fortunately, long before one reaches the “repulson” singularity in the geometry, there is a sensible dual heterotic string description (one of the duals in ref. (see also refs.)), where the geometry again has natural brane probes revealing the enhançon. So we see again that the enhançon mechanism resolves the physics of a supergravity situation, this time by driving eleven dimensional supergravity back to string theory.
In section 5, we show how to modify the discussion to make it pertinent to the orientifolded enhançon, using the Atiyah–Hitchin manifold combined with multi–Taub–NUT. On returning to string theory, we study the supergravity solution, and extract the expression for the enhançon radius and the leading $`1/N`$ correction which follows from the orientifold’s presence. We point out that an entire class of $`1/N`$ corrections can be concisely summarised in terms of the exponentially small differences between the smooth Atiyah–Hitchin manifold and the (negative mass) Taub–NUT solution. Unfortunately, we do not have such control over all of the corrections present in the geometry.
We also present the result of a probe computation analogous to that performed in ref. which yields the one–loop result for the metric on a subspace of the Coulomb branch of moduli space for the $`B,C,D`$ gauge theories. Again, they differ from the $`SU(N)`$ case by a global $`𝐙_2`$ action, and the leading behaviour for the $`1/N`$ corrections we computed. In all cases, there is a related monopole moduli space problem, in the spirit of refs..
## 2 The $`A`$–Type Enhançon
Consider wrapping $`N`$ coincident D6–branes on a $`K3`$ surface of volume $`V`$. This results in an effective 2+1 dimensional object, with $`N`$ units of D2–brane charge, due to the interaction
$$\frac{\mu _6}{48}C_{(3)}p_1()$$
(2.1)
on the D6–brane world–volume. The precise value $`N`$ comes about since $`\mu _6=(2\pi )^6\alpha ^{7/2}`$, $`\mu _2=(2\pi )^2\alpha ^{3/2}`$, $`=4\pi ^2\alpha ^{}R`$, and because for $`K3`$
$$p_1(R)\frac{1}{8\pi ^2}RR,$$
(2.2)
integrates to $`48`$. We will call this wrong–sign D2–brane a “D2\*–brane”. It preserves the same supersymmetries as a correct sign D2–brane with the same orientation, and therefore is not an anti–D2–brane. It is useful to think of it as a brane which is bound inside the D6–brane worldvolume, resulting from the curvature of the $`K3`$. It is quite analogous to the (correct sign) D2–brane which would be bound inside the worldvolume of a D6–brane if there was a field theory instanton configuration, due to the term
$$\frac{\mu _6}{2}C_{(3)},$$
(2.3)
where $`=2\pi \alpha ^{}F`$. An instanton in the 6+1 dimensional gauge theory has $`(8\pi ^2)^1FF=1`$, and consequently has the charge of a single D2–brane. In the limit where the instanton shrinks to zero size, there is a good description in the full string theory corresponding to a fully localised pointlike D2–brane.
Similarly, one would recover pointlike D2\*–branes from wrapping the D6–branes if $`K3`$’s curvature was located at a finite number of points, such as in an orbifold limit. This situation (no doubt) has a good string theory description, and is worth investigating. In the present case, the curvature of the $`K3`$ is distributed everywhere, and correspondingly the D2\*–branes are delocalised everywhere on it.
Imagine that the $`K3`$ surface lies in the $`x^6,x^7,x^8,x^9`$ directions, and that the remaining (unwrapped) part of the D6–brane lies in the $`x^0,x^1,x^2`$ directions. There is an $`SU(N)`$ gauge theory on the $`2+1`$ dimensional worldvolume, with eight supercharges. The gauge supermultiplet consists of a gauge field $`A_\mu `$ and three scalars $`\varphi _i`$, where $`i=3,4,5`$. The scalars parameterise the positions of the D6–D2\* system in the transverse directions, $`x^3,x^4,x^5`$. This vector supermultiplet transforms in the adjoint representation of $`SU(N)`$. The gauge theory has a scalar potential of the form $`\mathrm{Tr}[\varphi _i,\varphi _j]^2`$. Supersymmetric solutions of the theory, giving a moduli space of vacua, may be found by choosing vacuum expectation values (“vevs”) of the scalars such that they are in the Cartan subalgebra of $`SU(N)`$. This breaks $`SU(N)U(1)^{N1}`$, giving the “Coulomb branch” of moduli space.
Classically, the moduli space is
$$_{\mathrm{cl}}^N=\frac{\left(\mathrm{I}\mathrm{R}^3\times S^1\right)}{S_{N1}}^{N1},$$
(2.4)
where the $`S^1`$ factors represent the periodic scalars resulting from dualising the gauge fields (recall that we are in $`2+1`$ dimensions). The $`S_{N1}`$ is the Weyl group of $`SU(N)`$, which acts as permutations of the $`N1`$ eigenvalues of the $`\varphi `$’s, which are now in the Cartan subalgebra. $`U(1)^{N1}`$ is the gauge symmetry on $`N`$ separated, but wrapped D–branes, where the extra $`U(1)`$ we would naively expect corresponds to the overall centre of mass of the system.
We will focus on the situation where all of the branes are coincident, which is to say that the vev’s of all of the fields are given the same value, except for a complete set of four making a multiplet giving the location of a probe brane in the background of all the others. In the gauge theory, this is equivalent to focusing on a particular subspace of the relative moduli space. In another, equivalent description, it is the four dimensional $`(1,N1)`$ subspace representing relative moduli space of the full moduli space of $`N`$ $`SU(2)`$ monopoles; $`N1`$ of them are coincident, and one is separated. The classical moduli space is then
$$_{\mathrm{cl}}^{(1,N1)}=\mathrm{I}\mathrm{R}^3\times S^1.$$
(2.5)
One of the results of ref. (see below) is the computation of the one–loop result for the metric on this moduli space. Here, we shall compute a closely related version, representing a similar subspace of the Coulomb branch of the $`SO(2N+1)`$, $`USp(2N)`$, or $`SO(2N)`$ gauge theory. These will also have an interpretation as multimonopole moduli spaces, where the monopoles are of an $`SU(2)`$ gauge theory with a $`𝐙_2`$ identification.
For $`gN`$ large, ($`g`$ is the closed string coupling) we have a chance of obtaining a good description of the geometry of the system in terms of a ten dimensional type IIA supergravity solution, which is
$`ds^2`$ $`=`$ $`Z_2^{\frac{1}{2}}Z_6^{\frac{1}{2}}\eta _{\mu \nu }dx^\mu dx^\nu +Z_2^{\frac{1}{2}}Z_6^{\frac{1}{2}}dx^idx^i+V^{\frac{1}{2}}Z_2^{\frac{1}{2}}Z_6^{\frac{1}{2}}ds_{\mathrm{K3}}^2,`$
$`e^{2\mathrm{\Phi }}`$ $`=`$ $`g^2Z_{2}^{}{}_{}{}^{\frac{1}{2}}Z_{6}^{}{}_{}{}^{\frac{3}{2}},`$
$`C_{(\mathit{3})}`$ $`=`$ $`g^1(Z_2^11)dx^0dx^1dx^2,`$
$`C_{(\mathit{7})}`$ $`=`$ $`g^1(Z_6^11)dx^0dx^1dx^2dx^6dx^7dx^8dx^9.`$ (2.6)
Here, $`\mu ,\nu =0,1,2`$; $`i=3,4,5`$ and the $`x^6,x^7,x^8,x^9`$ directions contain $`ds_{\mathrm{K3}}^2`$, the (unknown) metric of a unit volume $`K3`$. The 345–harmonic functions representing the D2\*– and D6–branes respectively are:
$$Z_2=1+\frac{r_2}{r}\mathrm{and}Z_6=1+\frac{r_6}{r},$$
(2.7)
(recall that the D2\*’s are delocalised in $`K3`$), with
$$r=|\overline{r}|,\overline{r}\mathrm{I}\mathrm{R}_{3,4,5}^3,r_2=\frac{(2\pi )^4gN\alpha ^{5/2}}{2V}\mathrm{and}r_6=\frac{gN\alpha ^{1/2}}{2}.$$
(2.8)
The latter are written so as to give the masses of the BPS object which is formed when we wrap a D6–brane to make the D6–D2\* object:
$$\tau =\frac{N}{g}(\mu _6V\mu _2)=\frac{N}{g}\mu _6(VV_{})=\frac{N}{g}\mu _2\left(\frac{V}{V_{}}1\right),$$
(2.9)
where $`V_{}=(2\pi \sqrt{\alpha ^{}})^4`$.
There are a number of things to note about this supergravity solution. First, note that $`g`$ appears as the asymptotic value of the string coupling far away from the core of the solution ($`r\mathrm{}`$). The actual string coupling in the interior of the solution is given by the value of $`e^\mathrm{\Phi }`$, as usual, and varies with $`r`$. Similarly, the volume of $`K3`$ is a function of $`r`$: $`V(r)=VZ_2(r)/Z_6(r)`$, which approaches $`V`$ asymptotically, and decreases, becoming zero at the singularity $`r=|r_2|`$.
One of the key points noticed in ref. is that while a naive examination of the supergravity solution shows an unsettling naked singularity (the “repulson”) at $`r=|r_2|`$, this part of the geometry is actually non–physical. The geometry should only be taken at face value down to radius
$$r_\mathrm{e}=\frac{2V}{VV_{}}|r_2|.$$
(2.10)
This is the radius at which a number of special things happen:
* The volume of $`K3`$ is equal to the special value $`V_{}=(2\pi \sqrt{\alpha ^{}})^4`$.
* The 5+1 dimensional $`K3`$–compactified string theory has an R–R sector $`U(1)`$ which becomes enhanced to an $`SU(2)`$ gauge symmetry.
* A D6–D2\* probe is a monopole of this $`U(1)`$, and becomes massless at the enhanced symmetry point. It also ceases to be pointlike, and dissolves into the “enhançon” locus at $`r_\mathrm{e}`$, which is an $`S^2`$.
The interpretation of these and other facts uncovered in ref. is that there are no brane sources for $`r<r_\mathrm{e}`$, and therefore the supergravity solution inside that radius is simply the trivial flat solution with no R–R fields switched on. The smooth interpolating region between the two solutions in the neighbourhood of the enhançon radius is described by the relatively innocuous (but nonetheless interacting) monopole physics.
On the one hand, this situation represents another remarkable method by which string theory rids itself of potentially troublesome singularities, while on the other hand, it potentially teaches us something about gauge theories. For example, in seeking for a limit in which the gauge theory decouples from the rest of bulk physics, we take $`\alpha ^{}0`$ while holding fixed the 2+1 dimensional gauge coupling given by:
$$g_{\mathrm{YM}}^2=(2\pi )^4g\alpha ^{3/2}V^4$$
(2.11)
and hold fixed $`U=r/\alpha ^{}`$. In this limit, it was found that the metric on the moduli space, as seen by the D6–D2\* probe, can be read off from the effective Lagrangian for the monopole probe moving in the transverse space with coordinates $`(U,\theta ,\varphi ,\sigma )`$:
$$=f(U)\left(\dot{U}^2+U^2\dot{\mathrm{\Omega }}_2^2\right)+f(U)^1\left(\dot{\sigma }\frac{(N1)}{8\pi ^2}A_\varphi \dot{\varphi }\right)^2,$$
(2.12)
where<sup>1</sup><sup>1</sup>1Note that we have inserted $`N1`$ instead of the $`N`$ which appears in the supergravity solution (2.6) and also in the probe result exhibited in ref.. Strictly speaking, there are $`N1`$ D6–D2\* units being probed by one separated unit, giving $`N`$ in total. The difference is a $`1/N`$ effect, not considered in ref., but should be included here since we will later be discussing a family of corrections at that order.
$$f(U)=\frac{1}{8\pi ^2g_{\mathrm{YM}}^2}\left(1\frac{g_{\mathrm{YM}}^2(N1)}{U}\right),\dot{\mathrm{\Omega }}_2^2=\dot{\theta }^2+\mathrm{sin}^2\theta \dot{\varphi }^2,$$
(2.13)
with $`0\theta <0,\mathrm{\hspace{0.17em}0}\varphi <2\pi `$ and $`U_\mathrm{e}<U<\mathrm{}`$. Here, $`U_\mathrm{e}=g_{\mathrm{YM}}^2(N1)`$ and $`A_\varphi =\pm 1\mathrm{cos}\theta `$ is a $`U(1)`$ monopole potential. The metric in (2.12) is the Euclidean Taub–NUT metric, with a negative mass. It is a hyperKähler manifold, because $`f=\times A`$, where $`A=((N1)/8\pi ^2)A_\varphi d\varphi `$. The coordinate $`\sigma `$ is periodic with period $`4\pi `$, and is the dual of the $`U(1)`$ centre–of–mass gauge field on the 2+1 dimensional worldvolume of the monopole probe.
This result is completely in accord with the expectation from gauge theory, being the one–loop result for the metric on moduli space, in the special case where the $`N1`$ coordinates parameterising the Cartan subalgebra are chosen to be equal, corresponding to making all of the branes coincident. The enhançon is at $`U=U_\mathrm{e}`$, and corresponds to the Landau pole, representing in gauge theory the place where the one–loop correction makes the gauge coupling diverge.
In the equivalent monopole language, this is (an approximation to) the metric on a subspace $`^{(1,N1)}`$ (described above eqn. (2.5)) of the full moduli space of $`N`$ $`SU(2)`$ monopoles. This $`SU(2)`$ is the enhanced gauge symmetry from whence comes the name “enhançon”.
There are exponential corrections to this metric which will remove the singular behaviour and complete it into a smooth hyperKähler manifold, $`^{(1,N1)}`$. This space generalises the Atiyah–Hitchin manifold, which is the metric on the relative two–monopole moduli space $`^{(1,1)}`$ which governs the case of $`SU(2)`$ gauge theory.
A natural question arises about the nature of the story for the case where one studies gauge groups other than $`SU(N)`$. It is straightforward to construct gauge groups $`SO(2N)`$, $`SO(2N+1)`$, and $`USp(2N)`$ in perturbative string theory by combining D–branes with orientifolds. Studying the wrapping of such a system on $`K3`$ should therefore be our first step in answering the question. Let us do that.
## 3 Including Orientifolds
On general grounds, one expects a similar story to that which was constructed above, as all of the constituent features which are present to make the physics work as it should are still present after we insert an orientifold six–plane (O6–plane) parallel with the D6–branes. Of course, the details of precisely where the enhançon is located (corresponding to where in $`\mathrm{I}\mathrm{R}_{345}^3`$ the $`K3`$ volume reaches the value $`V_{}`$) will be modified, but only at subleading order in $`N`$. Globally, the orientifold will also place a $`𝐙_2`$ identification on $`\mathrm{I}\mathrm{R}_{345}^3`$ ($`𝐙_2`$ acts by multiplying each of $`x^3,x^4,x^5`$ by $`1`$), turning the $`S^2`$ of the enhançon into $`\mathrm{I}\mathrm{RP}^2S^2/𝐙_2`$. The basic problem is to understand the nature of the supergravity solution in the presence of the orientifolds, which we will do below in a particular case.
First, let us understand the physics of the perturbative string theory description, containing the weakly coupled gauge theory. For small $`gN`$, we have $`N`$ D6–branes, and an O6–plane parallel to them. This gives a gauge group $`SO(2N+1)`$, $`USp(2N)`$ or $`SO(2N)`$. In the latter case, the O6–plane has negative charge, equal to $`2\mu _6`$ and is often denoted O6<sup>-</sup>. We can obtain the former case by trapping a half D6–brane on the O6<sup>-</sup>–plane: this combination is often referred to as an $`\stackrel{~}{\mathrm{O6}}`$–plane, with charge $`3\mu _6/2`$. In the middle case, the O6–plane has charge $`+2\mu _6`$ and is written O6<sup>+</sup>. To be concise, we will use the symbol $`\alpha `$ to denote these O6–charges, measured in D2–brane units. It takes the values $`\alpha =3/2,+2,2`$, respectively.
We now wrap the whole system on $`K3`$. This results in the induced D2\*–branes as described above, but with an additional contribution. This is due to a curvature coupling, this time on the world–volume of the O6–plane, similar to eqn.(2.1). The couplings are different in each case $`\stackrel{~}{\mathrm{O6}}`$, O6<sup>+</sup>, O6<sup>-</sup>:
$$\frac{\mu _6}{32}C_3p_1(),\frac{5\mu _6}{48}C_3p_1(),\frac{\mu _6}{48}C_3p_1(),$$
(3.14)
which, after wrapping on $`K3`$ will induce some $`C_{(3)}`$ charge, $`\beta `$, which in D2–brane units, is respectively $`\beta =3/2,5,1`$. This will modify the contribution to the effective amount of D2\*–brane present<sup>2</sup><sup>2</sup>2The temptation to interpret these extra charges as induced wrong sign O2–planes should, we believe, be firmly resisted. First of all, the resulting charges are hard to interpret, given the existing types of O2–brane already known. Secondly, one would have to insert a $`𝐙_2`$ identification on the $`K3`$ part of the spacetime, which is hard to justify as the result of a smooth wrapping process. The most economical interpretation is the one presented here..
Note that we can introduce extra (correct sign) D2–branes parallel to the D2\*–branes into the story while preserving the eight supercharges. Open strings stretching between these new branes and the wrapped system play the role of extra hypermultiplets in the 2+1 dimensional gauge theory. For $`M`$ D2–branes, we have $`M`$ species of such hypermultiplets. For consistency, in the presence of the orientifold, these D2–branes will have the opposite orientifold projection acting on them from that acting on the D6’s, as follows from T–duality to the situations studied in refs..
So, for example, while there is an $`SO(2N)`$ (or $`USp(2N)`$) gauge symmetry on the D6–branes, the $`M`$–flavour sector has a $`USp(2M)`$ (or $`SO(2M)`$) symmetry. This is indeed correct from the perspective of gauge theory, and this fact has featured in the physics of orientifolds before. In ref., it was shown to correspond to the phenomenon that the orientifold must change its sign when it passes through an NS5–brane. Actually, one of the dual realizations of the enhançon story involves NS5–branes. We display it with the inclusion of the orientifold in figure 1, where our case here is $`p=2`$. The D3–branes in the interior are dual to the D6–D2\* units, while the those on the exterior (supplying matter multiplets) are dual to the ordinary D2–branes. The orientifold runs through the whole system, having a minus sign on the interior (giving $`SO(2N)`$) and a plus sign on the exterior, giving $`USp(2M)`$.
So we see that the sign flip of the O3–plane on either side of the NS5–brane is dual to the fact that the O6–plane has the opposite projection on the D2’s from that on the D6’s. Of course, this discussion clearly generalises to all D$`p`$–branes and D$`p`$\*–branes with orientifolds, and the dual situation involving D$`(p+1)`$–branes stretched between NS5–branes with an O$`(p+1)`$–plane passing through.
In this way, we see that we can consistently construct wrapped D–brane and orientifold systems giving gauge groups $`SO(2N+1)`$, $`USp(2N)`$, and $`SO(2N)`$. The amount of D2\*–branes induced from the wrapping is modified from $`N`$ (for the $`SU(N)`$ case) to $`N3/2`$, $`N5`$ and $`N1`$, respectively. The changes are to be thought of as $`1/N`$ corrections to the original case, and are different for each type of orientifold. We can now consider taking $`gN`$ large, and expect that the phenomena which occurred for $`SU(N)`$ will happen again, giving an enhançon for each case. We shall name the types of enhançon which can occur in each situation the $`A`$–type (for $`SU(N)`$), $`B`$–type (for $`SO(2N+1)`$), $`C`$–type (for $`USp(2N)`$) and $`D`$–type ($`SO(2N)`$).
Of course, there is no natural definition of an $`E`$–type, for (at least) two reasons: There is no known perturbative way to make $`E_{6,7,8}`$ gauge symmetry with D–branes, and furthermore, the enhançon as a smooth geometric object is a large $`N`$ phenomenon, which is incompatible with the fact that the maximum rank of the exceptional groups is eight.
The next step in our story will be to write down the geometry corresponding to the large $`gN`$ physics of the wrapped system of D–branes and orientifolds. The observation that we shall use to achieve this is the fact that for the O6<sup>-</sup> case (giving 2+1 dimensional $`SO(2N)`$), the supergravity geometry of the system can be written down accurately enough for us to make progress. Along the way, we will see that we can study cases $`SO(2N+1)`$ and $`USp(2N)`$ accurately enough for our purposes using similar techniques.
## 4 Uplifting the Enhançon
Before we proceed to the new types, let us pause for a moment to consider the $`A`$–type enhançon story in eleven dimensional terms. Recall that the metric of the Taub–NUT space, made into an eleven dimensional supergravity solution (by adding $`\mathrm{I}\mathrm{R}^{6,1}`$ for the world–volume directions) is (defining an eleventh direction $`\psi =x^{\mathrm{}}/16m`$):
$`ds_{11}^2=dt^2+dx_1^2+dx_2^2+dx_6^2+dx_7^2+dx_8^2+dx_9^2`$ (4.15)
$`+F(r)(dr^2+r^2d\mathrm{\Omega }_2^2)+F^1(r)\left(d\psi +C_\varphi d\varphi \right)^2,`$
where, with $`0\psi <4\pi `$,
$$d\mathrm{\Omega }_2^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2,F=1+\frac{4mN}{r},C_\varphi =4mN\mathrm{cos}\theta .$$
(4.16)
Reducing along the $`\psi `$–circle, the relation between eleven dimensional metrics and ten dimensional type IIA fields is:
$$ds_{11}^2=e^{\frac{2}{3}\varphi }ds_{10}^2+e^{\frac{4}{3}\varphi }(d\psi +C_{(1)})^2,$$
(4.17)
and so we recover the now standard fact that Taub–NUT corresponds to a familiar ten–dimensional solution:
$`ds_{10}^2=Z_6^{\frac{1}{2}}(dt^2+dx_1^2+dx_2^2+dx_6^2+dx_7^2+dx_8^2+dx_9^2)`$
$`+Z_6^{\frac{1}{2}}(dr^2+r^2d\mathrm{\Omega }_2^2)`$
$`Z_6=F(r);e^\varphi =Z_6^{\frac{3}{4}}(r),C_\varphi =4mN\mathrm{cos}\theta ,`$ (4.18)
which is precisely the D6–brane solution, if we identify $`4mN=r_6`$ (and set the asymptotic value of the dilaton to $`\mathrm{log}g`$.) The one–form potential $`C_{(1)}=C_\varphi d\varphi `$ can be Hodge–dualised in ten dimensions to give an electric source for $`C_{(7)}`$ of precisely the form given in eqn. (2.6).
Turning to the enhançon, by using the prescription of eqn. (4.17), supplemented with a direct uplift of the three–form potential $`C_{(3)}`$ to give the components of the eleven dimensional three–form $`A_{(3)}`$, it is easy to write an eleven dimensional solution for the uplifted D6–D2\* system:
$`ds_{11}^2=\stackrel{~}{Z}_2^{\frac{2}{3}}(dt^2+dx_1^2+dx_2^2)+\stackrel{~}{Z}_2^{\frac{1}{3}}\stackrel{~}{V}^{\frac{1}{2}}ds_{\mathrm{K3}}^2`$
$`+\stackrel{~}{Z}_2^{\frac{1}{3}}\left[\stackrel{~}{Z}_6(dr^2+r^2d\mathrm{\Omega }_2^2)+\stackrel{~}{Z}_6^1\left(d\psi +C_\varphi d\varphi \right)^2\right],`$
$`\mathrm{with}`$ $`A_{(3)}=\left(\stackrel{~}{Z}_2^11\right)dx^0dx^1dx^2,`$ (4.19)
$$\mathrm{where}\stackrel{~}{V}=g^2V,\stackrel{~}{Z}_2=gZ_2,\stackrel{~}{Z}_6=g^1Z_6.$$
(4.20)
It is interesting to contrast the interpretation of this solution with the ten dimensional discussion. Recall that from the point of view of ten dimensions, there is the geometry of $`K3`$, accompanied by D6–branes wrapped on it. The wrapping induced some D2\*–branes, completely delocalised in the $`K3`$. We were able to probe the geometry of the supergravity solution (4.19) with one of its basic constituents, a single D6–D2\* BPS object.
From the point of view of the eleven dimensional supergravity solution, everything is geometry: there are no branes here. The Taub–NUT part lies in the 345$`\mathrm{}`$ directions, while $`K3`$ lies in the 6789 directions. Together, they act as a source for the three–form potential $`A_{(3)}`$, due to the supergravity term:
$$A_{(3)}X_8,\mathrm{with}X_8=\frac{1}{24}\left(p_2\frac{1}{4}p_1p_1\right).$$
(4.21)
Given that for Taub–NUT of charge $`N`$, $`p_1=2N`$ and, as stated before, for $`K3`$ we know that $`p_1=48`$, we get $`N`$ units of $`A_{(3)}`$ charge, as the solution shows. This fits, as (4.21) is the M–theory ancestor of the brane world–volume term (2.1).
Sadly, there is no natural extended dynamical object which we can use as a candidate for the basic constituent of the geometry. Thus, we cannot perform a world–volume probe computation to deduce the true geometry. It is tempting to read the $`\stackrel{~}{Z}_2`$ part of the geometry as representing a “wrong sign” M2–brane which is otherwise dynamical, (perhaps restoring a $`1/r^3`$ behaviour to make it also localised in the $`x^{\mathrm{}}`$ direction.) Unfortunately, this cannot work. The putative M2\*–brane necessarily would have negative tension at all locations in $`\mathrm{I}\mathrm{R}_{345}^3`$, and since there is no larger wrapped brane with positive tension to combine it with to made a positive tension object, we cannot write a sensible worldvolume action. Of course, a probe computation with a correct sign M2–brane (which preserves the same amount of supersymmetry) gives a sensible result: simply the pure (with mass parameter of $`+4mN=r_6`$) Taub–NUT metric, as it should, with no sign of either enhançon or repulson. This is in accord with our expectation that the repulson (still present at $`r=|r_2|`$) is an artifact, while the enhançon should be invisible to an M2–brane because its world–volume theory does not relate to the $`SU(N)`$ gauge theory.
To get at the enhançon, there is a pertinent supergravity question to be asked all the same: Can we envision a supergravity mechanism by which the troublesome repulson singularity at $`r=|r_2|`$ is avoided? In the string theory situation, we saw that the $`K3`$ reached the natural value $`V_{}=(2\pi \sqrt{\alpha ^{}})^4`$, before it reached its singular value (zero), and the physics of the enhanced gauge symmetry took over. Is there a special value for the volume of $`K3`$ in this case? Here, the natural length scale is of course set by the Planck length, $`\mathrm{}_{11}=g^{1/3}\mathrm{}_s`$, which the system again reaches before the singular value of zero. Once $`K3`$ has shrunk to that size, we should search for a better description than eleven dimensional supergravity. The alternative to using the full (unknown) M–theory is to search for a dual description. Happily, eleven dimensional supergravity on such a small $`K3`$ is well described by the heterotic string on $`T^3`$. The distinguished $`\psi `$–circle which is fibred to make the Taub–NUT joins the rest to become a $`T^4`$, and the Taub–NUT structure becomes a “warped” (not “wrapped”) NS–fivebrane/Kaluza–Klein monopole structure giving rise to a monopole membrane whose mass goes to zero at an $`SU(2)`$ enhanced point of the torus (see also ref.). So we see that again, stringy physics (now heterotic) takes over before we get to the repulson radius, and repairs the geometry with the same $`SU(2)`$ physics that we saw in type IIA string theory.
## 5 The Orientifolded Enhançon
Just as the D6–brane has an interpretation as the Taub–NUT spacetime upon going to low energy M–theory (eleven dimensional supergravity), in a similar fashion, the O6<sup>-</sup>–plane becomes the Atiyah–Hitchin manifold, described by a metric:
$`ds_{11}^2=dt^2+dx_1^2+dx_2^2+dx_6^2+dx_7^2+dx_8^2+dx_9^2`$ (5.22)
$`+f(\rho )dr^2+8m^2\left(a^2(\rho )\sigma _1^2+b^2(\rho )\sigma _2^2+c^2(\rho )\sigma _3^2\right),`$
where
$`\sigma _1=\mathrm{sin}\psi d\theta +\mathrm{cos}\psi \mathrm{sin}\theta d\varphi ,\sigma _2=\mathrm{cos}\psi d\theta +\mathrm{sin}\psi \mathrm{sin}\theta d\varphi `$
$`\sigma _3=d\psi +\mathrm{cos}\theta d\varphi ,\rho ={\displaystyle \frac{r}{8m}},\psi ={\displaystyle \frac{x^{\mathrm{}}}{16m}},`$ (5.23)
and the functions $`f,a,b,c,d`$ are given in terms of elliptic integrals in ref.. There is also an identification by:
$$(r,\theta ,\varphi ,\psi )(r,\pi \theta ,\pi +\varphi ,\psi ),$$
(5.24)
which, in terms of the coordinates $`(x^3,x^4,x^5)`$ and the M–direction $`x^{\mathrm{}}`$, is simply a multiplication by a minus sign on all directions. The displayed metric (5.22) has a conical singularity at $`r=8\pi m`$. The space made by imposing the $`𝐙_2`$ identification (5.24) is the Atiyah–Hitchin space, and it is free of conical singularities.
While a closed form for the metric cannot be written, for large $`r`$ the metric becomes
$`ds_{11}^2=dt^2+dx_1^2+dx_2^2+dx_6^2+dx_7^2+dx_8^2+dx_9^2+`$
$`+G(r)(dr^2+r^2d\mathrm{\Omega }_2^2)+G^1(r)\left(d\psi +C_\varphi d\varphi \right)^2,`$ (5.25)
where
$$G=1\frac{16m}{r},C_\varphi =16m\mathrm{cos}\theta ,$$
(5.26)
which we recognise as the metric for Taub–NUT, but with a negative mass. Clearly, it can be reduced to ten dimensions in the same way as before, and we see that it has $`2`$ units of D6–brane charge, which is in accord with our knowledge of the charge of an O6<sup>-</sup>–plane from perturbative string theory. (The actual appearance of $`16m`$ in the metric instead of $`8m`$ follows from the fact that the displayed metric is the double cover of the actual solution: recall that we must divide by the $`𝐙_2`$ action.)
Now we are in a position to construct the geometry which gives rise the the $`D`$–type enhançon. We simply combine the geometry of the Atiyah–Hitchin manifold with that of $`N`$ coincident–centred Taub–NUT. The exact smooth metric certainly exists (the $`N=1`$ case is known, and is Dancer’s manifold), but we need not be able to write it exactly to get at the physics we require. The radius at which the enhançon appears can be tuned to be arbitrarily large by making $`N`$ as large as we like, so we can rest assured that if we take the approximate expression for the Atiyah–Hitchin manifold, we can capture the essential physics for large $`N`$.
Once we relax the condition of exactness, and focus on the large $`r`$ part of the solution, we can include the cases of the $`B`$– and $`C`$–type enhançons. While a precise relation to a cousin of the smooth Atiyah–Hitchin$`+`$Taub–NUT geometry is not known, at large $`r`$, the difference is immaterial, as only the leading behaviour is needed to characterise the enhançon at large enough $`N`$. We can simply use the same supergravity solution as before, but with different numbers inserted into the $`1/N`$ corrections to the harmonic functions.
It is clear therefore, that for all cases our solution can be written (for large enough $`r`$) in the covering space in the precise form of eqn. (4.19), but with the replacement of $`\stackrel{~}{Z}_2`$ and $`\stackrel{~}{Z}_6`$ by (respectively):
$$\stackrel{~}{Z}_2^{}=g\left(1\frac{2|r_2|(1\beta /N)}{r}\right)\mathrm{and}\stackrel{~}{Z}_6^{}=\frac{1}{g}\left(1+\frac{2r_6(1+\alpha /N)}{r}\right).$$
(5.27)
Here<sup>3</sup><sup>3</sup>3It is amusing to note that the sum $`\alpha +\beta `$ is the same in each case. We do not know if this has any physical significance. Later, in eqn. (5.31), we shall see that it is $`\alpha \beta `$ which controls the leading $`1/N`$ correction to the enhançon in each case. Were it the sum which appeared, we would have had a remarkably universal result. $`\beta =3/2,5,1`$, and $`\alpha =3/2,+2,2`$ for types $`B,C,D`$, respectively. We have deduced $`\stackrel{~}{Z}_2^{}`$’s asymptotic form<sup>4</sup><sup>4</sup>4While we know (in the $`D`$–type case) precisely how the harmonic function of $`\stackrel{~}{Z}_6^{}`$ gets corrected into the smooth Atiyah–Hitchin+Taub–NUT solution, we do not know how $`\stackrel{~}{Z}_2^{}`$, which owes its presence to the $`K3`$ part of the eleven dimensional geometry, gets corrected. In its current form, it must be there in order to measure the correct mass and charge at large $`r`$, but the small $`r`$ details are unknown to us. from the fact that it must give the correct induced D2\*–brane mass and charge at large $`r`$ in the string theory limit.
This should be taken to mean the metric on the covering space of our solution, and we must divide by the $`𝐙_2`$ action in order to reconstruct the correct solution, as before. This also accounts for the factors of two we have inserted into the harmonic functions. Notice that the contribution to the harmonic functions of (what will become) the orientifolds is simply a $`1/N`$ correction to the geometry. This will turn into part of the family of $`1/N`$ corrections to the location and shape of the enhançon locus, once we return to string theory.
The final step is clear. We return to type IIA string theory by reducing on the $`\psi `$–circle, recovering a supergravity solution representing the large $`gN`$ geometry of system of wrapped D6–brane and and O6<sup>-</sup>–plane, as promised in section 2. The solution is simply the geometry (2.6) with $`Z_2`$ and $`Z_6`$ replaced by their $`1/N`$ corrected counterparts in (5.27) with the factors of $`g`$ and $`1/g`$ removed. Crucially, there is a $`𝐙_2`$ identification on the $`(x^3,x^4,x^5)`$ directions, making it globally distinct from the $`A`$–type case, in addition to the different structure of the subleading behaviour in $`N`$.
Again, in string theory, the natural object to construct this geometry out of is the D6–D2\* at large $`gN`$, now in the presence of an orientifold, and we may examine the nature of the geometry as seen by the probe by a computation precisely along the lines of ref.. The structure of the computation is almost identical to that carried out there, and we refer the reader to that work for the details. A crucial difference is that we are working on the covering space of the actual geometry, and so we should insert a mirror image of the probe at the image position obtained by reflecting through the orientifold fixed point. The result is structurally identical:
$$=F^{}(r)\left(\dot{r}^2+r^2\dot{\mathrm{\Omega }}_2^2\right)+F^{}(r)^1\left(\dot{s}/2\mu _2C_\varphi \dot{\varphi }/2\right)^2,$$
(5.28)
where now
$$F^{}(r)=\frac{1}{2g}(\mu _6VZ_2^{}\mu _2Z_6^{}),$$
(5.29)
with
$`Z_2^{}`$ $`=`$ $`\left(1{\displaystyle \frac{2|r_2|(1(\beta +1)/N)}{r}}\right)\mathrm{and}`$
$`Z_6^{}`$ $`=`$ $`\left(1+{\displaystyle \frac{2r_6(1+(\alpha 1)/N)}{r}}\right),`$ (5.30)
where we have shifted $`N`$ to $`N1`$ to represent separating off the probe (see footnote 1). Here, $`s`$ is the fourth modulus obtained by dualising the world–volume centre of mass gauge field. The location $`r_\mathrm{e}^{}`$ of the $`D`$–type enhançon can be read off as the place in $`r`$ where the mass of the probe becomes zero (equivalent to $`V(r_\mathrm{e}^{})=V_{}`$):
$$r_\mathrm{e}^{}=\frac{2V}{VV_{}}|r_2|\left(1\frac{\gamma }{N}\right),$$
(5.31)
with
$$\gamma =\left(\frac{\alpha \beta 2}{2}\right)=1,\frac{5}{2},\frac{3}{2}$$
(5.32)
in each case $`B,C,D`$. (The analogous expression for the $`A`$ case —with the $`1/N`$ correction from separating off the probe (c.f. eqn. (2.10))— has $`\gamma =1`$. Note that for case $`B`$ the effect of the O6–plane is precisely cancelled by the effect of the D2\*–brane contribution which is produces from wrapping, giving the same leading $`1/N`$ contribution as for type $`A`$.)
Correspondingly, when we take the limit where we decouple the gauge theory with $`\alpha ^{}0`$ holding $`g_{\mathrm{YM}}^2`$ fixed, we recover the prediction for the metric on the moduli space of the gauge theory at large $`N`$ (in the coincident limit):
$$=f^{}(U)\left(\dot{U}^2+U^2\dot{\mathrm{\Omega }}_2^2\right)+f^{}(U)^1\left(\dot{\sigma }\frac{N(1\gamma /N)}{8\pi ^2}A_\varphi \dot{\varphi }\right)^2,$$
(5.33)
where
$$f^{}(U)=\frac{1}{8\pi ^2g_{\mathrm{YM}}^2}\left(1\frac{g_{\mathrm{YM}}^2N}{U}\left(1\frac{\gamma }{N}\right)\right).$$
(5.34)
This is the one–loop expression for the metric on moduli space for the $`SO(2N+1)`$, $`USp(2N)`$ or $`SO(2N)`$ 2+1 dimensional gauge theory. On general grounds, the classical moduli space has the geometry
$$_{\mathrm{cl}}=\frac{\left(\mathrm{I}\mathrm{R}^3\times S^1\right)}{S_N\times 𝐙_2}^N$$
(5.35)
where there is a natural $`𝐙_2`$ action reflecting the $`N`$ eigenvalues into (minus) themselves. In the subspace where we set all the vev’s (but four) to be equal, we are reduced to
$$_{\mathrm{cl}}=\frac{\mathrm{I}\mathrm{R}^3\times S^1}{𝐙_2}$$
(5.36)
for the classical moduli space. Our metric above, with the $`𝐙_2`$ action (imposed, recall, for smoothness of the Atiyah–Hitchin manifold representing the O6<sup>-</sup>–plane), is the one–loop expression for the metric on the full moduli space.
Finally, we point out that once again, these results have a dual interpretation as an approximate result for the metric on moduli space of $`N`$ monopoles. This time, they are monopoles of a spontaneously broken $`SU(2)`$ theory which has an identification by $`𝐙_2`$, which can be understood as follows<sup>5</sup><sup>5</sup>5See ref. for comments on such theories in a closely related stringy context.: The relation between the moduli space of $`2+1`$ dimensional gauge theories and that of monopoles is readily seen in the string realization of such theories by D3–branes stretched between NS5–branes. The spontaneously broken $`SU(2)`$ lives on the world–volume of the NS5–branes. The ends of the D3–branes in the NS5–brane worldvolumes are the monopoles. We need only look at the orientifolded version of that picture, drawn in figure 1, to see the origin of the $`𝐙_2`$ action on the $`SU(2)`$ theory. By passing through the world–volume of the NS5–branes, the O3–plane places a spacetime $`𝐙_2`$ identification on the $`SU(2)`$ gauge theory.
## 6 Closing Remarks
The enhançon locus which appears in the study of spacetime geometry associated to $`SU(N)`$ $`(p+1)`$–dimensional gauge theory (at large $`N`$) with eight supercharges has three natural counterparts: Those pertaining to $`SO(2N+1)`$, $`USp(2N)`$ and $`SO(2N)`$ gauge theory. The four classes deserve to be called types $`A`$, $`B`$, $`C`$, and $`D`$. (There is no natural $`E`$–type which has a smooth geometrical interpretation, since the rank of those groups cannot be made arbitrarily large.)
We presented the general scenario for the case of $`(p+1)`$–dimensions and exhibited and studied the orientifolded enhançon for the case of 2+1 dimensional gauge theory. Guided by the case of the $`D`$–type, where the fact that the O6<sup>-</sup>–plane has a known eleven dimensional supergravity description in terms of the Atiyah–Hitchin manifold, we were able to study aspects of all three new types: While the Atiyah–Hitchin manifold cannot be written explicitly, it reduces to (negative mass) Taub–NUT at large $`r`$ (up to exponentially small corrections in $`r`$) which was enough for us to study explicitly the relevant features of the supergravity solution which results from placing many D6–branes and an O6–plane on $`K3`$. This multi–Taub–NUT solution can also be reliably modified to capture the local asymptotic behaviour of the $`B`$– and $`C`$–type cases. The Atiyah–Hitchin structure imposes a global $`𝐙_2`$ identification on the entire geometry. Correspondingly, we found that there is a global $`𝐙_2`$ identification inherited by the enhançon locus, making the enhançon a natural $`\mathrm{I}\mathrm{RP}^2S^2/𝐙_2`$ geometry, in contrast to the $`S^2`$ geometry of the $`A`$–type. (We should also note that we observed that in all cases $`A,B,C,`$ or $`D`$, the apparent repulson singularity in eleven dimensional supergravity is naturally removed; not by M–branes, but by being forced back to ten dimensional heterotic string theory because the $`K3`$ becomes small. The heterotic string phenomena dual to the enhançon then take over the description.)
We displayed some leading $`1/N`$ corrections to the location of the all three types of orientifolded enhançon, as compared to the location of the $`A`$–type, and hence also the $`1/N`$ corrections to the one–loop metric on moduli space. Note that the $`A`$–type enhançon already has a series of exponential corrections of the form $`\mathrm{exp}(1/g_{\mathrm{YM}}^2)`$. On general grounds, our new types have a similar class of corrections, which can be phrased in terms of field theory instanton corrections, and equivalently, in terms of corrections from D1–brane world–sheets.
It is amusing to note that the $`1/N`$ corrections we studied here, which are of a different type, can all be written in terms of exponential corrections too. This follows from the fact that the part of the geometry of the Atiyah–Hitchin (–like) manifold that we neglected in writing the explicit supergravity solution is a series of exponential corrections in $`r`$. These corrections should also have an interpretation in terms of an instanton problem. Perhaps one can always organise the exponential corrections to these geometries in terms of structures reminiscent of the geometry of the Atiyah–Hitchin manifold, regardless of whether they are non–perturbative in $`g_{\mathrm{YM}}^2`$ or $`1/N`$.
## Acknowledgements
We would like to thank Peter Bowcock, George Papadopoulos, Simon Ross, Douglas Smith, Paul Sutcliffe and David Tong for comments and discussions. |
warning/0002/astro-ph0002327.html | ar5iv | text | # The 35-Day Evolution of the Hercules X-1 Pulse Profile: Evidence For A Resolved Inner Disk Occultation of the Neutron Star
## 1 Introduction
Her X-1 is a 1.24 second period accretion-powered X-ray pulsar in a 1.7 day circular orbit with a normal stellar companion, HZ Her. In addition to these basic periodicities this binary system has long been known to display an unusual 35-day long cycle of High and Low X-ray flux states. Within a single 35-day cycle are found a Main High and Short High X-ray flux state lasting roughly ten and five days each respectively and separated by ten day long Low states (see e.g. Scott & Leahy 1999). Pulsations are detected during the High states but cease during the intervening Low states. The High states are punctuated by deep X-ray eclipses every 1.7 days indicating a line-of-sight close to the binary plane.
A phenomenon currently known to exist only in Her X-1 is a repeating, systematic evolution of the pulse profile that occurs during the 35-day cycle. Observations with the Large Area Counters on Ginga and the Proportional Counter Array (PCA) on RXTE have allowed an unprecedented view of the evolution of the pulse profile in shape and energy spectrum across both the Main and the Short High states. The Ginga observations cover the energy range 1-37 keV and sample five Main and two Short High states. They are described in detail in Deeter et al. (1991), Scott (1993) and Deeter et al. (1998). The bulk of the data consists of a Main-Short-Main High state sequence in April-May-June of 1989. Lightcurves and softness ratios for the 1988 and 1989 observations as well as some pulse profiles can be found in Leahy (1995a). The RXTE observations cover a consecutive Main and Short High state in September-October of 1996 and a Short-Main-Short High state sequence in September-October 1997. The 1996 Observations with the PCA in the range 2-60 keV are presented in Scott et al. (1997a).
In this paper, we present a simple phenomenological model for the pulse evolution based upon the occultation of the central X-ray source by the inner edge of a tilted, precessing, accretion disk. This choice is motivated by the observed association of pulse shape changes with decreases in overall X-ray flux near the end of the Main High state (see Scott 1993; Scott et al. 1997a; Deeter et al. 1998; Joss et al. 1978 and Soong et al. 1990a) and the well known ability of a tilted, twisted, counter-precessing accretion disk to phenomenologically account for much of the complex 35-day optical and X-ray behavior displayed by the Her X-1/HZ Her system (Petterson 1975; Petterson 1977; Gerend & Boynton 1976; Crosa & Boynton 1980; Boynton, Crosa & Deeter 1980; Middleditch 1983). We stress that most of the observed 35-day behavior has been associated with the outer portion of the accretion disk whereas we will demonstrate that the pulse evolution must be an inner disk phenomenon.
Previous attempts to model the pulse shape evolution have relied on a combination of neutron star free precession and obscuration by the accretion disk (Trümper et. al. 1986; Kahabka 1987, 1989), obscuration by “flaps” of matter at the juncture of the accretion disk and pulsar magnetosphere (Petterson et al. 1991) and obscuration of the pulsar by a tilted precessing disk (Bai 1981; Averitsev et al. 1992). None of these previous models attempted to explain more than a few aspects of the pulse evolution. The model presented here refines the disk and pulsar beam geometry to qualitatively account for the observed pulse shape and its evolution during both the Main and Short High states.
A discussion of relevant aspects of the tilted, twisted, counter-precessing disk is presented in section 2. A summary of the main features observed in the pulse evolution during the 35-day cycle is presented in section 3. Section 4 discusses the absolute pulse phase alignment of the Main and Short High state pulses. Section 5 discusses the pulse evolution as a consequence of neutron free precession. Section 6 briefly discusses pulse evolution as a consequence of changing mass accretion patterns onto the neutron star. Section 7 presents a pulse evolution model based on an inner disk occultation and a simple X-ray beam configuration. In section 8 we discuss some implications of this inner disk occultation interpretation.
## 2 A tilted, twisted and counter-precessing disk in Her X-1
Her X-1 exhibits a rich variety of phenomena that appear to be well explained by an accretion disk that is tilted, counter-precessing and twisted. We review some of the observational arguments for such a disk (see Priedhorsky & Holt 1987 for an earlier review) and include some relevant new observations and interpretation. The occurrence of two distinct X-ray High states within the 35-day cycle has long been known. The Main High state was found soon after the discovery of Her X-1 (Giacconi et al. (1973)). The dimmer Short High state was first recognized in Copernicus observations (Fabian et al. 1973) and later in Ariel 5 and Uhuru observations (Cooke & Page 1975; Jones & Forman 1976). Few extensive observations of the Short High state were made until 1989 with Ginga (Deeter et al. 1998). In figure 1 we show the 1-37 keV lightcurve obtained with Ginga in 1989 and for the August 1991 Main High state. The August 1991 observation caught a turn-on to the Main High state which is confirmed by simultaneous monitoring at lower sensitivity with the Burst and Transient Source Experiment (BATSE) on board the Compton Gamma Ray Observatory (CGRO). The regularity of occurrence of the Short High state ten days after the end of the Main High state is now clearly demonstrated by the ongoing monitoring of Her X-1 with the RXTE All Sky Monitor (ASM) (see figure 1; Scott & Leahy 1999; and Shakura et al. 1999).
The companion star HZ Her is strongly heated by X-rays emanating from the neutron star and shows little variation in total magnitude, averaged over an orbital cycle, throughout the 35-day cycle (Gerend & Boynton 1976, Deeter et al. 1976). Optical pulsations produced by reprocessing of the primary X-ray flux have been observed during both the High and Low states (Middleditch & Nelson 1976, Middleditch 1983). These observations show that total X-ray production at the neutron star is relatively constant, eliminating gross periodic changes in accretion rate as the cause of the High-Low cycle. The strongest evidence for a tilted, counter-precessing disk in Her X-1 comes from the complex set of systematic changes in the optical orbital photometric light curve over the 35-day cycle that can be explained by a combination of disk emission and disk shadowing/occultation of the heated face of HZ Her (Gerend & Boynton 1976).
The occurrence of two distinct X-ray High states within the 35-day cycle has been attributed to obscuration of the central X-ray source by a tilted, twisted, counter-precessing accretion disk (Petterson 1975; Petterson 1977). Such a disk geometry can be idealized as a set of tilted, concentric rings in which the azimuth of the line-of-nodes of each successively smaller ring shifts smoothly as one moves radially inwards (Petterson 1977). The presence of two High states per 35-day cycle is naturally explained when the observer’s line-of-sight lies close to the binary plane. The appearance of pulsations with significant cold matter absorption at the onset of either High state (hereafter turn-on), is interpreted as the emergence of the pulsar from behind the tilted, outer rim of the accretion disk. The flux decline observed at the end of each High state shows little to no absorption effects and begins as the line-of-sight to the pulsar is approached by the hot, tilted inner edge of the precessing accretion disk. This behavior is illustrated in figure 2 and compared to a succession of High states observed with Ginga in 1989 and with an averaged 35-day lightcurve observed with the ASM on RXTE.
The direction of disk precession is apparently retrograde or “counter-precessing”. The evidence for this comes from X-ray and optical observations and theoretical considerations. X-ray absorption dips are observed in the Main High state just before eclipse that march toward earlier orbital phase as the High state progresses and at a frequency near but slightly lower than the sum of the orbital and 35-day frequencies (Crosa & Boynton 1980, Scott & Leahy 1999). The optical lightcurve of HZ Her exihibits a systematic pattern of changes and a harmonic decomposition revealed that nearly all power is confined to a discrete set of frequencies composed of sums of the orbital and 35-day frequencies (Deeter et al. 1976). A uniformly counter-precessing disk will repeat the same disk-star orientation at the sum of the orbital and 35-day frequencies as will any phenomena dependent on the orientation. Prograde precession should cause phenomena to repeat at the difference of the orbital and 35-day frequencies but no such phenomena are observed. Theoretically, a tilted disk should also precess in a retrograde fashion (Katz 1973).
### 2.1 The outer disk and early High state behavior
X-ray pulsations appear during the flux rise at the start of the High states (i.e. the turn-on) accompanied by cold matter absorption in the X-ray spectrum (e.g. Parmar, Sanford & Fabian 1980). The only Short High state turn-on that has been well observed to date is the May 1989 Short High state which is compared in figure 3 to the August 1991 Main High state turn-on. The turn-on’s are nearly identical in form and both show an increase in softness ratio characteristic of cold matter absorption. The turn-on’s are modeled in figure 3 by an X-ray point source emerging through an atmosphere with a gaussian density profile lying above the plane of a tilted, precessing disk. The disk angular velocity is $`2\pi /34.85`$ $`\mathrm{day}^1`$, the scale height is $`1/24`$ the disk radius and the disk is tilted at $`20^{}`$ with respect to the orbital plane of Her X-1. The optical depth at the base of the disk atmosphere is 30. The disk model is described in more detail in section 7. The two turn-on’s last about 3 hours in contrast to the eclipse egress which lasts only a few minutes (e.g. Leahy 1995b). The pulse profile exhibits no significant changes during the turn-on but merely increases in flux in different energy bands (Deeter et al. 1998). The X-ray observations coupled with the optical observations imply that the primary cause of the High-Low flux cycle is obscuration.
The flux after the beginning of the Main High state often shows a gradual increase by 20-50% over the next 1-4 days. The pulse profile is relatively constant in shape during this period. We propose that the X-ray flux rise is the result of viewing the neutron star through a hot dense lower corona lying just above the outer accretion disk surface. As the elevation of the observer’s line-of-sight increases with respect to the nominal outer disk plane the amount of obscuration will decrease. This effect should also be present in the Short High state and in the system LMC X-4 where a tilted precessing disk is also postulated to occur (Lang et al. 1981). The existence of another much larger and lower density scattering corona is indicated by the existence of Low state flux at 5% of the peak Main High state flux (e.g. Choi et al. 1997; Mihara et al. 1991). A two layer disk corona has been discussed theoretically by Schandl & Meyer (1994). The lower corona has a temperature of $`10^6`$ K while the upper corona has a temperature of $`10^{7.5}`$ K. The observed $`50\%`$ flux increase implies a column density of $`1\times 10^{24}`$ $`\mathrm{cm}^2`$ at the base of the lower corona for pure Thompson scattering. From figure 2, the duration of the flux increase implies a vertical angular thickness to the lower corona of $`4^{}14^{}`$ with respect to the nominal outer disk plane.
### 2.2 The inner disk and late High state behavior
If pure obscuration by a precessing, tilted, thin, and planar disk were the cause of the High-Low flux cycle then one might expect to observe: 1) rapid flux cutoffs at the ends of the High states equivalent to the High state turn-on’s 2) identical pulse profiles during the Main and Short High states and 3) nearly identical fluxes during the Main and Short High states except for variations caused by geometric differences in coronal obscuration. In contrast, both types of High state show gradual flux declines that last several days and without significant absorption effects. A tilted and twisted disk can explain the gradual flux decline if the disk is twisted in the direction of precession such that the azimuthal angle between the outer and inner disk line of nodes is $`>90^{}`$. In a tilted and twisted disk the line-of-nodes and tilt of individual contiguous disk rings varies smoothly as one moves from the outer to inner disk radii. The hot, inner region of the disk gradually covers the X-ray emitting region during a transition to a Low state. We illustrate this type of disk model in figure 2 and compare it to the lightcurve observed with Ginga and an average lightcurve from the RXTE/ASM. A tilt of $`\theta _{tilt,outer}=20^{}`$ for the outermost ring is determined from the assumption of an observer elevation of $`\alpha _{obs}=5^{}`$ and the observed 35-day phase separation of $`\mathrm{\Delta }\psi 0.58`$ for the Main and Short High state turn-on’s <sup>1</sup><sup>1</sup>1$`\theta _{tilt,outer}=\frac{\alpha _{obs}}{\mathrm{sin}(\pi (\mathrm{\Delta }\psi 0.5))}`$. The required outer disk tilt is independent of the outer disk thickness. We note that the geometry of the disk model proposed in Schandl & Meyer (1994) is inconsistent with the observations since it predicts cold matter absorption at the start of the Main High state with a gradual flux decline caused by a covering up of the neutron star by an inner disk corona and the same events, but in opposite sequence, during the Short High state (e.g. see their figure 12). However, the same sequence of events is observed in both the Main and Short High states.
The overall spectral changes during the Main High state were further explored by comparing the 20-50 keV pulsed flux average Main High state lightcurve observed with BATSE with a 2-12 keV flux average Main High state observed with RXTE/ASM. Following the procedure described in Scott & Leahy (1999), the BATSE pulsed flux light curve over the timespan MJD 49933 to MJD 50507, obtained from folded-on-board data (see Bildsten et al. 1997), was sorted into orbital phase “0.2” turn-on Main High states or “0.7” orbital phase turn-on’s and averaged. Seven Main High states were used to construct the average 0.2 turn-on Main High state lightcurve and eight for the average 0.7 turn-on Main High state. Likewise, for the 2-12 keV RXTE/ASM a similar sorting and folding was done to construct average light curves for the timespan MJD 50146 to MJD 50947 with 12 and 10 Main High states averaged to form, respectively, the average 0.2 turn-on and 0.7 turn-on Main High state lightcurves. During the flux decline at the end of the Main High state, the 20-50 keV pulsed flux dropped to the level of the background flux more than one day preceding a similar drop in the 2-12 keV flux in both turn-on type Main High states. In figure 4 we compare the softness ratio formed by the two lightcurves. The turn-on at 35-day phase 0.0 shows a rapid increase in softness ratio consistent with decreasing absorption. The pre-eclipse dips also show up as decreased softness consistent with absorption. However, the softness ratio shows a large increase during the flux decline near the end of the Main High state followed by a large decline. The softness ratio increase is incompatible with either absorption or an energy independent Thomson scattering of an unresolved point source by a corona or disk atmosphere (as in the Schandl & Meyer 1994 model) as the sole cause of the flux decline. Therefore the flux decline at the end of the Main High state cannot be the result of an occultation of a point source by either a cold or a hot disk edge.
The peak Short High state X-ray flux is only 30% that of the peak Main High state flux and exhibits a quite different pulse profile. Dramatic changes in the pulse profile are observed during the last few days of the Main High state and throughout the Short High state (see Deeter et al. 1998 and next section). If an obscuring region causes the gradual flux decline and the density scale height was much larger than the linear size of the pulse emitting region then minimal pulse shape changes would be observed as well as little difference in pulse profile between the Main and Short High states.
Two possible explanations for the High state flux declines, pulse shape evolution and the low Short High state flux are 1) Systematic changes in the X-ray beaming direction are occurring in addition to those caused by neutron star rotation and/or 2) The scale height of a precessing obscuring inner disk region is indeed comparable in size to the pulse emitting region. If explanation 1) is correct then we are observing a combination of an obscuration of a point source causing the flux declines and beaming changes that cause both pulse shape changes and the Main and Short High state flux difference. This possibility has been advocated by Trümper et al. (1986) and Kahabka (1989) using beam changes caused by free precession of the neutron star coupled with obscuration by a precessing disk. In 2) the flux declines, the Main and Short High state flux and pulse shape differences and the pulse shape changes are caused purely by progressive disk occultation of an extended source. We will discuss both these explanations in more detail after reviewing the phenomenology of the pulse evolution and phase alignment of the Main and Short High state pulses.
## 3 Phenomenology of the pulse evolution in Her X-1
We now review the salient features of the pulse profiles and their evolution presented in figures 5 and 6 and documented in Scott (1993) and Deeter et al. (1998). The early 1-37 keV Main High state pulse profile consists of a large main pulse and a smaller interpulse superposed on an underlying weakly pulsed component (see figure 5 for profiles and nomenclature of specific pulse features). The main pulse consists of two unequal shoulders (or peaks) at energies below 5 keV, but at higher energies a third central peak appears that grows with energy and dominates the main pulse profile above $`20`$ keV. These energy dependent features of the pulse are well known (see e.g. Soong et al. 1990b). We display a subset of the Ginga observations in figure 6 showing the evolution of the pulse profile during the Main and Short High states displayed in figure 1. From Ginga and other observations, we propose that the Main High state pulse profile evolution consists of three basic events listed in order of decreasing duration: 1) A deepening and widening of a “gap” in the underlying component near the pulse phase of the preinterpulse minima that begins early in the Main High state and continues until the end (i.e. pulse phase $`0.3`$ in figure 6, top panels). This phenomenon can also be observed in Uhuru and HEAO 1 pulse profiles (Joss et al. 1978; Soong et al. 1990a) and in recent RXTE observations (Scott et al. 1997a). The well known quasi-sinusoidal pulse profile that “appears” near the end of Main High state may simply be the uncovering of this already present gap due to the disappearance of the overlying main pulse. The quasi-sinusoidal profile is the last pulsed feature to disappear before the High states end. 2) The disappearance of the leading and then the trailing shoulder of the main pulse over a roughly two day period that starts six to seven days after the Main High state turn-on. The main pulse shows relatively little change before this point. The disappearance of the leading shoulder of the main pulse near the end of the Main High state has been noted many times previously (e.g. Soong et al. 1990a; Joss et al. 1978; Kahabka 1989; Sheffer et al. 1992; Scott 1993; Scott et al. 1997a; Deeter et al. 1998). 3) A rapid decay and disappearance of the spectrally hard central peak of the main pulse was observed with Ginga over an $`12`$ hour period that took place within the time span of the shoulder decay. A similar rapid decay of the main pulse was observed by HEAO 1 but at lower resolution (Soong et al. 1990a). In the context of the occultation model described below, this pulse shape evolution pattern suggests the presence, respectively, of three pulse emitting regions of decreasing size each roughly centered on the pulsar.
Two Short High states have been observed in enough detail to follow the pulse evolution. These are the Ginga observations shown in figure 1 and recent RXTE observations (see Scott et al. 1997). However with these two Short High state observations, combined with the fragmentary observations of other Short High states, we can describe the following evolution pattern. As the Short High state commences, the pulse profile is quite different in shape and lower in flux by $`70\%`$ relative to the Main High state pulse. Both Ginga and Exosat observations of the Short High state profile reveal a small hard peak and a larger soft peak separated by $`180^{}`$ in pulse phase and superposed on a quasi-sinusoid (see figure 6). During Exosat observations, the small hard peak was actually larger than the soft peak at energies above $`13`$ keV (Kahabka 1987, 1989). Ginga observations also showed the small hard peak increasing in amplitude relative to the soft peak with increasing energy, but not exceeding that of the soft peak. The RXTE observation of the November 1996 Short High state did not reveal the presence of the small hard peak. The Ginga observations show the small hard peak disappearing within one day of the turn-on as did the earlier Exosat observation (Kahabka 1987). The Ginga and RXTE observations showed that the soft peak also declined in flux, but more slowly, and disappeared three to four days after turn-on. A narrowing in width occurred in both cases.
The Short High state soft peak contains an even softer component on the trailing side of the peak indicated by a spectral softening at approximately pulse phase $`0.6`$ (see Fig. 5). The very soft component is apparent as long as the soft peak is present. The Exosat pulse profiles presented by Kahabka (1987, 1989) also show this very soft component on the trailing side of the soft peak.
The amplitude of the quasi-sinusoid first increases and then decreases as the May 1989 Short High state progresses, and is about $`180^{}`$ out of phase compared with the Main High state quasi-sinusoid (see section 4). A similar flux increase of the quasi-sinusoid can be seen in the Short High state profiles displayed in figure 4.3 of Kahabka (1987). The overall flux stays relatively constant for about four days following the Short High state turn-on, but this is due to a flux increase in the quasi-sinusoid just balancing a flux decrease in the small hard peak and the soft peak. The overall flux of the quasi-sinusoid is roughly half that of the Main High state quasi-sinusoid. As in the Main High state, the gap that defines the quasi-sinusoid shows a decrease in width and depth at increasing energies.
In summary, the Main High state pulse evolution involves a decay preferentially on the leading edge of the main pulse and interpulse that begins late in the High state. In addition, the formation and continuous slow evolution of an underlying quasi-sinusoidal profile may also be occurring. The Short High state involves the narrowing, decay and disappearance at very different rates of the two peaks in the profile superposed on an underlying quasi-sinusoid. Evolution of the quasi-sinusoid is much slower and involves only subtle changes in the profile. These changes are illustrated in figure 6. Overall, comparison of the Ginga observations with other observations of the pulse evolution are consistent with a repeating, stable and systematic pattern of change in the pulse profile. Future observations are needed to explore the details and stability of the pulse evolution during the Short High state and especially the Main High state flux decline.
## 4 Pulse Phase Alignment of the Main and Short High state pulses
The overall pulse evolution pattern can only be completely understood if the proper phase alignment of the Main and Short High state pulse profiles is known. To properly align the Main and Short High state profiles two methods might be tried: 1) extrapolating the pulse timing ephemeris between High states across the Low state where pulsations are unobservable or 2) matching pulse features that are common to each profile.
Figure 5 displays a Main High state pulse profile, a closeup of the interpulse of the same Main High state profile, and a Short High state profile. Both profiles are taken from an early point in their respective High states when the effects of the pulse profile evolution are minimal. The two High state profiles were phase aligned using the pulse timing extrapolation given in Deeter et al. (1998), in which a pulse phase ephemeris is extrapolated from the April and June 1989 Main High states into the June 1989 Short High state.
At first glance the Main and Short High state pulse profiles seem quite different but there are actually a number of features common to each profile. The bottom of each panel has a hardness ratio for the profile. Note that the hardness ratio of the Main High state profile shows a dip at pulse phase 0.6 and a spectral hardening near pulse phase 1.0. The Short High state profile also possesses a similar soft dip and spectral hardening separated by 0.4 in pulse phase.
The phase alignment suggested by the spectral features implies that the Main High state interpulse and the Short High state peak soft peak should be matched together. Examination of figure 5 shows that these two features exhibit very similar shapes, fluxes and energy dependence. The primary differences between the two features may simply be due to the different backgrounds upon which each feature rests, since the underlying quasi-sinusoid is shifted by 0.5 in phase between the two High states.
The spectral softening at pulse phase 0.6 in either High state profile is apparent in all the High states observed by Ginga with the exception of the anomalous June 1989 Main High state. The spectral softening is also readily apparent in both the Main and Short High state profiles of the Exosat observations (see Kahabka 1987, 1989). Energy dependent differencing of the Ginga pulse profiles was used to isolate this soft excess feature. A pulse profile in a high energy band was scaled until a fit to the surrounding pulse in a lower energy band was obtained and then the difference was taken. This process revealed that the soft excess feature is very similar in shape, flux and color in all the High states which suggests that it is the same feature in both the Main and Short High states (see Scott 1993). The presence of the same soft excess feature in both the Main High state interpulse and Short High state soft peak strongly supports the phase alignment of figure 5.
The phase alignment also suggests that the hard central peak of the Main High state pulse profile and the small hard peak of the Short High state profile are corresponding features. There is a significant flux difference between these two features but both hard peaks increase in width and amplitude at higher energies relative to the surrounding pulse. The width increase of the small hard peak is only marginally observable in figure 5, but is readily apparent in Kahabka (1989) and in RXTE/PCA observations. In the Exosat observation of the Short High state shown by Kahabka (1989), the small hard peak is actually larger than the soft peak at higher energies.
We can now describe the following similarities and differences of features in the Main and Short High state profiles based upon the phase alignment of figure 5. The Main High state interpulse and the Short High state soft peak are the same feature, as is the soft excess feature within each profile. The Short High state small hard peak is a greatly reduced version of the hard central peak of the Main High state. The large soft two peaked component underlying the hard central peak in the Main High state is absent in the Short High state profile. The quasi-sinusoidal profile of the Short High state is reduced in flux by about 50% compared to the Main High state quasi-sinusoid and shifted in phase by roughly 0.5. The Short High state profile can be viewed as a modified version of the Main High state profile rather than a completely distinct pulse.
The initial appearance and probable phase alignment of the Short High state pulse profile require an explanation of 1) The disappearance of the soft shoulders of the Main High state main pulse. 2) The large drop in flux of the Main High state hard central peak 3) The nearly equivalent amplitudes of the Main High state interpulse and the Short High state soft peak. 4) The $`50\%`$ drop in flux and $`180^{}`$ phase difference between the two High state quasi-sinusoidal profiles.
## 5 Free precession of the neutron star
Force-free precession of the neutron star with a fixed beam geometry has been proposed as the cause for the pulse shape change between the Main and Short High states by Trümper et al. (1986) and was later elaborated in Kahabka (1987), Ögelman & Trümper (1988) and Kahabka (1989). They noted that free precession alone cannot explain the overall High-Low cycle since the beams needed to model the observed Main High state pulse are too wide to disappear from view and cause a Low state. For example, the unpulsed and quasi-sinusoidal components of the profile must be produced by unbeamed emission or very wide beams, so free precession cannot cause these pulse components to disappear at the end of a High state. A model based solely on free precession also has great difficulty explaining the cold matter absorption seen during turn-on in both kinds of High states. Therefore, a periodic obscuration of the X-ray source by the accretion disk is also assumed to explain the disappearance of the pulses and the overall High-Low cycle. A composite model with a precessing disk acting as an occulting body and a freely precessing neutron star in which both cause pulse shape changes has been proposed by Kahabka (1989). The period and phase of both the precessing disk and the freely precessing neutron star must be closely locked together to prevent longterm drift between the pulse evolution pattern and the High-Low intensity cycle (for example, the leading edge decay of the main pulse has only been observed at the end of the Main High state). The observed random walk in the phase of the Main High state poses some difficulties for phase locking since a freely precessing neutron star should be a relatively good clock (see Boynton, Crosa & Deeter 1980; Trümper et al. 1986 and Baykal et al. 1993).
A spinning body can exhibit periodic, force-free precession if two of its three principal axes of inertia are unequal (i.e. a symmetric top with $`I_1=I_2I_3`$). A review of the basic equations of free precession can be found in Bisnovatyi-Kogan, Mersov & Sheffer (1990). Free precession can cause major changes in the observed pulse profile that repeat over the course of successive precession cycles. The pulse profile at a given pulse phase changes due to the slow variation in the angle between the magnetic dipole axis and the observer’s line-of-sight. Beams emanating from the magnetic polar caps or from any fixed location on the neutron star (other than the figure axis) will appear to a distant observer to “wobble” in rotational latitude over the precession cycle. Either end of the dipole axis, when facing the observer, will oscillate sinusoidally in latitude during one free precession cycle. The observed free precession light curve is highly dependent on the angle between the angular momentum axis and the observer’s line-of-sight and the emission beam geometry and will be symmetric about the precession phases of the latitude extrema in the dipole motion. The sensitivity to the beam profile depends mostly on the beam width. A narrow beam produces pulse components that can appear and disappear from view as the wobbling takes place while components from wider beams will merely show variations in amplitude.
Another potentially observable signature of free precession is a characteristic variation in the pulse period over the free precession cycle (Bisnovatyi-Kogan, Mersov & Sheffer 1990; Bisnovatyi-Kogan & Kahabka 1993). Measurement of this effect was initially a major motivation for observing Her X-1 with Ginga. We can estimate the size of the period change expected between the Main and Short High state using equation 7 of Bisnovatyi-Kogan & Kahabka (1993) to be $`5.0\times 10^7`$ s. However the known random walk in pulse frequency (see Bildsten et al. 1997; Deeter 1981) also causes pulse period changes. The expected period change between the Main and Short High state is given by:
$$\mathrm{\Delta }P^2^{\frac{1}{2}}=P_0^2(St)^{\frac{1}{2}}$$
(1)
and has a value of $`8.0\times 10^7`$ s for a noise strength $`S=1.8\pm 0.8\times 10^{19}`$ $`\mathrm{Hz}^2\mathrm{s}^1`$ and $`t=17.5`$ days. The expected period changes are thus comparable and would be difficult to separate unless many Main-Short-Main High state pulse period measurements could be made. Currently, the most observable potential manifestation of free precession is the pulse evolution observed between and within the High states.
The ability of free precession to account for the change in pulse profile between the Main and Short High states was tested by Kahabka (1987) using Exosat observations. The Main High state pulse profile was modeled using a gaussian pencil beam for the hard central peak and a coaxial gaussian fan beam for the soft leading and trailing peaks. The interpulse is produced by an identical antipodal fan beam. The opening angle of the fan beam was determined to be about $`45^{}`$. We show a sequence of pulse profiles evolving over the 35-day cycle using the parameters determined by Kahabka (1987, 1989) in figure 7. The pulses occur one day apart over the 35-day cycle and do not include the unpulsed or quasi-sinusoidal component. The Main High state must be centered roughly about day zero with the Short High state occurring 18 days later. From the figure several features can be noted: (1) The interpulse and the soft two peak component are always visible. (2) The interpulse flux increases (or decreases) as the main pulse flux decreases (or increases). (3) No leading edge decay of either the main pulse or interpulse occurs. (4) The pulse profile evolves smoothly and slowly over the 35-day cycle, with no rapid profile changes. (5) The decrease in pulsed intensity is only about 10%-20% between the Main and Short High state. All of these features are inconsistent with the evolution actually observed as described in section 3. Related criticisms of free precession have been made by Bisnovatyi-Kogan, Mersov & Sheffer (1990). In Kahabka (1989), a precessing disk was assumed to cause the leading edge decay of the main pulse during the termination of the Main High state with some contribution from free precession. However, this model has difficulty in explaining the actual Short High state profile and the subsequent disappearance of both peaks in the profile as noted in points (1) and (2) above.
A freely precessing neutron star with approximately axisymmetric fan and pencil beams has several general problems in explaining the observed pulse evolution in Her X-1. (1) The pencil beam cannot disappear from view without the fan beam disappearing as well. The small hard peak of the Short High state appears to be the same pulse component as the hard central peak of the Main High state (see section 4), but without the surrounding soft two peaked component that is presumably produced by a fan beam. (2) The preferential decay of one side of a fan beam cannot occur due to free precession. The two cuts across a fan beam, which result in a two-peaked pulse component, should change in amplitude and width simultaneously as precession occurs. (3) Rapid profile changes, such as those observed near the end of the Main High state or the beginning of the Short High state, are difficult to explain unless the beam profile has sharp edges. For example, the flux of the hard central peak dropped by more than 70% in 6.7 hours near the end of the April 1989 Main High state. If we assume that the hard central peak is produced by a pencil beam and that the precessional motion of the pulsar is responsible for its amplitude decrease, then we can estimate the latitudinal width of the pencil beam. We conservatively estimate that the hard central peak flux drops to zero in approximately 12 hours. In this period, the beam axis moves approximately through an angle $`4\mathrm{\Phi }\mathrm{\Delta }t/P_{35}`$ away from the observer’s line-of-sight. With $`\mathrm{\Phi }12.5^{}`$, this produces a latitudinal width over this portion of the beam of about $`0.7^{}`$. We can also estimate the longitudinal width of this portion of the beam to be $`40^{}`$ from the observed pulse shape as the beam sweeps across our line-of-sight. The beam responsible for the hard central peak would have to be $`60`$ times broader in longitude than in latitude, a quite improbable situation.
To explain a rapid change in the pulse profile observed with HEAO 1 near the end of a Main High state (Soong et al. 1987) with a free precession model, Shakura, Postnov & Prokhorov (1998) have invoked a temporary transition from a symmetric top to a triaxial neutron star caused by a “quake” during the Main High state flux decline. The quake induces a rapid shift in the rotational latitude of the magnetic axis and thus induces rapid changes in the observed pulse profile. This explanation may work for one case of observed rapid pulse profile change but the rapid decline of the pulse during the Main High state now appears to be a normal and repeating phenomenon rather than an anomalous event (see Deeter et al. 1998). It seems rather improbable that neutron star quakes could be arranged to regularly occur during the 35-day phase of the Main High state flux decline but not to occur at other times. In addition, there appears to be no evidence for an expected phase shift caused by the quake based on Ginga and RXTE observations. In conclusion, we find little evidence supporting neutron star free precession as the cause of the pulse shape changes in Her X-1.
## 6 Other causes of systematically recurring pulse evolution
If the neutron star is not undergoing free precession then the precessional motion of the inner disk becomes the primary candidate causing pulse evolution. Pulse shape changes might be caused by modulation of matter flow onto the magnetic field lines by a changing aspect between the magnetosphere and a tilted, precessing accretion disk. The cycle of High and Low states can be attributed to obscuration by the precessing accretion disk while the pulse shape evolution results from variations in the accretion column structure induced by the changing pattern of matter entry onto the magnetic field lines. As a naive example of a possible variation, assume the accreting matter attaches directly to a dipole magnetic field at a radius $`R_m`$ and then travels along the field lines onto the magnetic poles. The angular extent of the accretion cap will be given by $`\mathrm{sin}^2(\theta _C)=R/R_m`$ where $`R`$ is the radius of the neutron star (Lamb, Pethick & Pines 1973). An increase in $`R_m`$ will cause a decrease in the accretion column radius at the neutron star surface and concentrate the energy released by the accretion into a smaller area. A decrease in $`R_m`$ will have the opposite effect. A full evaluation of the coupling between a tilted, precessing accretion disk, the neutron star magnetosphere and the net effect on the pulsar beam involves complex physics that has not been undertaken as yet and is beyond the scope of this discussion.
It seems plausible that some degree of systematic pulse evolution should be caused by a changing disk orientation, but we argue that this is unlikely to be the primary cause of the pulse shape changes observed in Her X-1. Nonlinear changes in the flow rate and pattern could be responsible for rapid changes observed in the pulse profile, but these rapid changes must be scheduled to occur just as the overall flux declines at the end of the Main High state and not at other times. This mechanism should also cause both increases and decreases in the pulse profile. In the visible High states only decreases in pulse features are ever observed.<sup>2</sup><sup>2</sup>2During the Short High state the small hard peak and the soft peak disappear, but this is accompanied by a compensating small flux increase in the quasi-sinusoidal component that occurs early in the Short High state. The flux of the quasi-sinusoid subsequently declines. The early flux increase of the quasi-sinusoid may simply be part of an overall increase in flux as the line-of-sight passes through a decreasing density of coronal disk material at the start of the High state. The most complicated pulses are apparent at the beginning of the High states with a subsequent disappearance of features as the High state progresses. Nature would have to conspire so that increases in the amplitude of pulse profile components occur only during the Low state.
Petterson et al. (1991) have qualitatively presented a time dependent disk obscuration model for the pulse evolution that relies on obscuration provided by “flaps” of matter at the points where a tilted accretion disk meets the neutron star’s magnetosphere. These “flaps” rotate at the pulse period and move in pulse phase as the tilted disk precesses. It is not clear whether the physics of disk-magnetosphere coupling will produce obscuring flaps of this type, but in any case the model has some qualitative problems when confronted by the observed pulse profiles. During the Short High state, the interpulse flux for the “flaps” model is larger than during the Main High state (compare Figs. 2b and 3b of Petterson et al. 1991). The Ginga observations show that the soft peak flux is the same or smaller during the Short High state than that of the Main High state interpulse with which it is identified (see figure 5). Similar behavior in the Exosat observations can also be seen in Fig. 1 of Petterson et al. (1991). During the Short High state the “flaps” model predicts that the main pulse flux should decrease due to increasing “flap” obscuration, while the interpulse flux should increase, contrary to the Ginga observations. In the Main High state, the main pulse flux and the interpulse flux should also change in an opposing manner according to the “flaps” model, but both fluxes are observed to decrease.
## 7 Pulse evolution resulting from a resolved occultation of the neutron star by a tilted, precessing accretion disk
The apparent stability and repeating nature of the pulse evolution may be the result of a resolved occultation of the pulse emission region by the inner edge of a tilted, precessing accretion disk. Two different occultation sequences, and hence two different sets of pulse shape changes will occur, as the observer sees the inner disk edge sweeping alternately ‘upwards’ and ‘downwards’ across the pulse emitting region and these events will be $`180^{}`$ apart in disk precession phase. The pulse shape changes observed during both High state types require the “scale height” of the disk (or more generally the “occulter”) to be roughly the size of the dominant pulse emitting region i.e. a few neutron star radii, which will then naturally produce pulse shape changes associated with decreases in X-ray flux. To pursue the occultation idea further a simple geometric model is developed and qualitatively compared to the observations. The model proposed below is a considerably refined version of the model originally proposed by Bai (1981).
A disk occultation model requires at least two components. A model for the tilted and precessing disk itself is necessary, and a model for the pulsar emission geometry. In the simplest approximation, the pulsar beams do not depend on the azimuth of the disk and the disk simply occults the pulse emitting region. This is probably not entirely true since the disk is coupled to the neutron star through the magnetosphere and the azimuthal progression of a tilted disk may cause some changes in the accretion column and hence the pulse shape. However, these complexities will be ignored for now and the beams from the pulsar are assumed to be decoupled from the disk. With the choice of a physical location for the beam emission region with respect to the neutron star, the apparent spatial location of each beam component as seen by the observer can be computed. Likewise, the occulting disk may have also have a complex geometry due to interaction with the magnetosphere, among other effects, but we will assume a simple planar disk shape. The sequence of pulse shape changes can be modeled by sweeping the disk over the pulse emission region and comparing the predicted pulse shape changes with the observations.
The least obscured pulse should occur early in the Main High state after the pulsar emerges from behind the outer disk rim. The main and interpulse profiles are clearly asymmetric about their maxima at this time but we will provisionally make several assumptions to simplify the modeling process. We will first assume that the beam is axisymmetric, the magnetic field is purely dipole and that identical beam emission regions exist at the ends of an axis that is close to but not necessarily identical with the magnetic dipole axis. We attribute the hard central peak of the main pulse to a pencil beam directed along the beam axis and the softer flanking shoulders to a surrounding fan beam (similar to Kahabka 1987). The interpulse is produced as the edge of the fan beam emanating from the antipodal magnetic pole grazes the observer’s line-of-sight.
We assume gaussian intensity profiles for the beams since Kahabka (1987) has shown that the pulse profile of Her X-1 can be well fit with a small number of gaussian components and this is confirmed by fitting the Ginga profiles. A simple gaussian beam model for both the fan and pencil beams is used and is given by:
$$I(\theta )=I_{pen}\mathrm{exp}(\theta ^2/\sigma _p^2)+I_{fan}\mathrm{exp}((\theta _{cone}\theta )^2/\sigma _f^2)+I_D+I_E+I_{Low}$$
(2)
where $`\theta `$ is the angle from the beam axis, $`I_{pen}`$ and $`I_{fan}`$ are the pencil and fan beam amplitudes and $`\sigma _p`$ and $`\sigma _f`$ are the beam widths. The fan beam opening angle is given by $`\theta _{cone}`$. Three constant flux components are present, two due to magnetospheric emission ($`I_D`$ and $`I_E`$) and one due to Low state coronal emission $`I_{Low}`$. Values for the beam components are given in Table 1. The soft two shoulder component in the Main High state main pulse results from a cut of the line-of-sight across the two edges of the fan beam “cone”, while the interpulse results from a grazing cut along the edge of the “cone”. The large difference in the amplitudes of the main pulse and the interpulse require the line-of-sight to be $`20^{}40^{}`$ from the neutron star’s rotational equator. Since previous optical and X-ray observations show that the observer is offset by $`510^{}`$ relative to the binary plane (Gerend & Boynton 1976; Middleditch 1983; Deeter et al. 1991) the neutron star rotation axis must be inclined by $`10^{}50^{}`$ to the orbital axis.
To model the physical location of the primary pulsar beams requires some additional assumptions. Theoretically, beam models have been divided into slab and column geometries depending on the physical mechanism assumed to decelerate the infalling plasma. Column models assume that either a radiation pressure shock or a collisionless shock lies above the neutron star surface and decelerates the flow. Brainerd & Mészáros (1991) have shown that the radiation pressure for the luminosity of the Her X-1 is too weak to significantly decelerate the infalling plasma. It thus seems prudent to assume a slab geometry for the magnetic polar cap of Her X-1. In the slab model the infalling plasma is decelerated at the neutron star surface and the emitting region is a thin cap only a few meters in height (Mészáros & Nagel 1985). In more complex models, X-ray radiation emitted from the foot of the accretion column is backscattered as it rises through the accretion column (e.g. Brainerd & Mészáros 1991). Thus the beam consists of direct emission from the polar cap and a backscattered component.
We will first consider a simple pulsar model in which a single emitting point lying on the neutron star surface is chosen to approximate the beam emission region. The pulse emission region will have a width of $`2R_{ns}`$, where $`R_{ns}`$ is the neutron star radius. An observer will see the emitting point sweep out an ellipse as the neutron star rotates with the surface blocking the emitting point for a portion of the rotation period. We will refer to this beam geometry as a “direct fan beam” geometry. The beam pattern originating from the emitting point consists of a pencil beam surrounded by a concentric fan beam as in equation 2 (see figure 8, top panels). This was the same beam pattern and emission geometry used in the free precession model in section 5.
We will also consider a second beam/emission model in which the fan beam emission is produced by backscattered radiation from the accretion column that is focused around the neutron star and beamed in the antipodal direction, a “reversed fan beam” geometry. This beaming configuration will produce a similar pulse profile but the spatial locations of the pulse components will be significantly altered (see figure 8, bottom panels). The fan beam components will now be observed to originate at some distance above the neutron star surface and $`\theta _{cone}>90^{}`$. The interpulse will now be emitted by the same magnetic pole that produces the hard central peak while the soft shoulders of the main pulse come from the opposite pole. A similar fan and pencil beam configuration has been discussed theoretically by Brainerd & Mészáros (1991). In their model, a fan beam is produced from magnetic polar cap radiation that is preferentially backscattered by the incoming accretion flow and then gravitationally focused around the neutron star. The accretion column is calculated to be optically thin to Thomson scattering while the fan beam photons are produced by cyclotron resonance scattering. Only photons at energies less than or equal to the surface cyclotron frequency will be scattered in the accretion column. Support for this pulse profile interpretation is provided by the observed cyclotron absorption feature in Her X-1, which reaches a maximum absorption depth at the pulse phase of the hard central peak and by the disappearance of the main pulse shoulders above the $`38`$ keV cyclotron line energy (Soong et al. (1990b)).
The location of the reversed fan beam emission was modeled by assuming the emission occurred from a point at a height of $`1.5R_{ns}`$. This is the cyclotron scattering height of 10 keV photons for a surface field strength corresponding to a cyclotron line energy of 40 keV, a simple dipole field and a neutron mass of $`1.4\mathrm{M}_{}`$. Softer photons will scatter from higher up and harder photons from lower down so the apparent emission region location is energy dependent. The pulses displayed in figure 8 approximately model the Main High state pulse profile observed in the Ginga $`9.314`$ keV band. Gravitational lightbending of the photon trajectories causes the photons to appear to be emitted from a region higher above the neutron star surface than for the case when no lightbending is present. We illustrate this effect in figure 9 where trajectories are calculated using the equations taken from Brainerd & Mészáros (1991) and Riffert & Mészáros (1989). A reversed fan beam with an opening angle of $`\theta _{cone}=140^{}`$ with respect to the accretion column and a profile with a half width of $`\sigma _f=20^{}`$ was used to calculate the nonattenuated fan beam intensity since this was a good approximation to the beam pattern shown in figure 9. The location of the fan beam emitting point was calculated using a sinusoidal approximation to the photon impact parameter for the case of lightbending in figure 9.
The disk is modeled as an infinite plane with a circular hole centered on the neutron star. The circular hole has a radius $`R_{inner}`$ and a gaussian density profile in the direction perpendicular to the disk plane characteristic of an $`\alpha `$disk of Shakura & Sunyaev (1973). The vertical disk density profile is described by:
$$\rho (z)=\rho _0\mathrm{exp}(z^2/\sigma _d^2)$$
(3)
where $`z`$ is the vertical distance above the inner disk midplane, and requires the specification of the free parameters $`\rho _0`$, the midplane density of the disk and $`\sigma _d`$, the disk “scale height”. To model the occultation, the disk is assumed to be simply a linear edge with a minimum distance $`R_{occ}`$ from the neutron star as seen by the observer.
The disk selectively obscures the pulse emitting region. When the observer’s line of sight lies at an angle $`\theta _{occ}`$ above the disk midplane the disk edge will appear to be at a distance $`R_{occ}=R_{inner}\mathrm{sin}(\theta _{occ})`$ from the neutron star. The optical depth of the disk material will be caused by Thomson scattering due to the complete ionization of hydrogen and helium in the inner disk region. The optical depth will therefore follow the density distribution and is assumed to scale as:
$$\tau (z)=\tau _{disk}\mathrm{exp}(z^2/\sigma _d^2)$$
(4)
where $`\sigma _d`$ is the optical depth scale height, and $`\tau _{disk}`$ is the optical depth at $`z`$=0.0.
The total attenuation caused by the disk for any emitting point depends only on its height above (or below) the disk midplane and the angle $`\theta _{occ}`$. The observer’s line-of-sight to the emitting point will pass through a range of heights above the disk plane, so the total optical depth must be found by integrating along this path. The total optical depth is calculated using:
$$\tau (z_0)=0.5(\tau _{mid}/\mathrm{sin}\theta _{occ})_{z_0}^{\mathrm{}}\mathrm{exp}(\frac{z}{\sigma _d\mathrm{cos}\theta _{occ}})^2(\frac{dz}{\mathrm{cos}\theta _{occ}})$$
(5)
where $`z_0`$ is the height above the disk midplane where the ray connecting the observer and the emitting point intersects the inner disk edge. The extinction to the emitting point is then given by:
$$A(z_0)=e^{\tau (z_0)}$$
(6)
The effect of the twisted disk is taken into account by having the optical depth increase to infinity as $`z_0`$ becomes increasingly negative (see figure 2).
The geometric orientation of the neutron star and the inner disk are determined by tilt and azimuth angles. The tilt is specified with respect to an axis that will be identified as (but is not required to be) the stellar binary axis. Reasonable values for the disk tilt lie in the range of $`1020^{}`$ in order to fit the overall High state light curve with a tilted, twisted, counter-precessing disk (see figure 2). If one assumes the disk edge is cutting across the face of the neutron star at 35-day phases 0.23 and 0.58 (from observing the pulse evolution) then a tilt of $`11^{}`$ can be derived.
### 7.1 Comparison with Observations
#### 7.1.1 Direct Fan Beam
Can this simple model reproduce, qualitatively, the features observed in the Main High state pulse evolution? Figure 8 (top panel) shows the approximate spatial locations of the pencil beam and fan beam emission regions for the direct fan beam model. The neutron star rotation is prograde in figure 8 so a counter-precessing inner disk edge will sweep across the neutron star face from right-to-left in the figure. We define case 1 disk orientation as illustrated in the top panel of figure 8, that is, the disk covers the neutron star from top-to-bottom as well as right-to-left. Case 2 disk orientation is defined as the case that the disk covers the neutron star face from right-to-left, bottom-to-top, as illustrated in the lower panel of figure 8. A straightforward consideration of the twisted disk geometry shows that during a full 35-day precession period, case 1 and case 2 both occur, separated by one-half a precession period.
The relatively long period during the Main High state when the pulse profile shows only small changes, simply means that the line-of-sight to the pulsar is far from the obscuring inner disk edge. As the disk edge approaches the pulsar, it is clear from figure 8 (top panel), that in both case 1 and case 2, the trailing shoulder (Bb) of the main pulse will disappear before the pencil beam (C) or the leading shoulder (Ba). In case 2, the interpulse (A) will disappear before the pencil beam (C) and in case 1 they disappear at nearly the same time. Neither case reproduces the pulse evolution behavior observed in either High state.
Reorienting the neutron star spin axis about the line-of-sight will alter the occultation sequences. If the upper spin pole in figure 8 is tipped toward the right (toward the oncoming disk edge) then case 1 will produce a pulse profile sequence in which the order of disappearance will be 1) the leading and trailing soft shoulders (Ba and Bb), 2) the hard central peak (C), 3) the interpulse (A). Case 1 does not resemble the observed Main High state pulse evolution sequence. But it can can reproduce the Short High state sequence if the emergence from the turn-on only reveals the occultation after the disappearance of the two soft shoulders. If case 1 gives the Short High state sequence then case 2 must give the Main High state sequence. However, in case 2 the trailing shoulder (Bb) disappears first, then the hard central peak (C) will disappear, both before the leading shoulder (Ba), contrary to the observations.
One can consider all the different possible combinations of neutron star spin and motions of the disk edge. This has been done and the results (of whether a satisfactory occultation sequence is obtained) are summarized in Table 2. Only two of the eight possibilities result in occultation sequences in the correct order for both the Main and Short High states. One requires a reverse neutron star spin in one case and a prograde disk precession in the other case. Is a retrograde rotation of the neutron star or a prograde precession of the disk possible? The long term spinup trend observed in Her X-1 (Nagase 1989) and the frequency behavior of optical pulsations from the lobes of the companion star HZ Her (Middleditch & Nelson 1976) both strongly argue for a prograde pulsar spin. Likewise, the preeclipse dip recurrence period and the predominance of integer combinations of the sum of the orbital and 35-day frequency in the power density spectrum of the optical light curve strongly indicate a counter-precession for the accretion disk (Deeter et al. 1976; Crosa and Boynton 1980). Thus we discard the direct fan beam model, since it cannot give the correct occultation sequence and keep prograde spin and retrograde precession.
We also note that those occultation sequences from the direct fan beam model, even though in the correct order, have a difficulty. For the Short High state the observed sequence starts with a weak hard pulse (C) but has no evidence whatsoever for any trace of the leading (Ba) or trailing (Bb) shoulders. This is very hard to achieve in the direct fan beam model since the hard pulse (C) and soft shoulders (Ba, Bb) are produced in directly adjacent locations. The above qualitative difficulties in matching the observed sequence of pulse profile changes can be resolved using a “reversed” fan beam, as we show next.
#### 7.1.2 Reversed Fan Beam
Both the observed Main and Short High state pulse evolution can be reproduced qualitatively if the neutron star is tilted as shown in the lower panels of figure 8. Case 1 (disk tilted as in the upper panel) reproduces the Main High state decay of the leading soft shoulder (Ba) followed by the hard central peak (C) and the trailing soft shoulder (Bb). The Short High state evolution is produced by the case 2 disk orientation as shown in the lower panel of figure 8. The observed evolution requires complete occultation of the soft shoulders (Ba, Bb) of the main pulse before the Short High state turn-on. This can be readily accomplished since the shoulders (Ba, Bb) are from regions well separated from the location of the hard central peak (C). The shorter length of the Short High state and the required initial partial occultation of the pulse emitting region can be produced by the same offset of the observer’s line-of-sight from the binary plane. In figure 10a and 10b, a disk occultation and the resulting pulse profile changes are illustrated using the model disk and beam profiles described earlier. The parameters used are displayed in table 1.
The Main High state pulse evolution is modeled in figure 10a. The leading edge decay of both the fan beam and the interpulse occurs as the disk edge sweeps across the neutron star face. The actual timing of the decay of the soft shoulders and the interpulse put important constraints on the apparent tilt of the neutron star relative to the disk edge, which is also quite sensitive to the actual locations of the emitting regions. For the model shown in figure 10a, the interpulse decays away a bit early relative to the soft shoulders compared to the observations. In addition the decay of the trailing soft shoulder is delayed compared to the observations since the model disk edge appears to be too sharp. However, making the disk edge fuzzier will tend to wash out the Short High state sequence, assuming a fixed inner disk ring tilt and radius. This model predicts that the pencil beam (hard central peak, C), should show a trailing edge decay, but observing this effect requires a known pulse ephemeris during this decay phase. Lastly, as the disk edge cuts across the neutron star face, short term variations in the pulse profile should occur in addition to the longer term systematic changes, due to variations in the disk opacity as material is accreted onto the neutron star.
The Short High state pulse evolution is modeled in figure 10b. The disk motion will cause the pencil beam (small hard peak, C) to decay away first followed by the interpulse (A). The more edge-on view of the inner disk plane at the start of the Short High state will lengthen the traversal of the disk edge across the neutron star face relative to the Main High state and this is consistent with the slow disappearance of the interpulse (soft peak, A) relative to the same event during the Main High state. In contrast to the Main High state, no significant leading edge decay of the pulse components is predicted during the Short High state and none is observed. The Short High state must begin with the accretion disk partially obscuring the pulse emission region to account for both the different pulse shape and the lower Short High state flux. The placement of the inner disk edge at the turn-on is probably somewhat variable due to changes in the total disk twist and elevation so this model predicts that the Short High state pulse profile should be quite variable as well at the start of the turn-on. Depending on the exact placement of the inner disk edge, the small hard peak (C) may be larger, smaller, or completely attenuated compared to the interpulse in figure 6. Likewise, the disk edge may be advanced enough that the interpulse is gone as well at the turn-on.
The evolution of the quasi-sinusoidal components is modeled by the occultation of emission regions D and E in figure 10. The deepening of the pre-interpulse minima early in the Main High state and the persistence of the quasi-sinusoidal pulse after the disappearance of the main pulse suggests an occultation of an emission region much larger in size than the neutron star. Radiation scattered from the inflowing magnetospheric material many neutron star radii above the surface may produce this emission. If we assume 1) the magnetospheric flow is optically thin to Thomson scattering as in the Brainerd & Mészáros model and 2) that accretion onto a magnetic pole occurs over a range of magnetic azimuth less than $`180^{}`$, then two equal and constant flux pulse components may be produced by scattering in the accretion flows. The locus of scattering points illuminated by the hard pencil beam (C) over a neutron star rotation period are indicated by D and E in the rightmost panels of figure 8. D and E are much further away from the neutron star than the emission components A, Ba, Bb and C (middle column of panels in figure 8): the small circle between D and E indicates the neutron star surface. Since the emission from E is due to the magnetospheric flow onto the the opposite pole on the neutron star from the flow producing the emission at D, the instantaneous location of scattered emission on E will be $`180^{}`$ different in pulse phase from the instantaneous location of the emission on D. D and E are still close enough to the neutron star that light travel time delays are negligible. In the absence of any occultation, the instantaneous scattering from D and E are visible at all pulse phases so the total emission from D and E appears unmodulated. A third, small, constant-flux contribution should also be present from the much larger accretion disk corona that produces the Low state flux.
The Main and Short High state quasi-sinusoidal pulse evolution and the preinterpulse deepening can together be interpreted as an occultation of the X-ray illuminated, magnetospheric flow (regions D and E). Both the upper and lower magnetospheric flows will sweep out cones as the neutron star rotates, which cross the “plane of the sky” at pulse phases 0.25 and 0.75. For the neutron star orientation show in figure 8 (lower panels), a deepening is predicted in the pulse profile during the Main High state near pulse phase 0.25 as the outer portion of the upper magnetospheric flow rotates behind the oncoming disk edge. The widening will continue as the occultation progresses but the deepening will stop since the flux at pulse phase 0.25 will then be produced mostly by the unocculted, antipodal, magnetospheric flow.
Near the middle of the Main High state, a second deepening is predicted at pulse phase 0.75 as the disk begins to occult the lower magnetospheric flow. At the Short High state turn-on, the lower magnetospheric flow will already be occulted, accounting for the $`50\%`$ drop in flux observed in the “quasi-sinusoid”. The constant flux produced by the upper magnetospheric flow will now be preferentially occulted near pulse phase 0.75 as the occultation progresses, producing a dip there, and an $`180^{}`$ phase shift between the two High state “quasi-sinusoids”. The dip should widen as the Short High state progresses. The initial increase and then decrease of the overall quasi-sinusoidal flux occurs simply because the line-of-sight is moving away from the outer disk edge and through decreasing accretion disk coronal density as the inner disk occultation progresses.
A comparison of the model in figure 10 with the observations readily shows that while many features of the quasi-sinusoidal evolution can be modeled, there is a problem with component E. At the end of the Main High state, the quasi-sinusoid peaks at pulse phase 0.75, whereas component D is near minima, implying that it is out-of-phase by 0.5. The Short High state provides no constraint on E since it is completely occulted. The “fix” to the simple model needed to produce the asymmetry in the soft shoulder peaks should also affect the location of pulse component E and may solve this problem. In any case, using two rotating rings to model the quasi-sinusoids is probably a gross simplification of the actual situation since the interaction geometry of the accretion column with the disk should be quite complex. In addition, another weakly pulsed component may be present due to scattering or reemission from the inner disk edge. This may be the cause of the low energy sinusoid that appears to be in phase with the higher energy quasi-sinusoid (see e.g. McCray et al. 1982). We note that the pulse evolution of the soft energy quasi-sinusoid ($`<1`$ kev) is at present unknown and that its observation may provide valuable clues about the inner disk. Recent BeppoSAX observations of a Short High state flux decline (Oosterbroek et al. 2000) discovered an increase in relative absorption that can be explained by assuming seperate scattering and absorption regions. This is in qualitative agreement with a gradual inner disk occultation of an extended scattering region associated with the accretion column.
The radius of the inner disk edge can be estimated from the duration of the Main High state occultation event. Let $`R_d`$ be the inner disk radius and $`R_e`$ the radius of the pulse emitting region. From the reversed beam model $`R_e`$ is about $`2R_{ns}`$ for the region emitting the main pulse. The velocity of the disk edge is given approximately by $`V_d=R_d\omega _d\mathrm{sin}\theta _t`$ where $`\omega _d`$ is the angular velocity of the disk precession (equal to $`\frac{2\pi }{P_{35}}`$) and $`\theta _t`$ is the tilt angle of the inner disk. The inner disk radius can now be estimated from:
$$R_d=\frac{2R_e}{T_{occ}\omega _d\mathrm{sin}\theta _t}$$
where $`T_{occ}`$ is the duration of the occultation of the main pulse emitting region. From the April 1989 Main High state, $`T_{occ}3`$ days. The inner disk tilt angle is estimated to be between $`1020^{}`$. These parameter values produce an estimate of $`R_d2040R_{ns}`$. This is much smaller than the corotation radius of $`157R_{ns}`$ where the orbital angular velocity equals that of the neutron star, assuming a 12 km neutron star radius and $`M_{NS}=1.3\mathrm{M}_{}`$.
In conclusion, the reverse fan beam model with prograde neutron star spin and retrograde disk precession naturally accounts for the major features of the observed Main and Short High state pulse evolution. It also accounts for a number of other features, as described above, e.g. the $`180^{}`$ phase shift between the two High state “quasi-sinusoids”.
## 8 Discussion
The model proposed here for the pulse evolution cycle in Her X-1 consists of an occultation of the neutron star by the inner edge of a tilted and precessing disk. The leading edge decay of the main and interpulses during the Main High state preclude a beam geometry consisting of a pencil beam surrounded by a fan beam and emitting from the neutron star surface. However, reversing the fan beam so that is emitted in the antipodal direction with respect to the pencil beam and at some distance above the stellar surface allows the leading edge decay to be reproduced by an occultation in a natural fashion. The observed Short High state evolution pattern then arises naturally with this geometry.
Cyclotron resonant scattering in the accretion column is an attractive mechanism for producing a pencil beam surrounded by an reversed fan beam. The energy of the cyclotron resonance is directly proportional to the local magnetic field strength and this decreases with altitude above the neutron star surface. When the energy of an upward traveling photon equals the local cyclotron energy, scattering will occur. This creates a natural energy dependent filtering process for photons emitted in the accretion cap and divides the beam into three components. Hard photons (especially those above the surface cyclotron frequency) will escape in a pencil beam. Softer photons will be backscattered and gravitationally focused around the neutron star in an antipodally directed fan beam. Finally, photons scattered from the accretion column isotropically will produce a constant pulse component. The softest photons will scatter from the highest altitudes in the accretion column. The neutron star rotation will cause the highest altitude emission (D and E in figure 8) to rotate with pulse phase. The occultation of this high altitude emission by the precessing disk edge creates gaps in the constant emission profile and therefore a “quasi-sinusoidal” profile. A hard pencil beam (C), a softer fan beam (Ba and Bb) and a quasi-sinusoidal component (D and E) qualitatively match the observed energy dependence in the Main High state profile. This type of beam model may be applicable to other X-ray pulsars as well, for example Vela X-1, which displays a pulse profile consisting of two doubled peaked components at soft X-ray energies, possibly superposed on a quasi-sinusoidal component, that fills in at harder X-rays in a way quite similar to the Her X-1 main pulse (White, Swank & Holt 1983). In the case of Vela X-1, the observer’s line-of-sight would have to be located much closer to the neutron star spin equator than in Her X-1.
While an occultation model apparently has many attractive features, a roughly $`45^{}`$ tilt is required between the neutron star spin axis and the binary axis of the system. Is this plausible? A large tilt to the neutron star may have originated in the supernova explosion that gave it birth. The accretion of matter from the companion will cause angular momentum to be accreted by the neutron star. Since the direction of the time averaged accreted angular momentum is along the binary axis of the system, the neutron star’s spin axis should become coaligned with the binary axis. This event will take some time and the accreted angular momentum will spin up the neutron star. In fact, the alignment timescale and the spin-up timescale should be comparable. The measured spin-up time of Her X-1 is about $`10^5`$ years. This timescale is the same as that predicted for the entire X-ray emitting phase of Her X-1 (Savonije 1978). Therefore, if Her X-1 is currently tilted that tilt will in all likelyhood be maintained throughout the rest of the X-ray emitting phase. However the important question is what is the current ratio of accreted angular momentum to that at the start of the X-ray emitting phase?
Her X-1 may be near the start of the X-ray emitting phase. The X-ray emitting phase ends when the accretion flow becomes great enough to smother the neutron star and prevent pulsations. Historical optical observations of HZ Her show that the X-ray heating and hence the mass accretion rate has ceased occasionally for years to decade-long periods over the last hundred years (Jones, Forman & Liller 1973; Hudec & Wenzel 1986). The current state of mass transfer depends on the X-ray heating of HZ Her and is inhibited by the X-radiation pressure. Thus HZ Her must be close to, but not quite, filling its Roche lobe. This state of affairs is more consistent with Her X-1 being at the beginning rather than the end of its X-ray phase. If so, then a highly tilted neutron star can plausibly exist in the Her X-1 system. Another point in favor of such an interpretation is the 1.24 second pulse period of Her X-1. This period is typical of the radio pulsar population. A period longer than 3 seconds or much shorter than 1 second would unequivocally show that Her X-1 has been spun down or spun up by a significant history of interaction with circumstellar matter. Finally, a tilted neutron star should cause a persistent asymmetry in the optical orbital lightcurve about orbital phase 0.5. No such asymmetry was reported by Deeter et al. (1976) but later studies of the orbital optical lightcurve with more data indicate the presence of just such an asymmetry (e.g. Voloshina, Lyuti & Sheffer 1990; Thomas et al. 1983).
The properties required of the inner disk to fit the observed pulse evolution are a scale height comparable to the neutron star diameter and a small inner radius ($`2040R_{ns}`$). For comparison, the corotation radius is $`157R_{ns}`$ for Her X-1. The predicted magnetospheric radius, $`R_m`$, for disk accretion depends on assumptions on the boundary conditions. Using the known parameters for Her X-1 (Leahy & Scott 1998), one obtains $`R_m=3.5\times 10^8\alpha ^{2/61}`$ cm for the model of Kiraly & Mészáros (1988); $`R_m=3.8\times 10^8\gamma ^{2/7}`$ cm for the model of Lamb (1988); and $`R_m=4.9\times 10^8`$ K cm for the formula of Finger et al. (1996) (which has the special case of $`K=0.47`$ for the model of Ghosh & Lamb, 1978). The model of Aly (1980) for a highly conducting disk, gives only a slightly smaller value for the disk inner radius: $`R_m=3.0\times 10^8\alpha ^{2/7}\mathrm{cos}(\chi )^{4/7}`$ cm (for $`\chi `$ near $`\pi /2`$, the dependence on $`\chi `$ is $`\mathrm{sin}(\chi )^{40/69}`$, so the minimum $`R_m`$ is $`1/3`$ of this for $`\chi =0`$). The disk viscosity parameter is $`\alpha `$, $`\gamma `$ is defined in Lamb (1988), K is a dimensionless parameter, which Finger et al. (1996) find to be $`1`$ for A0535+26, and $`\chi `$ is the tilt of the dipole axis from the equatorial plane. In all cases the predicted disk radius greatly exceeds the inner disk radius required by the occultation model.
Another observation relevant to the magnetospheric radius is the spinup rate of Her X- 1. It is the smallest among the accretion-powered pulsars and indicates the slowest rate of net angular momentum accretion. During the giant outbursts of Be-transients, an X-ray flux vs. spin-up rate correlation has been observed in which spin-up occurs at a rate consistent with the fiducial torque (e.g. Finger et al. (1996); Bildsten et al. (1997); Nelson et al. (1997); Scott et al. (1997)). The formula for the fiducial accretion torque $`N_f`$ is given by $`N_f=\dot{M}(GM_XR_s)^{1/2}`$ where $`R_s`$ is the magnetospheric radius for spin-up, $`\dot{M}`$ is the mass accretion rate and $`M_X`$ is the mass of the neutron star. This formula assumes all the angular momentum of the accreting matter at $`R_s`$ is given to the neutron star. The known spin-up rate gives a value for $`R_s=8.5`$ km: clearly too small to be physical. The small disk radius from the occultation model is much larger than this. This implies that the actual torque on the neutron star must be smaller than the fiducial torque, to allow a magnetospheric radius as large as $`2040R_{ns}`$. In the Ghosh & Lamb (1978, 1979) model of an aligned rotator, a residual portion of the magnetic field is not screened by currents near the magnetospheric radius and will interact with material orbiting beyond the corotation radius. The interaction adds a negative torque component to the total torque. A smaller net torque will be exerted if the magnetospheric radius approaches the corotation radius<sup>3</sup><sup>3</sup>3Note that Ghosh and Lamb prefer to call the corotation radius the ‘centrifugal radius’ and the magnetospheric radius the ‘corotation radius’. of the neutron star. Thus nearly any spin-up rate can be obtained within a small range of magnetospheric radii near to but inside the corotation radius. So it is not surprising that the predicted magnetospheric radius from the Ghosh and Lamb model for Her X-1 is nearly equal to the corotation radius given above. However, it is well known that the Ghosh and Lamb model is not consistent with the spin behaviour of many X-ray pulsars (e.g. Bildsten et al. 1997; Nelson et al. 1997). In summary, different models give a magnetospheric radius from the spin-up rate of Her X-1 in a wide range between $`R_{ns}`$ and $`R_c=157R_{ns}`$.
How does one account for an inner disk edge which is at $`2040R_{ns}`$? It is much smaller than the predicted neutron star magnetosphere radius, yet may be much larger than the radius predicted from simple spin-up. We speculate on what may be the real physical origin of the $`2040R_{ns}`$ disk inner edge. One possibility is that the dipole magnetic field is much smaller than that indicated by the observed X-ray cyclotron line so that the predicted magnetospheric radius is much smaller. The magnetospheric radius decreases from $`200R_{ns}`$ to $`30R_{ns}`$ if the dipole component of B is a factor of 25 less than deduced from the cyclotron line. The cyclotron line would be explained as arising in non-dipolar magnetic fields in the accretion region at the neutron star surface. An emission region concentrated in a non-dipolar surface pocket of the field and producing a pencil beam would likely be difficult to distinguish from a similar emission region in the case of a pure dipole. Another explanation has been put forth by Baushev & Bisnovatyi-Kogan (1999) in which a magnetic field of $`46\times 10^{10}`$ G is estimated for Her X-1 from the observed cyclotron line energy by assuming a large anisotropy exists in the electron momentum distributions parallel and perpendicular to the magnetic field lines. Whether either of these possibilities is viable will probably require much more research.
An alternative to a reduced dipole field is a stable, thin inner disk which penetrates far into the magnetosphere. However, this seems very unlikely: the magnetosphere is very stiff due to the steep gradient of magnetic pressure and energy density (as $`r^6`$) for a dipole field. Models which have the disk penetrate as deeply as possible (e.g. Aly 1980), do not have an inner radius very much less than other models (tilting the magnetic axis gives a reduction at most by a factor 3 in the inner disk radius). So a disk penetrating far into the magnetosphere appears to be unfeasible.
Another possibility is that current models for determining the magnetospheric radius are inadequate. For example, Miller & Stone (1997) use magnetohydrodynamic calculations to show that the Balbus-Hawley instability and magnetic braking have dramatic effects on the magnetospheric boundary. They also result in outflowing winds along field lines opened up by reconnection. It is quite possible that the small inner disk in Her X-1 is due to such effects. We note that mass outflows are quite likely as there is extended X-ray emission from a large corona in Her X-1.
The model described in this paper is phenomenological and highly idealized. The purpose was to show that an inner disk occultation can explain the observations and has reasonable physical grounds for support. Considerable refinement of the model remains to be done. The decay of the soft shoulders of the main pulse occurs earlier in the Main High state than an axisymmetric beam model predicts. The soft shoulders of the Main High state main pulse are unequal in amplitude. These problems might indicate an offset in the dipole axis.
An occultation of the neutron star by the inner accretion disk explains the pulse evolution cycle of Her X-1 in a natural fashion. Many of the details of the observed evolution can be accounted for by invoking a reversed fan beam geometry around a neutron star significantly inclined to the binary axis of the system. The leading edge decay of the main pulse and interpulse during the Main High state is properly predicted. During the Short High state, the rapid disappearance of the small hard peak as well as the slower decay of the soft peak are predicted. Note that the soft peak shows little decay of either the leading or trailing edge and this is also predicted by the occultation model. In summary, most of the curious pulse shape changes observed during the Main and Short High states are tied together with a single simple occultation model. We note that no previous model for the pulse shape changes, including free precession, has been able to account for the observed details of the pulse evolution.
DMS acknowledges John E. Deeter and Paul E. Boynton for useful discussions and Fumiaki Nagase and Takashi Aoki for assistance while visiting ISAS.
Fig. 1—Light curves (1–37 keV) for Ginga LAC observations of Her X-1 in the April, May and June 1989 Main, Short and Main High states, respectively. The 35-Day phase is calculated using $`P_{35}=34.8534`$ days ($`20.5P_{orb}`$) and an epoch $`T_0=48478.631P_{35}=47398.14`$ MJD with phase 0.0 corresponding to the Main High state turn-on. Eclipse ingress and egress are denoted by dashed vertical lines. Ticks and numbers at the top of each panel indicate the 35-day phase. Solid vertical lines mark predicted turn-on time. The bottommost plot also shows the 20-50 keV pulsed flux lightcurve obtained from BATSE monitoring of Her X-1, with the flux scale shown on the righthand side.
Fig. 2—(upper panel) View of the disk as seen from the neutron star. The outermost ring (filled diamonds) is tilted by $`20^{}`$, the innermost ring (hollow diamonds) by $`11^{}`$ and leads in precession phase by $`138.6^{}`$. The elevation of the observer of Her X-1 ($`5^{}`$) is marked by the horizontal dashed line. Bold solid vertical lines mark turn-on’s. The observer sees Her X-1 emerge as a point source from the outer disk rim but an extended source at the inner disk rim. The vertical dotted line shows the approximate point where the disk cuts across the neutron star face during the Main High state. (Middle panel) The 1989 1–37 keV Her X-1 light curve observed with Ginga. (Bottom panel) Nineteen “0.2 turn-on” 35-day cycles of the 2–12 keV RXTE ASM light curve folded at a period of 20.5 $`P_{orb}`$.
Fig. 3—(upper left panel) Ginga observation of 1-37 keV count rate during the May 1989 Short High state turn-on. Vertical dashed line marks predicted eclipse egress. Solid curve models flux of point source rising through outer disk edge (see text). Errors are plotted, unless smaller than plotted point size. (lower left panel) Same for August 1991 Main High state turn-on. (upper right panel) Softness ratios of two sets of energies during the May 1989 Short High state turn-on. Note that harder flux increases first consistent with cold matter absorption. (lower right panel) Same for August 1991 Main High state turn-on.
Fig. 4—The softness ratio for the average Main High state light curve. The hard color is the BATSE 20-50 kev pulsed flux and the soft color is the RXTE/ASM 2-12 keV flux. The top panel shows the average softness ratio for 0.2 turn-on Main High states, the bottom panel for 0.7 turn-on Main High states. Turn-on is 35-day phase 0. Vertical dashed lines mark eclipse ingress/egress boundaries.
Fig. 5—Sample Ginga pulse profiles in five energy bands, after subtracting background and correcting for collimator transmission. The bands increase in energy from top to bottom: 1.0–4.6, 4.6–9.3, 9.3–14, 14–23, 23–37 keV. The bottommost panels display a hardness ratio (9.3–23 keV band divided by 1.0–4.6 keV band). Pulse features discussed in the text are labeled. (a)—Leftmost panel set. Main High state observation on MJD 47643. Total exposure time is 8713 seconds. Main pulse occupies phase interval 0.75–1.25 and the interpulse occupies phase interval 0.3–0.7. (b)—Center panel set. Same profiles as in (a), but with close-up of interpulse. (c)—Rightmost panel set. Short High state observation on MJD 47662, less than a day after Short High state turn-on. Total exposure time is 10118 seconds. At same scale as center panel set.
Fig. 6—Time evolution of the pulse profile. Upper right and upper left panels show evolution during the April 1989 Main High state in energy bands 1.0-4.6 keV and 9.3-14 keV. The flux of the smallest amplitude pulse in each panel is correct while offsets of 50 counts/sec have been added between pulses for clarity. The 35d phase of the pulses increase with decreasing amplitude as: 0.05, 0.162, 0.216, 0.243, 0.245, 0.247, 0.249, 0.250. Lower left and lower right panels. Same for the May 1989 Short High state but with offsets of 20 counts/sec added. The 35d phase increases with decreasing amplitude as 0.590, 0.595, 0.613, 0.641, 0.698.
Fig. 7—Sequence of pulse profiles resulting from a freely precessing neutron star with an axisymmetric pencil and fan beam using the parameters determined by Kahabka (1987,1989). The pencil beam amplitude is three times the fan beam amplitude and both beam components have half-widths of about $`20^{}`$. One full 35-day precession cycle is shown with successive pulses occuring one day apart and spaced by one flux unit. On the right is the pulsed flux light curve over the precession cycle normalized to a maximum value of 10.
Fig. 8—Top left panel. Model pulse comprised of identical beam configurations at each pole consisting of a central pencil beam (C) and surrounding by a fan beam (B) with an opening angle of $`40^{}`$. Three constant components are present, two produced by isotropic emission high in the accretion column (D,E) and one from coronal emission. Top center. Emission locations of pulse components A, B and C. Location of disk edge for 35-day phase 0.23 shown moving from right-to-left, top-to-bottom across figure (case 1). Top right. Larger view showing emission locations of pulse components D and E. Bottom left. Same as above but with reversed fan beam (opening angle of $`140^{}`$) and light bending taken into account. Bottom center. The ‘A’ and ‘B’ components are now emitted when the accretion column is on the “back” side of the neutron star with respect to the observer. Location of disk edge for 35-day phase 0.58 shown moving from right-to-left, bottom-to-top across figure (case 2).
Fig. 9—Top panel. Photon trajectories for $`10`$ keV photons backscattered from a cyclotron resonance at a height of $`1.5R_{ns}`$ above a $`1.4\mathrm{M}_{}`$ neutron star. The scattering height assumes a surface magnetic field strength corresponding to a $`40`$ keV cyclotron line and a simple dipole field. Bottom panel. The photon impact parameter and apparent emission angle for the case of no light bending (diamonds) and with light bending (triangles). Distance units in Schwarzschild radii.
Fig. 10a—Model of Main High state pulse evolution. Left panel sets shows a sequence of pulse profiles corresponding to the figure 8b, reversed fan beam case. Middle panel set shows closeup of pulse emission region and gradual covering by disk. Rightmost panel set shows larger view of same region. 35-day phase is shown at far right.
Fig. 10b—Model of Short High state pulse evolution. Left panel sets shows a sequence of pulse profiles corresponding to the figure 8b, reversed fan beam case. Middle panel set shows closeup of pulse emission region and gradual covering by disk. Rightmost panel set shows larger view of same region. 35-day phase is shown at far right. |
warning/0002/hep-ph0002140.html | ar5iv | text | # Measurement of entropy of a multiparticle system: a do-list
## Abstract
An algorithm for measurement of entropy in multiparticle systems, based on the recently published proposal of the present authors , is given. Dependence on discretization of the system and effects of multiparticle correlations are discussed in some detail.
It was suggested recently that studying event-by-event fluctuations may be used for determination of entropy of multiparticle systems created in high-energy collisions. A generalization of this idea and a specific proposal for measurement of entropy were formulated in . In the present note we spell out explicitly the steps to be taken to implement effectively the method proposed in . Importance of the dependence of measurements on discretization of particle momenta and the role of (multi)particle correlations are emphasized.
1. Selection of the phase-space region.
As the first step in the process of measurement one has to select a phase-space region in which measurements are to be performed. This of course depends on the detector acceptance as well as on the physics one wants to investigate. The region cannot be too large because for large systems the method is difficult to apply (the requirements on statistics become too demanding). With the statistics of $`10^6`$ events, the region containing (on the average) $`100`$ or less particles should be possible to investigate. A reasonable procedure seems to be to start from a small region and then increase it until the errors become unbearable.
Comment 1: The proposed measurement is not restricted to systems with very large number of particles. It can be applied to any multiparticle system, e.g., in $`e^+e^{}`$ annihilation, hadron-hadron collisions or peripheral nucleus-nucleus collisions.
2. Discretization of the spectrum.
The selected phase-space region should now be divided into bins of equal size in momentum space. The number of bins cannot be too large if one wants to keep errors under control. On the other hand, as argued below, it is important to study the dependence of results on the size (and thus the number) of the bins. Therefore, large statistics is essential for a success of the measurement.
Comment 2: If one chooses the bins which are not of equal size in momentum space, the original expression for entropy requires a correction which follows from an appropriate change of variables . This correction is, in general, not easy to calculate. Nevertheless it may be interesting to study the dependence on the shape of the binning, as well.
3. Description of an event.
Using this procedure, an event is characterized by the number of particles in each bin, i.e. by a set of integer numbers $`sm_i^{(j)}`$, where $`i=1,\mathrm{}M`$ ($`M`$ is the total number of bins) and the superscript $`(j)`$ runs over all kinds of particles present in the final state. These sets represent different states of the multiparticle system which were realized in a given experiment. The number of possible different sets is, generally, very large (for 5 bins and 100 particles one obtains $`10^6`$ sets). This is, in fact, the main difficulty in application of the proposed method. It simply reflects the fact that the system we are dealing with has very many states.
Comment 3: It should be realized that, in practice, such a description is never complete, i.e., it never describes fully the event. Most often some of the variables are summed over. This is the case, e.g., when one measures only charged particles. Then all the variables (i.e. multiplicities and momenta) related to neutral particles are summed over. It may be thus interesting to study reduced events, when even some of the measured variables (e.g. particle identity) are summed over (i.e. ignored).
4. Measurement of coincidence probabilities.
As explained in , this is the basis of the method and therefore the most important step in the whole procedure.
The measurement consists of the simple counting how many times ($`n_s`$) any given set $`s`$ appears in the whole sample of events<sup>1</sup><sup>1</sup>1 Since the number of different sets is very large, most of them shall appear only once or not at all.. Once the numbers $`n_s`$ are known for all sets, one forms the sums:
$$N_k=\underset{s}{}n_s(n_s1)\mathrm{}(n_sk+1),ł1$$
(1)
with $`k=1,2,3,\mathrm{}`$. The sum formally runs over all sets s recorded in a given experiment, but nonvanishing contributions give only those which were recorded at least $`k`$ times. One sees that $`N_k`$ is the total number of observed coincidences of $`k`$ configurations. The coincidence probability of $`k`$ configurations is thus given by
$$C_k=\frac{N_k}{N(N1)\mathrm{}(Nk+1)},ł2$$
(2)
where $`N`$ is the total number of the events in the sample<sup>2</sup><sup>2</sup>2As explained in this ratio is equal to the $`(k1)`$-th moment of the probability distribution $`C_k=_s(p_s)^k`$. The proof follows closely the argument of ..
One sees that $`N_k`$ given by (LABEL:1) are simply factorial moments of the distribution of $`n_s`$ . It is also clear that, since $`_sn_s=N`$, $`C_1=1`$. Finally, one sees that only states with $`n_sk`$ contribute to $`N_k`$ (and thus also to $`C_k`$).
5. Errors.
The error of $`C_k`$ is determined by the error of the numerator in (LABEL:2). This error can be estimated by standard methods used in evaluation of the moments of a distribution.
6. Renyi entropies and Shannon entropy.
Once the coincidence probabilities $`C_k`$ $`(k=1,2,\mathrm{})`$ are measured, it is convenient to calculate the Renyi entropies defined by
$$H_k\frac{\mathrm{log}C_k}{k1}.ł4$$
(3)
The Shannon entropy $`S`$ (i.e. the standard statistical entropy) is formally equal to the limit of $`H_k`$ as $`k1`$ and thus can only be obtained by extrapolation from a series of measured values: $`H_k=H_2,H_3,\mathrm{}.`$ to $`k=1`$ <sup>3</sup><sup>3</sup>3 Obviously, one cannot just put $`k=1`$ in the formula (LABEL:4) for that purpose: since $`C_1=1`$, the R.H.S. of (LABEL:4) for $`k=1`$ represents the undefined symbol $`0/0`$. . Of course such an extrapolation procedure is not unique and introduces uncertainty. The main point is, as usual, to choose the ”best” extrapolation formula, i.e. the functional dependence of $`H_k`$ on $`k`$ which will be used to reach the point $`k=1`$ from the measured points $`k=2,3,\mathrm{}`$. This form can only be guessed on the basis of physics arguments (or prejudicies).
In it was suggested to use
$$H_k=a\frac{\mathrm{log}k}{k1}+a_0+a_1(k1)+a_2(k1)^2+\mathrm{}.,ł5$$
(4)
where the number of terms is determined by the number of measured Renyi entropies. This formula turned out to be very effective in reproducing the correct value of entropy for some typical distributions encountered in high-energy collisions.
Another possibility is to use
$$H_k=a_0+a_1/k+a_2/k^2+\mathrm{}.ł6$$
(5)
suggested by the formula for the free gas of massless bosons<sup>4</sup><sup>4</sup>4For the free gas of massless bosons the Renyi entropies are given by $`H_k=(1+1/k+1/k^2+1/k^3)S/4`$ where $`S`$ is the Shannon entropy..
It will be interesting to compare the results from these two formulae.
Comment 4: The measured values of the Renyi entropies give valuable information about the system and thus are of great interest, independently of the accuracy of the extrapolation.
7. Dependence on discretization; Scaling.
As the result of the procedure explained in Sections 1 to 6, we obtain the Renyi entropies $`H_k`$, (k=2,3,…) and the Shannon entropy $`S`$ of a given phase-space region. These entropies still depend on the method of discretization of the momentum spectrum, in particular on the size of the binning. If the bins are small enough and if the system is close to thermal equilibrium (i.e. if fluctuations are small), one expects the following scaling law to hold
$$H_k(lM)=H_k(M)+\mathrm{log}lS(lM)=S(M)+\mathrm{log}lł7$$
(6)
($`M`$ and $`lM`$ are numbers of bins in two different discretizations). If the scaling law is verified, one can determine the part of entropy which is independent of binning.
The rule (LABEL:7) is not expected to hold if the system is far from thermal equilibrium and the fluctuations of the particle distribution are large. In particular, the effects of intermittency and erraticity as implied, e.g., by a cascading mechanism of particle production are expected to violate (LABEL:7). Thus testing the dependence of entropies on the number of bins may reveal interesting features of the system.
8. Comparison of different regions; Additivity.
Measurements of the entropies $`H_k`$ and $`S`$, as described above, can be performed independently (and - in fact- simultaneously) in different phase-space regions. The results should give information on the entropy density and its possible dependence on the position in phase-space (e.g., it seems likely that the results in the central rapidity region may be rather different from those in the projectile or target fragmentation). Furthermore, it is important to verify to what an extent the obtained entropies are additive, i.e., whether the entropies measured in a region $`R`$ which is the sum of two regions $`R_1`$ and $`R_2`$ satisfy
$$H_k(R)=H_k(R_1)+H_k(R_2)S(R)=S(R_1)+S(R_2).ł8$$
(7)
Eq. (LABEL:8) should be satisfied if there are no strong correlations between the particles belonging to the regions $`R_1`$ and $`R_2`$. Thus verification of (LABEL:8) gives information about the correlations between different phase-space regions.
Comment 5: It may be worth to point out that the scaling law (LABEL:7) and the additivity (LABEL:8) can be more precisely tested for Renyi entropies ($`H_k`$) than for the Shannon entropy ($`S`$) where the extrapolation procedure (described in Section 6) introduces always an additional uncertainty.
9. Conclusions
In conclusion, one sees that the measurement of entropy in limited regions of phase-space is feasible. Moreover, even the simplest tests of the general scaling and additivity rules can provide essential information on fluctuations and on correlations in the system. It should be emphasized that for these tests the Renyi entropies turn out to be more useful than the standard Shannon entropy.
Acknowledgements We thank Y.Foka for suggesting preparation of this note, to K.Fialkowski, A.Ostruszka and J.Wosiek for discussions and M.Gazdzicki for correspondence.
This investigation was supported in part by the KBN Grant No 2 P03B 086 14 and by Subsydium FNP 1/99. |
warning/0002/math0002132.html | ar5iv | text | # Difference Equations Compatible with Trigonometric KZ Differential Equations
## 1. Introduction
The trigonometric KZ equations associated with a Lie algebra $`𝔤`$ depend on a parameter $`\lambda 𝔥`$ where $`𝔥𝔤`$ is the Cartan subalgebra. We suggest a system of dynamical difference equations with respect to $`\lambda `$ compatible with the trigonometric KZ differential equations. The dynamical equations are constructed in terms of intertwining operators of $`𝔤`$-modules.
Our dynamical difference equations are a special example of the difference equations introduced by Cherednik. In \[Ch1, Ch2\] Cherednik introduces a notion of an affine R-matrix associated with the root system of a Lie algebra and taking values in an algebra $`F`$ with certain properties. Given an affine R-matrix, he defines a system of equations for an element of the algebra $`F`$.
In this paper we construct an example of an affine R-matrix and call the corresponding system of equations the dynamical equations. In our example, $`F`$ is the algebra of functions of complex variables $`z_1,\mathrm{},z_n`$ and $`\lambda 𝔥`$ taking values in the tensor product of $`n`$ copies of the universal enveloping algebra of $`𝔤`$. The fact that our dynamical difference equations are compatible with the trigonometric KZ differential equations is a remarkable property of our affine R-matrix.
There is a similar construction of dynamical difference equations compatible with the qKZ difference equations associated with a quantum group. The dynamical difference equations in that case are constructed in the same way in terms of interwining operators of modules over the quantum group. We will describe this construction in a forthcoming paper.
There is a degeneration of the trigonometric KZ differential equations to the standard (rational) KZ differential equations. Under this limiting procedure the dynamical difference equations constructed in this paper turn into the system of differential equations compatible with the standard KZ differential equations and described in \[FMTV\]. In \[FMTV\] we proved that the standard hypergeometric solutions of the standard KZ equations \[SV, V\] satisfy also the dynamic differential equations of \[FMTV\].
The trigonometric KZ differential equations also have hypergeometric solutions, see \[Ch3, EFK\]. We conjecture that the hypergeometric solutions of the trigonometric KZ differential equations also solve the dynamical difference equations of this paper.
In Section 2 we study relations between intertwining operators of $`𝔤`$-modules and the Weyl group $`𝕎`$ of $`𝔤`$. For any finite dimensional $`𝔤`$-module $`V`$ and $`w𝕎`$ we construct a rational function $`𝔹_{w,V}:\mathrm{End}(V)`$. The operators $`𝔹_{w,V}(\lambda )`$ are used later to construct an affine R-matrix and dynamical equations.
In Section 3 we define the dynamical difference equations for $`𝔤=sl_N`$ in terms of operators $`𝔹_{w,V}(\lambda )`$ directly (without introducing affine R-matrices). For $`𝔤=sl_N`$, we prove that the dynamical equations are compatible with the trigonometric KZ differential equations. We give a formula for the determinant of a square matrix solution of the combined system of KZ and dynamical equations.
In Section 4 we review \[Ch1, Ch2\] and construct the dynamical difference equations for any simple Lie algebra $`𝔤`$. We show that the dynamical equations are compatible with the trigonometric KZ equations if the Lie algebra $`𝔤`$ has minuscle weights, i.e. is not of type $`E_8,F_4,G_2`$. We conjecture that the dynamical difference equations and trigonometric KZ equations are compatible for any simple Lie algebra.
We thank I.Cherednik for valuable discussions and explanation of his articles \[Ch1, Ch2\] and P.Etingof who taught us all about the Weyl group and intertwining operators.
## 2. Intertwining Operators
### 2.1. Preliminaries
Let $`𝔤`$ be a complex simple Lie algebra with root space decomposition $`𝔤=𝔥(_{\alpha \mathrm{\Sigma }}𝔤_\alpha )`$ where $`\mathrm{\Sigma }𝔥^{}`$ is the set of roots.
Fix a system of simple roots $`\alpha _1,\mathrm{},\alpha _r`$. Let $`\mathrm{\Gamma }`$ be the corresponding Dynkin diagram, and $`\mathrm{\Sigma }_\pm `$ — the set of positive (negative) roots. Let $`𝔫_\pm =_{\alpha \mathrm{\Sigma }_\pm }𝔤_\alpha `$. Then $`𝔤=𝔫_+𝔥𝔫_{}`$.
Let $`(,)`$ be an invariant bilinear form on $`𝔤`$. The form gives rise to a natural identification $`𝔥𝔥^{}`$. We use this identification and make no distinction between $`𝔥`$ and $`𝔥^{}`$. This identification allows us to define a scalar product on $`𝔥^{}`$. We use the same notation $`(,)`$ for the pairing $`𝔥𝔥^{}`$.
We use the notation: $`Q=_{i=1}^r\alpha _i`$ \- root lattice; $`Q^+=_{i=1}^r_0\alpha _i`$; $`Q^{}=_{i=1}^r\alpha _i^{}`$ \- dual root lattice, where $`\alpha ^{}=2\alpha /(\alpha ,\alpha )`$; $`P=\{\lambda 𝔥|(\lambda ,\alpha _i^{})\}`$ \- weight lattice; $`P^+=\{\lambda 𝔥|(\lambda ,\alpha _i^{})_0\}`$ \- cone of dominant integral weights; $`\omega _iP^+`$ \- fundamental weights: $`(\omega _i,\alpha _j^{})=\delta _{ij}`$; $`\rho =\frac{1}{2}_{\alpha \mathrm{\Sigma }_+}\alpha =_{i=1}^r\omega _i`$; $`P^{}=_{i=1}^r\omega _i^{}`$ \- dual weight lattice, where $`\omega _i^{}`$ -dual fundamental weights: $`(\omega _i^{},\alpha _j)=\delta _{ij}`$.
Define a partial order on $`𝔥`$ putting $`\mu <\lambda `$ if $`\lambda \mu Q^+`$.
Let $`s_i:𝔥𝔥`$ denote a simple reflection, defined by $`s_i(\lambda )=\lambda (\alpha _i^{},\lambda )\alpha _i`$; $`𝕎`$ \- Weyl group, generated by $`s_1,\mathrm{},s_r`$. The following relations are defining:
$`s_i^2=1,(s_is_j)^m=1\text{for}m=2,3,4,6,`$
where $`m=2`$ if $`\alpha _i`$ and $`\alpha _j`$ are not neighboring in $`\mathrm{\Gamma }`$, otherwise, $`m=3,4,6`$ if 1,2,3 lines respectively connect $`\alpha _i`$ and $`\alpha _j`$ in $`\mathrm{\Gamma }`$. For an element $`w𝕎`$, denote $`l(w)`$ the length of the minimal (reduced) presentation of $`w`$ as a product of generators $`s_1,\mathrm{},s_r`$.
Let $`U𝔤`$ be the universal enveloping algebra of $`𝔤`$; $`U𝔤^n`$ \- tensor product of $`n`$ copies of $`U𝔤`$; $`\mathrm{\Delta }^{(n)}:U𝔤U𝔤^n`$ \- the iterated comultiplication (in particular, $`\mathrm{\Delta }^{(1)}`$ is the identity, $`\mathrm{\Delta }^{(2)}`$ is the comultiplication); $`U𝔤_0^n=\{xU𝔤^n|[\mathrm{\Delta }^{(n)}(h),x]=0\text{for any }h𝔥\}`$ \- subalgebra of weight zero elements.
For $`\alpha \mathrm{\Sigma }`$ choose generators $`e_\alpha 𝔤_\alpha `$ so that $`(e_\alpha ,e_\alpha )=1`$. For any $`\alpha `$, the triple
$`H_\alpha =\alpha ^{},E_\alpha ={\displaystyle \frac{2}{(\alpha ,\alpha )}}e_\alpha ,F_\alpha =e_\alpha `$
forms an $`sl_2`$-subalgebra in $`𝔤`$, $`[H_\alpha ,E_\alpha ]=2E_\alpha ,[H_\alpha ,F_\alpha ]=2F_\alpha ,[E_\alpha ,F_\alpha ]=H_\alpha `$.
A dual fundamental weight $`\omega _i^{}`$ is called minuscule if $`(\omega _i^{},\alpha )`$ is 0 or 1 for all $`\alpha \mathrm{\Sigma }_+`$, i.e. for any positive root $`\alpha =_{i=1}^rm_i\alpha _i`$, the coefficient $`m_i`$ is either 0 or 1. For a root system of type $`A_r`$ all dual fundamental weights are minuscule. There is no minuscule dual fundamental weight for $`E_8,F_4,G_2`$. For a minuscule dual fundamental weight $`\omega _i^{}`$, define an element $`w_{[i]}=w_0w_0^i𝕎`$ where $`w_0`$ (respectively, $`w_0^i`$) is the longest element in $`𝕎`$ (respectively, in $`𝕎^i`$ generated by all simple reflections $`s_j`$ preserving $`\omega _i^{}`$).
###### Lemma 1.
Let $`\alpha `$ be a positive root. Then $`w_{[i]}(\alpha )\mathrm{\Sigma }_+`$ if $`(\omega _i^{},\alpha )=0`$ and $`w_{[i]}(\alpha )\mathrm{\Sigma }_{}`$ if $`(\omega _i^{},\alpha )=1`$.
Let $`𝔾`$ be the simply connected complex Lie group with Lie algebra $`𝔤`$, $`𝔾`$ the Cartan subgroup corresponding to $`𝔥`$, $`N()=\{x𝔾|xx^1=\}`$ the normalizer of $``$. Then the Weyl group is canonically isomorphic to $`N()/`$. The isomorphism sends $`x`$ to Ad$`{}_{x}{}^{}|_{𝔥}^{}`$.
Let $`V`$ be a finite dimensional $`𝔤`$-module with weight decomposition $`V=_{\mu 𝔥}V[\mu ]`$. $`𝔾`$ acts on $`V`$ so that $``$ acts trivially on $`V[0]`$. Thus the action of $`𝕎`$ on $`V[0]`$ is well defined. For any $`n`$, the Weyl group in the same way acts also on $`U𝔤_0^n`$.
###### Lemma 2.
For $`\alpha \mathrm{\Sigma }`$ and $`k_0`$, consider $`e_\alpha ^ke_\alpha ^kU𝔤_0`$ and $`e_\alpha e_\alpha U𝔤_0^{\mathrm{\hspace{0.17em}2}}`$. Then for any $`w𝕎`$,
$`w(e_\alpha ^ke_\alpha ^k)=e_{w(\alpha )}^ke_{w(\alpha )}^k,w(e_\alpha e_\alpha )=e_{w(\alpha )}e_{w(\alpha )}.`$
###### Proof.
Let $`xN()`$ be a lifting of $`w`$. Ad$`{}_{x}{}^{}:𝔤𝔤`$ is an automorphism of $`𝔤`$ preserving the invariant scalar product and sending $`𝔤_\beta `$ to $`𝔤_{w(\beta )}`$ for all $`\beta `$. Thus, Ad$`{}_{x}{}^{}e_{\beta }^{}=c_{x,\beta }e_{w(\beta )}`$ for suitable numbers $`c_{x,\beta }`$ and $`c_{x,\alpha }c_{x,\alpha }=1`$. ∎
Let $`x_1,\mathrm{},x_r`$ be an orthonormal basis in $`𝔥`$, set
$`\mathrm{\Omega }^0={\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_ix_i,\mathrm{\Omega }^+=\mathrm{\Omega }^0+{\displaystyle \underset{\alpha \mathrm{\Sigma }_+}{}}e_\alpha e_\alpha ,\mathrm{\Omega }^{}=\mathrm{\Omega }^0+{\displaystyle \underset{\alpha \mathrm{\Sigma }_+}{}}e_\alpha e_\alpha .`$
Define the Casimir operator $`\mathrm{\Omega }`$ and the trigonometric R-matrix $`r(z)`$ by
$`\mathrm{\Omega }=\mathrm{\Omega }^++\mathrm{\Omega }^{},r(z)={\displaystyle \frac{\mathrm{\Omega }^+z+\mathrm{\Omega }^{}}{z1}}.`$
For any $`xU𝔤`$, we have $`\mathrm{\Delta }(x)\mathrm{\Omega }=\mathrm{\Omega }\mathrm{\Delta }(x)`$. We will use a more symmetric form of the trigonometric R-matrix: $`r(z_1/z_2)`$.
The Weyl group acts on $`r(z_1/z_2),\mathrm{\Omega }U𝔤_0^{\mathrm{\hspace{0.17em}2}}`$. $`\mathrm{\Omega }`$ is Weyl invariant. For any $`w𝕎`$,
$`w(r(z_1/z_2))={\displaystyle \frac{1}{z_1z_2}}({\displaystyle \frac{z_1+z_2}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_ix_i+{\displaystyle \underset{\alpha \mathrm{\Sigma }_+}{}}(z_1e_{w(\alpha )}e_{w(\alpha )}+z_2e_{w(\alpha )}e_{w(\alpha )})).`$
###### Lemma 3.
For a minuscule dual fundamental weight $`\omega _i^{}`$,
$`z_1^{(\omega _i^{})^{(1)}}z_2^{(\omega _i^{})^{(2)}}r(z_1/z_2)z_1^{(\omega _i^{})^{(1)}}z_2^{(\omega _i^{})^{(2)}}=w_{[i]}^1(r(z_1/z_2)).`$
Proof. Using Lemma 1 it is easy to see that both sides of the equation are equal to
$`{\displaystyle \frac{1}{z_1z_2}}({\displaystyle \frac{z_1+z_2}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_ix_i+{\displaystyle \underset{\alpha \mathrm{\Sigma }_+,(\alpha ,\omega _i^{})=0}{}}(z_1e_\alpha e_\alpha +z_2e_\alpha e_\alpha )+`$
$`{\displaystyle \underset{\alpha \mathrm{\Sigma }_+,(\alpha ,\omega _i^{})=1}{}}(z_1e_\alpha e_\alpha +z_2e_\alpha e_\alpha )).\mathrm{}`$
### 2.2. The Trigonometric KZ Equations
Let $`V=V_1\mathrm{}V_n`$ be a tensor product of $`𝔤`$-modules. For $`\kappa `$ and $`\lambda 𝔥`$, introduce the KZ operators $`_i(\lambda ,\kappa ),i=1,\mathrm{},n,`$ acting on functions $`u(z_1,\mathrm{},z_n)`$ of $`n`$ complex variables with values in $`V`$ and defined by
$`_i(\lambda ,\kappa )=\kappa z_i{\displaystyle \frac{}{z_i}}{\displaystyle \underset{j,ji}{}}r(z_i/z_j)^{(i,j)}\lambda ^{(i)}.`$
Here $`r^{(i,j)}`$, $`\lambda ^{(i)}`$ denote $`r`$ acting in the $`i`$-th and $`j`$-th factors of the tensor product and $`\lambda `$ acting in the $`i`$-th factor.
The trigonometric KZ equations are the equations
(1) $`_i(\lambda ,\kappa )u(z_1,\mathrm{},z_n,\lambda )=\mathrm{\hspace{0.17em}0},i=1,\mathrm{},n,`$
see \[EFK\]. The KZ equations are compatible, $`[_i,_j]=0`$.
### 2.3. Intertwining Operators, Fusion Matrices, \[ES, EV1\]
For $`\lambda 𝔥`$, let $`M_\lambda `$ be the Verma module over $`𝔤`$ with highest weight $`\lambda `$ and highest weight vector $`v_\lambda `$. We have $`𝔫_+v_\lambda =0`$, and $`hv_\lambda =(h,\lambda )v_\lambda `$ for all $`h𝔥`$. Let $`M_\lambda =_{\mu \lambda }M_\lambda [\mu ]`$ be the weight decomposition. The Verma module is irreducible for a generic $`\lambda `$. Define the dual Verma module $`M_\lambda ^{}`$ to be the graded dual space $`_{\mu \lambda }M_\lambda ^{}[\mu ]`$ equipped with the $`𝔤`$-action: $`u,av=au,v`$ for all $`a𝔤,uM_\lambda ,vM_\lambda ^{}`$. Let $`v_\lambda ^{}`$ be the lowest weight vector of $`M_\lambda ^{}`$ satisfying $`v_\lambda ,v_\lambda ^{}=1`$.
Let $`V`$ be a finite dimensional $`𝔤`$-module with weight decompostion $`V=_{\mu 𝔥}V[\mu ]`$. For $`\lambda ,\mu 𝔥`$ consider an intertwining operator $`\mathrm{\Phi }:M_\lambda M_\mu V`$. Define its expectation value by $`\mathrm{\Phi }=\mathrm{\Phi }(v_\lambda ),v_\mu ^{}V[\lambda \mu ]`$. If $`M_\mu `$ is irreducible, then the map Hom$`{}_{𝔤}{}^{}(M_\lambda ,M_\mu V)V[\lambda \mu ],\mathrm{\Phi }\mathrm{\Phi }`$, is an isomorphism. Thus for any $`vV[\lambda \mu ]`$ there exists a unique intertwining operator $`\mathrm{\Phi }_\lambda ^v:M_\lambda M_\mu V`$ such that $`\mathrm{\Phi }_\lambda ^v(v_\lambda )v_\lambda v+_{\nu <\mu }M_\mu [\nu ]V`$.
Let $`V,W`$ be finite-dimensional $`𝔤`$-modules and $`vV[\mu ],wW[\nu ]`$. Consider the composition
$`\mathrm{\Phi }_\lambda ^{w,v}:M_\lambda \stackrel{\mathrm{\Phi }_\lambda ^v}{}M_{\lambda \mu }V\stackrel{\mathrm{\Phi }_{\lambda \mu }^w}{}M_{\lambda \mu \nu }WV.`$
Then $`\mathrm{\Phi }_\lambda ^{w,v}\mathrm{Hom}_𝔤(M_\lambda ,M_{\lambda \mu \nu }WV)`$. Hence, for a generic $`\lambda `$ there exists a unique element $`u(VW)[\mu +\nu ]`$ such that $`\mathrm{\Phi }_\lambda ^u=\mathrm{\Phi }_\lambda ^{w,v}`$. The assignment $`(w,v)u`$ is bilinear, and defines an $`𝔥`$-linear map
$$J_{WV}(\lambda ):WVWV.$$
The operator $`J_{WV}(\lambda )`$ is called the fusion matrix of $`W`$ and $`V`$. The fusion matrix $`J_{WV}(\lambda )`$ is a rational function of $`\lambda `$. $`J_{WV}(\lambda )`$ is strictly lower triangular, i.e. $`J=1+L`$ where $`L(W[\nu ]V[\mu ])_{\tau <\nu ,\mu <\sigma }W[\tau ]V[\sigma ]`$. In particular, $`J_{WV}(\lambda )`$ is invertible.
If $`V_1,\mathrm{}V_n`$ are $`𝔥`$-modules and $`F(\lambda ):V_1\mathrm{}V_nV_1\mathrm{}V_n`$ is a linear operator depending on $`\lambda 𝔥`$, then for any homogeneous $`u_1,\mathrm{},u_n`$, $`u_iV_i[\nu _i]`$, we define $`F(\lambda h^{(i)})(u_1\mathrm{}u_n)`$ to be $`F(\lambda \nu _i)(u_1\mathrm{}u_n)`$.
There is a universal fusion matrix $`J(\lambda )U𝔤_0^{\mathrm{\hspace{0.17em}2}}`$ such that $`J_{WV}(\lambda )=J(\lambda )|_{WV}`$ for all $`W,V`$. The universal fusion matrix $`J(\lambda )`$ is the unique solution of the \[ABRR\] equation
$`J(\lambda )(1(\lambda +\rho {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_i^2))=(1(\lambda +\rho {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_i^2)+{\displaystyle \underset{\alpha \mathrm{\Sigma }_+}{}}e_\alpha e_\alpha )J(\lambda ).`$
such that $`\left(J(\lambda )1\right)𝔟_{}(U𝔟_{})(U𝔟_+)𝔟_+`$ where $`𝔟_\pm =𝔥𝔫_\pm `$.
We transform this equation to a more convenient form. The equation can be written as
$`J(\lambda )(\lambda +\rho {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_i^2)^{(2)}=((\lambda +\rho {\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_i^2)^{(2)}{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_ix_i+\mathrm{\Omega }^{})J(\lambda ).`$
We make a change of variables: $`\lambda \lambda \rho +\frac{1}{2}(h^{(1)}+h^{(2)})`$. Then the equation takes the form
$`J(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))(\lambda +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}){\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_i^2)^{(2)}=`$
$`((\lambda +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}){\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_i^2)^{(2)}{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{r}{}}}x_ix_i+\mathrm{\Omega }^{})J(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})).`$
Notice that $`(h^{(1)}+h^{(2)})^{(2)}=_{i=1}^rx_i^{(2)}(x_i^{(1)}+x_i^{(2)})`$. Now the equation takes the form
(2) $`J(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))(\lambda ^{(2)}+\mathrm{\Omega }^0)=(\lambda ^{(2)}+\mathrm{\Omega }^{})J(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})).`$
For $`w𝕎`$, let $`w(J(\lambda ))`$ be the image of $`J(\lambda )`$ under the action of $`w`$. Let $`xN()`$ be a lifting of $`w`$. Let $`W,V`$ be finite dimensional $`𝔤`$-modules. Then
(3) $`w(J(\lambda ))|_{WV}=xJ_{WV}(\lambda )x^1,`$
and RHS does not depend on the choice of $`x`$.
### 2.4. Main Construction, I
Introduce a new action of the Weyl group $`𝕎`$ on $`𝔥`$ by
$$w\lambda =w(\lambda +\rho )\rho .$$
Remind facts from \[BGG\].
Let $`M_\mu ,M_\lambda `$ be Verma modules. Two cases are possible: a) Hom$`{}_{𝔤}{}^{}(M_\mu ,M_\lambda )=0`$,
b) Hom$`{}_{𝔤}{}^{}(M_\mu ,M_\lambda )=`$ and every nontrivial homomorphism $`M_\mu M_\lambda `$ is an embedding.
Let $`M_\lambda `$ be a Verma module with dominant weight $`\lambda P^+`$. Then Hom$`{}_{𝔤}{}^{}(M_\mu ,M_\lambda )=`$ if and only if there is $`w𝕎`$ such that $`\mu =w\lambda `$.
Let $`w=s_{i_k}\mathrm{}s_{i_1}`$ be a reduced presentation. Set $`\alpha ^1=\alpha _{i_1}`$ and $`\alpha ^j=(s_{i_1}\mathrm{}s_{i_{j1}})(\alpha _{i_j})`$ for $`j=2,\mathrm{},k`$. Let $`n_j=(\lambda +\rho ,(\alpha ^j)^{})`$. For a dominant $`\lambda P^+`$, $`n_j`$ are positive integers.
###### Lemma 4.
The collection of integers $`n_1,\mathrm{}n_k`$ and the product $`(e_{\alpha _{i_k}})^{n_k}\mathrm{}(e_{\alpha _{i_1}})^{n_1}`$ do not depend on the reduced presentation.
###### Proof.
It is known that $`\alpha ^1,\mathrm{},\alpha ^k`$ are distinct positive roots and $`\{\alpha ^1,\mathrm{},\alpha ^k\}=\{\alpha \mathrm{\Sigma }_+|w(\alpha )\mathrm{\Sigma }_{}\}`$. Hence, the collection $`n_1,\mathrm{}n_k`$ does not depend on the reduced presentation.
The vector $`(e_{\alpha _{i_k}})^{n_k}\mathrm{}(e_{\alpha _{i_1}})^{n_1}v_\lambda `$ is a singular vector in $`M_\lambda `$. If $`w=s_{i_k^{}}\mathrm{}s_{i_1^{}}`$ is another reduced presentation, then the vectors $`(e_{\alpha _{i_k}})^{n_k}\mathrm{}(e_{\alpha _{i_1}})^{n_1}v_\lambda `$ and
$`(e_{\alpha _{i_k^{}}})^{n_k^{}}\mathrm{}(e_{\alpha _{i_1^{}}})^{n_1^{}}v_\lambda `$ are proportional. Since $`M_\lambda `$ is a free $`𝔫_{}`$-module, we have
$`(e_{\alpha _{i_k^{}}})^{n_k^{}}\mathrm{}(e_{\alpha _{i_1^{}}})^{n_1^{}}=c(e_{\alpha _{i_k}})^{n_k}\mathrm{}(e_{\alpha _{i_1}})^{n_1}`$ in $`𝔫_{}`$ for a suitable $`c`$. $`c=1`$ since the monomials are equal when projected to the commutative polynomial algebra generated by $`e_{\alpha _1},\mathrm{},e_{\alpha _r}`$. ∎
Define a singular vector $`v_{w\lambda }^\lambda M_\lambda `$ by
(4) $`v_{w\lambda }^\lambda ={\displaystyle \frac{(e_{\alpha _{i_k}})^{n_k}}{n_1!}}\mathrm{}{\displaystyle \frac{(e_{\alpha _{i_1}})^{n_1}}{n_k!}}v_\lambda .`$
This vector does not depend on the reduced presentation by Lemma 4.
For all $`\lambda P^+`$, $`w𝕎`$, fix an embedding $`M_{w\lambda }M_\lambda `$ sending $`v_{w\lambda }`$ to $`v_{w\lambda }^\lambda `$.
Let $`V`$ be a finite dimensional $`𝔤`$-module, $`V=_{\nu 𝔥}V[\nu ]`$ the weight decomposition, $`P(V)=\{\nu 𝔥|V[\nu ]0\}`$ the set of weights of $`V`$. We say that $`\lambda P^+`$ is generic with respect to $`V`$ if
1. For any $`\nu P(V)`$ there exist a unique intertwining operator $`\mathrm{\Phi }_\lambda ^v:M_\lambda M_{\lambda \nu }V`$ such that $`\mathrm{\Phi }_\lambda ^v(v_\lambda )=v_{\lambda \nu }v+`$ lower order terms.
2. For any $`w,w^{}𝕎,ww^{}`$, and any $`\nu P(V)`$, the vector $`w\lambda w^{}(\lambda \nu )`$ does not belong to $`P(V)`$.
It is clear that all dominant weights lying far inside the cone of dominant weights are generic with respect to $`V`$.
###### Lemma 5.
Let $`\lambda P^+`$ be generic with respect to $`V`$. Let $`vV[\nu ]`$. Consider the intertwining operator $`\mathrm{\Phi }_\lambda ^v:M_\lambda M_{\lambda \nu }V`$. For $`w𝕎`$, consider the singular vector $`v_{w\lambda }^\lambda M_\lambda `$. Then there exists a unique vector $`A_{w,V}(\lambda )(v)V[w(\nu )]`$ such that
$`\mathrm{\Phi }_\lambda ^v(v_{w\lambda }^\lambda )=v_{w(\lambda \nu )}^{\lambda \nu }A_{w,V}(\lambda )(v)+\text{lower order terms}.`$
###### Proof.
$`\mathrm{\Phi }_\lambda ^v(v_{w\lambda }^\lambda )`$ is a singular vector in $`M_{\lambda \nu }V`$. It has to have weight components of the form $`v_{w^{}(\lambda \nu )}^{\lambda \nu }u`$ for suitable $`w^{}𝕎`$ and $`uV`$. Since $`\lambda `$ is generic, we have $`w=w^{}`$ and $`\mathrm{\Phi }_\lambda ^v(v_{w\lambda }^\lambda )`$ is of the required form for a suitable $`A_{w,V}(\lambda )(v)V[w(\nu )]`$. ∎
For generic $`\lambda P^+`$, Lemma 5 defines a linear operator $`A_{w,V}(\lambda ):VV`$ such that $`A_{w,V}(\lambda )(V[\nu ]))V[w(\nu )]`$ for all $`\nu P(V)`$. It follows from calculations in Section 2.5 that $`A_{w,V}(\lambda )`$ is a rational function of $`\lambda 𝔥`$.
The following Lemmas are easy consequences of definitions.
###### Lemma 6.
If $`w_1,w_2𝕎`$ and $`l(w_1w_2)=l(w_1)+l(w_2)`$, then
$`A_{w_1w_2,V}(\lambda )=A_{w_1,V}(w_2\lambda )A_{w_2,V}(\lambda ).`$
###### Lemma 7.
Let $`W,V`$ be finite dimensional $`𝔤`$-modules. Let $`w𝕎`$. Then
$`A_{w,WV}(\lambda )J_{WV}(\lambda )=J_{WV}(w\lambda )(A_{w,W}(\lambda h^{(2)})A_{w,V}(\lambda )).`$
Let $`x_wN()𝔾`$ be a lifting of $`w𝕎`$. For a finite dimensional $`𝔤`$-module $`V`$, define an operator
$`B_{x_w,V}(\lambda ):VV,vx_w^1A_{w,V}(\lambda )v.`$
$`B_{x_w,V}`$ preserves the weight of elements of $`V`$.
Lemma 7 implies
$`B_{x_w,WV}(\lambda )J_{WV}(\lambda )=(x_w^1J_{WV}(w\lambda )x_w)(B_{x_w,W}(\lambda h^{(2)})B_{x_w,V}(\lambda )),`$
cf. (3).
The operator $`B_{x_w,V}`$ depends on the choice of $`x_w`$. If $`x_wg,g`$, is another lifting of $`w`$, then $`B_{x_wg,V}=g^1B_{x_w,V}`$.
The operators $`B_{x_w,V}(\lambda )`$, $`w𝕎`$, are defined now for generic dominant $`\lambda `$ and depend on the choice of liftings $`x_w`$. In the next two Sections we fix a normalization $`B_{w,V}(\lambda )`$ of $`B_{x_w,V}(\lambda )`$ so that $`B_{w,V}(\lambda )\mathrm{\hspace{0.17em}1}`$ as $`\lambda \mathrm{}`$. We show that for any $`w𝕎`$, there is a universal $`B_w(\lambda )U𝔤_0`$ such that $`B_w(\lambda )|_V=B_{w,V}(\lambda )`$ for every finite dimensional $`𝔤`$-module $`V`$. For any $`w𝕎`$, we present $`B_w(\lambda )`$ as a suitable product of operators $`B_{s_i}(\lambda )`$ corresponding to simple reflections.
### 2.5. Operators $`B_{x_w,V}(\lambda )`$ for $`𝔤=sl_2`$
Consider $`sl_2`$ with generators $`H,E,F`$ and relations $`[H,E]=2E,[H,F]=2F,[E,F]=H`$. Let $`\alpha _1`$ be the positive root. Identifying $`𝔥`$ and $`𝔥^{}`$, we have $`\alpha _1=\alpha _1^{}=H`$, $`\omega _1=\omega _1^{}=H/2`$, $`𝕎=\{1,s_1\}`$.
Let $`\lambda =l\omega _1`$, $`l_0`$, be a dominant weight. Then $`s_1\lambda =(l+2)\omega _1`$. For any dominant weight $`\lambda `$, fix an embedding
$`M_{s_1\lambda }M_\lambda ,v_{s_1\lambda }v_{s_1\lambda }^\lambda ={\displaystyle \frac{F^{(\lambda ,\alpha _1)+1}v_\lambda }{((\lambda ,\alpha _1)+1)!}}`$
as in Section 2.4.
For $`m_0`$, let $`L_m`$ be the irreducible $`sl_2`$ module with highest weight $`m\omega _1`$. $`L_m`$ has a basis $`v_0^m,\mathrm{},v_m^m`$ such that
$$Hv_k^m=(m2k)v_k^m,Fv_k^m=(k+1)v_{k+1}^m,Ev_k^m=(mk+1)v_{k1}^m.$$
For $`𝔤=sl_2`$, we have $`𝔾=SL(2,)`$. Then $`𝔾`$ is the subgroup of diagonal matrices. Fix a lifting $`xN()`$ of $`s_1`$, set $`x=(x_{ij})`$ where $`x_{11}=x_{22}=0`$, $`x_{12}=1`$, $`x_{21}=1`$. Then the action of $`x`$ in $`L_m`$ is given by $`v_k^m(1)^kv_{mk}^m`$ for any $`k`$. We have $`x=\text{exp}(E)\text{exp}(F)\text{exp}(E)`$.
For $`t`$, introduce
(5) $`p(t;H,E,F)={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}F^kE^k{\displaystyle \frac{1}{k!}}{\displaystyle \underset{j=0}{\overset{k1}{}}}{\displaystyle \frac{1}{(tHj)}}.`$
$`p(t;H,E,F)`$ is an element of $`U(sl_2)_0`$.
###### Theorem 8.
Let $`\lambda `$ be a dominant weight for $`sl_2`$. Let $`L_m,x`$ be as above. Let $`B_{x,L_m}(\lambda ):L_mL_m`$ be the operator defined in Section 2.4. Then for $`k=0,\mathrm{},m`$,
(6)
$`B_{x,L_m}(\lambda )v_k^m={\displaystyle \frac{((\lambda ,\alpha _1^{})+2)((\lambda ,\alpha _1^{})+3)\mathrm{}((\lambda ,\alpha _1^{})+k+1)}{((\lambda ,\alpha _1^{})m+k+1)((\lambda ,\alpha _1^{})m+k+2)\mathrm{}((\lambda ,\alpha _1^{})m+2k)}}v_k^m`$
and
(7) $`p((\lambda ,\alpha _1^{});H,E,F)|_{L_m}=B_{x,L_m}(\lambda ).`$
###### Corollary 9.
$`B_{x,L_m}(\lambda )`$ is a rational function of $`(\lambda ,\alpha _1^{})`$. $`B_{x,L_m}(\lambda )`$ tends to $`1`$ as $`(\lambda ,\alpha _1^{})`$ tends to infinity.
The Theorem is proved by direct verification. First we calculate explicitly $`\mathrm{\Phi }_\lambda ^{v_k^m}(v_\lambda )`$, $`\mathrm{\Phi }_\lambda ^{v_k^m}(\frac{F^{(\lambda ,\alpha _1^{})+1}}{((\lambda ,\alpha _1^{})+1)!}v_\lambda )`$, and then get an expression for $`B_{x,L_m}(\lambda )v_k^m`$ as a sum of a hypergeometric type. Using standard formulas from \[GR\] we see that $`B_{x,L_m}(\lambda )v_k^m`$ is given by (6). Similarly we check that $`p((\lambda ,\alpha _1^{});H,E,F)v_k^m`$ gives the same result. Thus we get (7). $`\mathrm{}`$
Formula (6) becomes more symmetric if $`\lambda `$ is replaced by $`\lambda \rho +\frac{1}{2}\nu `$ where $`\nu =m\omega _1k\alpha _1`$ is the weight of $`v_k^m`$, then
(8) $`p((\lambda +{\displaystyle \frac{1}{2}}\nu ,\alpha _1^{})1;H,E,F)v_k^m={\displaystyle \underset{j=0}{\overset{k1}{}}}{\displaystyle \frac{(\lambda ,\alpha _1^{})+\frac{m}{2}j}{(\lambda ,\alpha _1^{})\frac{m}{2}+j}}v_k^m.`$
###### Theorem 10.
$`p(t2;H,F,E)p(t;H,E,F))={\displaystyle \frac{tH+1}{t+1}}.`$
To prove this formula it is enough to check that RHS and LHS give the same result when applied to $`v_k^mL_m`$, which is done using (8). $`\mathrm{}`$
Notice that $`p(t;H,F,E)=s_1(p(t;H,E,F))`$.
Remark. Let $`J(\lambda )=_ia_ib_i`$ be the universal fusion matrix of $`sl_2`$. Following \[EV2\] introduce $`S(Q)(\lambda )U(sl_2)_0`$ as $`S(Q)(\lambda )=_iS(a_i)b_i`$ where $`S(a_i)`$ is the antipode of $`a_i`$. The action of $`S(Q)(\lambda )`$ in $`L_m`$ was computed in \[EV2\]. Comparing the result with Theorem 8, one sees that $`p((\lambda ,\alpha _1^{});H,E,F)`$ is equal to $`(S(Q)(\lambda ))^1`$ up to a simple change of argument $`\lambda `$.
###### Corollary 11.
Let $`A_{s_1,L_m}(\lambda ):L_mL_m`$ be the operator defined in Section 2.4. Then $`A_{s_1,L_m}(\lambda )=xp((\lambda ,\alpha _1^{});H,E,F)|_{L_m}`$. $`A_{s_1,L_m}(\lambda )`$ is a rational function of $`(\lambda ,\alpha _1^{})`$. $`A_{s_1,L_m}(\lambda )`$ tends to $`x`$ as $`(\lambda ,\alpha _1^{})`$ tends to infinity.
### 2.6. Main Construction, II
Return to the situation considered in Section 2.4.
For any simple root $`\alpha _i`$, the triple $`H_{\alpha _i},E_{\alpha _i},F_{\alpha _i}`$ defines an embedding $`sl_2𝔤`$ and induces an embedding $`SL(2,)𝔾`$. Denote $`x_i𝔾`$ the image under this embedding of the element $`xSL(2,)`$ defined in Section 2.5.
###### Lemma 12.
For $`i=1,\mathrm{},r`$, we have $`x_iN()`$ and Ad$`{}_{x_i}{}^{}:𝔤𝔤`$ restricted to $`𝔥`$ is the simple reflection $`s_i:𝔥𝔥`$.
###### Proof.
Since $`x_i=\mathrm{exp}(E_{\alpha _i})`$ $`\mathrm{exp}(F_{\alpha _i})`$ $`\mathrm{exp}(E_{\alpha _i})`$, we have that Ad$`{}_{x_i}{}^{}(H_{\alpha _i})=H_{\alpha _i}`$ and Ad$`{}_{x_i}{}^{}(h)=h`$ for any $`h𝔥`$ orthogonal to $`\alpha _i`$. Hence $`x_iN()`$ and Ad$`{}_{x_i}{}^{}|_{𝔥}^{}=s_i`$. ∎
For $`i=1,\mathrm{},r`$ and $`\lambda 𝔥`$, set
$`B_{s_i}(\lambda )=p((\lambda ,\alpha _i^{});H_{\alpha _i},E_{\alpha _i},F_{\alpha _i})`$
where $`p(t;H,E,F)`$ is defined in (5). Set
$`A_{s_i}(\lambda )=x_iB_{s_i}(\lambda ).`$
For any $`\nu P(V)`$, we have $`A_{s_i}(\lambda )(V[\nu ])V[s_i(\nu )]`$.
Let $`V`$ be a finite dimensional $`𝔤`$-module. For $`w𝕎`$, let $`w=s_{i_k}\mathrm{}s_{i_1}`$ be a reduced presentation. For a generic dominant $`\lambda P^+`$, consider the operator $`A_{w,V}(\lambda ):VV`$ defined in Section 2.4.
###### Lemma 13.
$`A_{w,V}(\lambda )=A_{s_{i_k}}((s_{i_{k1}}\mathrm{}s_{i_1})\lambda )|_VA_{s_{i_{k1}}}((s_{i_{k2}}\mathrm{}s_{i_1})\lambda )|_V\mathrm{}A_{s_{i_1}}(\lambda )|_V.`$
###### Proof.
See Corollary 11 and Lemma 6. ∎
###### Corollary 14.
The operator $`A_{w,V}(\lambda )`$ is a rational function of $`\lambda `$. $`A_{w,V}(\lambda )`$ tends to $`x_{i_k}\mathrm{}x_{i_1}`$ as $`\lambda `$ tends to infinity in a generic direction. In particular, the product $`x_{i_k}\mathrm{}x_{i_1}`$ does not depend on the choice of the reduced presentation.
Set $`x_w=x_{i_k}\mathrm{}x_{i_1}`$. $`x_wN()`$ is a lifting of $`w`$. Consider the operator $`B_{x_w,V}(\lambda ):VV`$ defined in Section 2.4 for this lifting $`x_w`$. Denote this operator $`B_{w,V}(\lambda )`$.
###### Corollary 15.
$`B_{w,V}(\lambda )=`$
$`(s_{i_{k1}}\mathrm{}s_{i_1})^1(B_{s_{i_k}}((s_{i_{k1}}\mathrm{}s_{i_1})\lambda ))|_V(s_{i_{k2}}\mathrm{}s_{i_1})^1(B_{s_{i_{k1}}}((s_{i_{k2}}\mathrm{}s_{i_1})\lambda ))|_V\mathrm{}B_{s_{i_1}}(\lambda )|_V.`$
$`B_{w,V}(\lambda )`$ is a rational function of $`\lambda `$. $`B_{w,V}(\lambda )`$ tends to $`1`$ as $`\lambda `$ tends to infinity in a generic direction.
For any notrivial element $`w𝕎`$ and $`\lambda 𝔥`$, define an element $`B_w(\lambda )U𝔤_0`$ by
$`B_w(\lambda )=`$
$`(s_{i_{k1}}\mathrm{}s_{i_1})^1(B_{s_{i_k}}((s_{i_{k1}}\mathrm{}s_{i_1})\lambda ))(s_{i_{k2}}\mathrm{}s_{i_1})^1(B_{s_{i_{k1}}}((s_{i_{k2}}\mathrm{}s_{i_1})\lambda ))\mathrm{}B_{s_{i_1}}(\lambda ).`$
Set $`B_w(\lambda )=1`$ if $`w`$ is the identity in $`𝕎`$. We have $`B_w(\lambda )|_V=B_{w,V}(\lambda )`$, and $`B_w(\lambda )`$ does not depend on the choice of the reduced presentation of $`w`$.
Properties of $`B_w(\lambda )`$.
1. If $`w_1,w_2𝕎`$ and $`l(w_1w_2)=l(w_1)+l(w_2)`$, then
$`B_{w_1w_2}(\lambda )=(w_2)^1(B_{w_1}(w_2\lambda ))B_{w_2}(\lambda ).`$
2. Let $`i=1,\mathrm{},r`$, $`\omega 𝔥`$, and $`(\alpha _i,\omega )=0`$, then
$`B_{s_i}(\lambda +\omega )=B_{s_i}(\lambda ).`$
3. For $`i=1,\mathrm{},r`$,
$`s_i(B_{s_i}(s_i\lambda ))B_{s_i}(\lambda )={\displaystyle \frac{(\lambda ,\alpha _i^{})H_{\alpha _i}+1}{(\lambda ,\alpha _i^{})+1}}.`$
4. Every relation $`(s_is_j)^m=1`$ for $`m=2,3,4,6`$ in $`𝕎`$ is equivalent to a homogeneous relation $`s_is_j\mathrm{}=s_js_i\mathrm{}`$. Every such a homogeneous relation generates a relation for $`B_{s_i}(\lambda ),B_{s_j}(\lambda )`$. Namely, for $`m=2`$, the relation is
$`(s_j)^1(B_{s_i}(s_j\lambda ))B_{s_j}(\lambda )=(s_i)^1(B_{s_j}(s_i\lambda ))B_{s_i}(\lambda ),`$
for $`m=3`$, the relation is
$`(s_js_i)^1(B_{s_i}((s_js_i)\lambda ))(s_i)^1(B_{s_j}(s_i\lambda ))B_{s_i}(\lambda )=`$
$`(s_is_j)^1(B_{s_j}((s_is_j)\lambda ))(s_j)^1(B_{s_i}(s_j\lambda ))B_{s_j}(\lambda ),`$
and so on.
5. $`\mathrm{\Delta }(B_w(\lambda ))J(\lambda )=w^1(J(w\lambda ))(B_w(\lambda h^{(2)})B_w(\lambda )).`$
The operators $`B_w(\lambda )`$ are closely connected with extremal projectors of Zhelobenko, see \[Zh1, Zh2\].
### 2.7. Operators $`𝔹_{w,V}`$
In order to study interrelations of operators $`B_{w,V}(\lambda )`$ with KZ operators it is convenient to change the argument $`\lambda `$.
Let $`V`$ be a finite dimensional $`𝔤`$-module. For $`w_1,w_2𝕎`$ and $`\lambda 𝔥`$, define $`w_1(𝔹_{w_2,V}(\lambda )):VV`$ as follows. For any $`\nu P(V)`$ and $`vV[\nu ]`$, set
$`w_1(𝔹_{w_2,V}(\lambda ))v=w_1(B_{w_2}(\lambda \rho +{\displaystyle \frac{1}{2}}\nu ))|_Vv.`$
In particular,
$`𝔹_{w,V}(\lambda )v=B_{w,V}(\lambda \rho +{\displaystyle \frac{1}{2}}\nu )v.`$
$`w_1(𝔹_{w_2,V}(\lambda ))`$ is a meromorphic function of $`\lambda `$, $`w_1(𝔹_{w_2,V}(\lambda ))`$ tends to 1 as $`\lambda `$ tends to infinity in a generic direction.
Properties of $`𝔹_{w,V}(\lambda )`$.
1. If $`w_1,w_2𝕎`$ and $`l(w_1w_2)=l(w_1)+l(w_2)`$, then
$`𝔹_{w_1w_2,V}(\lambda ))=w_2^1(𝔹_{w_1,V}(w_2(\lambda )))𝔹_{w_2,V}(\lambda ).`$
2. If $`i=1,\mathrm{},r`$, $`w𝕎`$, $`vV[\nu ]`$, then
$`𝔹_{s_i,V}(\lambda )v=p((\lambda +{\displaystyle \frac{1}{2}}\nu ,\alpha _i^{})1;H_{\alpha _i},E_{\alpha _i},F_{\alpha _i})v`$
and
$`w(𝔹_{s_i,V}(w^1(\lambda )))v=p((\lambda +{\displaystyle \frac{1}{2}}\nu ,w(\alpha _i^{}))1;H_{w(\alpha _i)},E_{w(\alpha _i)},F_{w(\alpha _i)})v`$
where $`p(t;H,E,F)`$ is defined in (5).
For $`\alpha \mathrm{\Sigma },\lambda 𝔥`$, define a linear operator $`𝔹_V^\alpha (\lambda ):VV`$ by
$$𝔹_V^\alpha (\lambda )v=p((\lambda +\frac{1}{2}\nu ,\alpha ^{})1;H_\alpha ,E_\alpha ,F_\alpha )v$$
for any $`vV[\nu ]`$.
1. $$𝔹_V^\alpha (\lambda )𝔹_V^\alpha (\lambda )v=\frac{(\lambda \frac{1}{2}\nu ,\alpha ^{})}{(\lambda +\frac{1}{2}\nu ,\alpha ^{})}v$$
for any $`vV[\nu ]`$.
2. Let $`\alpha \mathrm{\Sigma }`$, $`\omega 𝔥`$, and $`(\alpha ,\omega )=0`$, then
$`𝔹_V^\alpha (\lambda +\omega )=𝔹_V^\alpha (\lambda ).`$
3. Every relation $`(s_is_j)^m=1`$ for $`m=2,3,4,6`$ in $`𝕎`$ is equivalent to a homogeneous relation $`s_is_j\mathrm{}=s_js_i\mathrm{}`$. Every such a homogeneous relation generates a relation for $`𝔹_{s_i,V}(\lambda ),𝔹_{s_j,V}(\lambda )`$. Namely, for $`m=2`$, the relation is
$`(s_j)^1(𝔹_{s_i,V}(s_j(\lambda )))𝔹_{s_j,V}(\lambda )=(s_i)^1(𝔹_{s_j,V}(s_i(\lambda )))𝔹_{s_i,V}(\lambda ),`$
for $`m=3`$, the relation is
$`(s_js_i)^1(𝔹_{s_i,V}((s_js_i)(\lambda )))(s_i)^1(𝔹_{s_j,V}(s_i(\lambda )))𝔹_{s_i,V}(\lambda )=`$
$`(s_is_j)^1(𝔹_{s_j,V}((s_is_j)(\lambda )))(s_j)^1(𝔹_{s_i}(s_j(\lambda )))B_{s_j}(\lambda ),`$
and so on.
These relations can be written in terms of operators $`𝔹_V^\alpha (\lambda )`$.
1. For $`\alpha ,\beta \mathrm{\Sigma }`$, denote $`\alpha ,\beta `$ the subspace $`\alpha +\beta 𝔥`$. Then
$`𝔹_V^\alpha (\lambda )𝔹_V^\beta (\lambda )`$ $`=`$ $`𝔹_V^\beta (\lambda )𝔹_V^\alpha (\lambda ),`$
$`𝔹_V^\alpha (\lambda )𝔹_V^{\alpha +\beta }(\lambda )𝔹_V^\beta (\lambda )`$ $`=`$ $`𝔹_V^\beta (\lambda )𝔹_V^{\alpha +\beta }(\lambda )𝔹_V^\alpha (\lambda ),`$
$`𝔹_V^\alpha (\lambda )𝔹_V^{\alpha +\beta }(\lambda )𝔹_V^{\alpha +2\beta }(\lambda )𝔹_V^\beta (\lambda )`$ $`=`$ $`𝔹_V^\beta (\lambda )𝔹_V^{\alpha +2\beta }(\lambda )𝔹_V^{\alpha +\beta }(\lambda )𝔹_V^\alpha (\lambda ),`$
$`𝔹_V^\alpha (\lambda )𝔹_V^{3\alpha +\beta }(\lambda )𝔹_V^{2\alpha +\beta }(\lambda )𝔹_V^{3\alpha +2\beta }(\lambda )𝔹_V^{\alpha +\beta }(\lambda )𝔹_V^\beta (\lambda )=`$
$`𝔹_V^\beta (\lambda )𝔹_V^{\alpha +\beta }(\lambda )𝔹_V^{3\alpha +2\beta }(\lambda )𝔹_V^{2\alpha +\beta }(\lambda )𝔹_V^{3\alpha +\beta }(\lambda )𝔹_V^\alpha (\lambda )`$
under the assumption that $`\alpha ,\beta =\{\pm \gamma \}`$ where $`\gamma `$ runs over all indices in the corresponding identity.
2. $`𝔹_{w,WV}(\lambda ))=`$ $`x_w^1(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))x_w`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))J(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1`$
###### Lemma 16.
Let $`W,V`$ be finite dimensional $`𝔤`$-modules, $`\lambda 𝔥`$, $`w𝕎`$. Then
$`\mathrm{\Omega }𝔹_{w,WV}(\lambda )=𝔹_{w,WV}(\lambda )\mathrm{\Omega }`$
and
$`(w^1(\mathrm{\Omega }^{})+\lambda ^{(2)})𝔹_{w,WV}(\lambda )=𝔹_{w,WV}(\lambda )(\mathrm{\Omega }^{}+\lambda ^{(2)}).`$
###### Proof.
The first equation holds since $`\mathrm{\Omega }`$ commutes with the comultiplication. Now
$`𝔹_{w,WV}(\lambda )(\mathrm{\Omega }^{}+\lambda ^{(2)})=x_w^1(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))x_w`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))J_{WV}(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1(\mathrm{\Omega }^{}+\lambda ^{(2)})=`$
$`x_w^1(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))x_w`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))(\mathrm{\Omega }^0+\lambda ^{(2)})J_{WV}(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1=`$
$`x_w^1(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))x_w(\mathrm{\Omega }^0+\lambda ^{(2)})`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))J_{WV}(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1=`$
$`x_w^1(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))(\mathrm{\Omega }^0+(w(\lambda ))^{(2)})x_w`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))J_{WV}(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1=`$
$`x_w^1(\mathrm{\Omega }^{}+(w(\lambda ))^{(2)})(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))x_w`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))J_{WV}(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1=`$
$`(w^1(\mathrm{\Omega }^{})+\lambda ^{(2)})x_w^1(J_{WV}(w(\lambda )\rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)})))x_w`$
$`(𝔹_{w,W}(\lambda {\displaystyle \frac{1}{2}}h^{(2)})𝔹_{w,V}(\lambda +{\displaystyle \frac{1}{2}}h^{(1)}))J_{WV}(\lambda \rho +{\displaystyle \frac{1}{2}}(h^{(1)}+h^{(2)}))^1=`$
$`(w^1(\mathrm{\Omega }^{})+\lambda ^{(2)})𝔹_{w,WV}(\lambda ).`$
## 3. Difference Equations Compatible with KZ Equations for $`𝔤=sl_N`$
### 3.1. Statement of Results
Let $`e_{i,j}`$, $`i,j=1,\mathrm{}N`$, be the standard generators of the Lie algebra $`gl_N`$,
$$[e_{i,j},e_{k,l}]=\delta _{j,k}e_{i,l}\delta _{i,l}e_{j,k}.$$
$`sl_N`$ is the Lie subalgebra of $`gl_N`$ such that $`sl_n=𝔫_+𝔥𝔫_{}`$ where
$$𝔫_+=_{1i<jN}e_{i,j},𝔫_{}=_{1j<iN}e_{i,j},$$
and $`𝔥=\{\lambda =_{i=1}^N\lambda _ie_{i,i}|\lambda _i,_{i=1}^N\lambda _i=0\}`$.
The invariant scalar product is defined by $`(e_{i,j},e_{k,l})=\delta _{i,l}\delta _{j,k}`$. The roots are $`e_{i,i}e_{j,j}`$ for $`ij`$. $`\alpha ^{}=\alpha `$ for any root. For a root $`\alpha =e_{i,i}e_{j,j}`$, we have $`H_\alpha =e_{i,i}e_{j,j},E_\alpha =e_{i,j},F_\alpha =e_{j,i}`$. The simple roots are $`\alpha _i=e_{i,i}e_{i+1,i+1}`$ for $`i=1,\mathrm{},N1`$. $`𝕎`$ is the symmetric group $`S^N`$ permutting coordinates of $`\lambda 𝔥`$. The (dual) fundamental weights are $`\omega _i=\omega _i^{}=_{j=1}^i(1\frac{i}{N})e_{j,j}_{j=i+1}^N\frac{i}{N}e_{j,j}`$ for $`i=1,\mathrm{},N1`$. All dual fundamental weights are minuscule. For $`i=1,\mathrm{},N1`$, the permutation $`w_{[i]}^1S^N`$ is $`({}_{i+1}{}^{1}{}_{i+2}{}^{2}{}_{\mathrm{}}{}^{\mathrm{}}{}_{N}{}^{Ni}{}_{1}{}^{Ni+1}{}_{\mathrm{}}{}^{\mathrm{}}{}_{i}{}^{N})`$.
For any finite dimensional $`sl_N`$-module $`V`$ and $`wS^N`$ consider the operators $`𝔹_{w,V}(\lambda ):VV`$.
Let $`V=V_1\mathrm{}V_n`$ be a tensor product of finite dimensional $`sl_N`$-modules. For $`\kappa `$ and $`\lambda 𝔥`$, consider the trigonometric KZ equations with values in $`V`$,
(9) $`_j(\lambda ,\kappa )u(z_1,\mathrm{},z_n,\lambda )=\mathrm{\hspace{0.17em}0},j=1,\mathrm{},n.`$
Here $`u(z_1,\mathrm{},z_n,\lambda )V`$ is a function of complex variables $`z_1,\mathrm{},z_n`$ and $`\lambda 𝔥`$.
Introduce the dynamical difference equations on a $`V`$-valued function $`u(z_1,\mathrm{},z_n,\lambda )`$ as
(10)
$`u(z_1,\mathrm{},z_n,\lambda +\kappa \omega _i^{})=K_i(z_1,\mathrm{},z_n,\lambda )u(z_1,\mathrm{},z_n,\lambda ),i=1,\mathrm{},N1`$
where
$$K_i(z_1,\mathrm{},z_n,\lambda )=\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}𝔹_{w_{[i]},V}(\lambda ).$$
The operator $`_{k=1}^nz_k^{(\omega _i^{})^{(k)}}`$ is well defined if the argument of $`z_1,\mathrm{},z_n`$ is fixed. The dynamical difference equations are well defined on functions of $`(z,\lambda )`$ where $`\lambda 𝔥`$ and $`z`$ belongs to the universal cover of $`(^{})^n`$. Notice that the KZ equations are well defined for $`V`$-valued functions of the same variables.
The KZ operators $`_j(\lambda ,\kappa )`$ and the operators $`K_i(z_1,\mathrm{},z_n,\lambda )`$ preserve the weight decomposition of $`V`$.
###### Theorem 17.
The dynamical equations (10) together with the KZ equations (9) form a compatible system of equations.
### 3.2. Proof
First prove that
$$\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}𝔹_{w_{[i]},V}(\lambda )_j(\lambda ,\kappa )=_j(\lambda +\kappa \omega _i^{},\kappa )\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}𝔹_{w_{[i]},V}(\lambda )$$
for all $`i`$ and $`j`$. Multiplying both sides from the left by $`_{k=1}^nz_k^{(\omega _i^{})^{(k)}}`$ and using Lemma 3, we reduce the equation to
$`𝔹_{w_{[i]},V}(\lambda )({\displaystyle \underset{k,kj}{}}r(z_j/z_k)^{(j,k)}+\lambda ^{(j)})=({\displaystyle \underset{k,kj}{}}w_{[i]}^1(r(z_j/z_k))^{(j,k)}+\lambda ^{(j)})𝔹_{w_{[i]},V}(\lambda ).`$
###### Lemma 18.
For $`j=1,\mathrm{},n`$ and $`w𝕎`$, we have
$`𝔹_{w,V}(\lambda )({\displaystyle \underset{k,kj}{}}r(z_j/z_k)^{(j,k)}+\lambda ^{(j)})=({\displaystyle \underset{k,kj}{}}w^1(r(z_j/z_k))^{(j,k)}+\lambda ^{(j)})𝔹_{w,V}(\lambda ).`$
###### Proof.
It is sufficient to check the equation for the residues of both sides at $`z_j=z_k,kj`$, and for the limit of both sides as $`z_j\mathrm{}`$. The residue equation $`[𝔹_{w,V}(\lambda ),\mathrm{\Omega }^{(j,k)}]=0`$ is true since the Casimir operator commutes with the comultiplication. The limit equation
$`𝔹_{w,V}(\lambda )({\displaystyle \underset{k,kj}{}}(\mathrm{\Omega }^+)^{(j,k)}+\lambda ^{(j)})=({\displaystyle \underset{k,kj}{}}w^1(\mathrm{\Omega }^+)^{(j,k)}+\lambda ^{(j)})𝔹_{w_{[i]},V}(\lambda )`$
is a corollary of Lemma 16. ∎
The Theorem is proved for $`sl_N,N=2`$. For $`N>2`$, it remains to prove that
(11) $`K_i(z,\lambda +\kappa \omega _j^{})K_j(z,\lambda )=K_j(z,\lambda +\kappa \omega _i^{})K_i(z,\lambda )`$
for all $`i,j`$, $`0<i<j<N`$. We prove (11) for $`N=3`$. For arbitrary $`N`$ the proof is similar. Another proof see in Section 4. For $`N=3`$, $`i=1,j=2`$, equation (11) takes the form
(12) $`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _1^{})^{(k)}}𝔹_V^{\alpha _1+\alpha _2}(\lambda +\kappa \omega _2^{})𝔹_V^{\alpha _1}(\lambda +\kappa \omega _2^{}){\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _2^{})^{(k)}}𝔹_V^{\alpha _1+\alpha _2}(\lambda )𝔹_V^{\alpha _2}(\lambda )=`$
$`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _2^{})^{(k)}}𝔹_V^{\alpha _1+\alpha _2}(\lambda +\kappa \omega _1^{})𝔹_V^{\alpha _2}(\lambda +\kappa \omega _1^{}){\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _1^{})^{(k)}}𝔹_V^{\alpha _1+\alpha _2}(\lambda )𝔹_V^{\alpha _1}(\lambda ).`$
We have $`𝔹_V^{\alpha _1}(\lambda +\kappa \omega _2^{})=𝔹_V^{\alpha _1}(\lambda )`$ since $`(\omega _2^{},\alpha _1)=0`$. We have $`[𝔹_V^{\alpha _1}(\lambda ),_{k=1}^nz_k^{(\omega _2^{})^{(k)}}]=0`$ since $`𝔹_V^{\alpha _1}(\lambda )`$ is a power series in $`E_{\alpha _1},F_{\alpha _1}`$. Similarly, $`𝔹_V^{\alpha _2}(\lambda +\kappa \omega _1^{})=𝔹_V^{\alpha _2}(\lambda )`$ and $`[𝔹_V^{\alpha _2}(\lambda ),_{k=1}^nz_k^{(\omega _1^{})^{(k)}}]=0`$. Using these remarks and the relation
$$𝔹_V^{\alpha _2}(\lambda )𝔹_V^{\alpha _1+\alpha _2}(\lambda )𝔹_V^{\alpha _1}(\lambda )=𝔹_V^{\alpha _1}(\lambda )𝔹_V^{\alpha _1+\alpha _2}(\lambda )𝔹_V^{\alpha _2}(\lambda )$$
we reduce (12) to
$`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _1^{}\omega _2^{})^{(k)}}𝔹_V^{\alpha _1+\alpha _2}(\lambda +\kappa \omega _2^{})=𝔹_V^{\alpha _1+\alpha _2}(\lambda +\kappa \omega _1^{}){\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _1^{}\omega _2^{})^{(k)}}.`$
This equation holds since $`𝔹_V^{\alpha _1+\alpha _2}(\lambda +\kappa \omega _2^{})=𝔹_V^{\alpha _1+\alpha _2}(\lambda +\kappa \omega _1^{})`$, each of these operators is a power series in $`E_{\alpha _1+\alpha _2},F_{\alpha _1+\alpha _2}`$, and $`(\omega _1^{}\omega _2^{},\alpha _1+\alpha _2)=0`$.
### 3.3. An Equivalent Form of Dynamical Equations for $`sl_N`$
For $`j=1,\mathrm{},N`$, set $`\delta _j=\omega _j^{}\omega _{j1}^{}`$ where $`\omega _0^{}=\omega _N^{}=0`$. Then the system of equations (10) is equivalent to the system
$`u(z_1,\mathrm{},z_n,\lambda +\kappa \delta _i)=`$ $`\left(𝔹_V^{e_{i1,i1}e_{i,i}}(\lambda +\kappa \delta _i)\right)^1\mathrm{}\left(𝔹_V^{e_{1,1}e_{i,i}}(\lambda +\kappa \delta _i)\right)^1\times `$
$`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\delta _i)^{(k)}}𝔹_V^{e_{i,i}e_{n,n}}(\lambda )\mathrm{}𝔹_V^{e_{i,i}e_{i+1,i+1}}(\lambda )u(z_1,\mathrm{},z_n,\lambda )`$
where $`i=1,\mathrm{},N`$.
Notice that the inverse powers can be eliminated using property III in Section 2.7.
### 3.4. Application to Determinants
Let $`𝔤`$ be a simple Lie algebra, $`V`$ a finite dimensional $`𝔤`$-module, $`V[\nu ]`$ a weight subspace. For a positive root $`\alpha `$ fix the $`sl_2`$ subalgebra in $`𝔤`$ generated by $`H_\alpha ,E_\alpha ,F_\alpha `$. Consider $`V`$ as an $`sl_2`$-module. Let $`V[\nu ]_\alpha V`$ be the $`sl_2`$-submodule generated by $`V[\nu ]`$,
$$V[\nu ]_\alpha =_{k_0}W_k^\alpha L_{\nu +k\alpha }$$
the decomposition into irreducible $`sl_2`$-modules. Here $`L_{\nu +k\alpha }`$ is the irreducible module with highest weight $`\nu +k\alpha `$ and $`W_k^\alpha `$ the multiplicity space. Let $`d_k^\alpha =`$ dim $`W_k^\alpha `$. Set
$`X_{\alpha ,V[\nu ]}(\lambda )={\displaystyle \underset{k_0}{}}\left({\displaystyle \underset{j=1}{\overset{k}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(1\frac{(\lambda \frac{1}{2}(\nu +j\alpha ),\alpha )}{\kappa }\right)}{\mathrm{\Gamma }\left(1\frac{(\lambda +\frac{1}{2}(\nu +j\alpha ),\alpha )}{\kappa }\right)}}\right)^{d_k^\alpha },`$
cf. formula (8). Here $`\mathrm{\Gamma }`$ is the standard gamma function.
Let $`V=V_1\mathrm{}V_n`$ be a tensor product of finite dimensional $`𝔤`$-modules. Set $`\mathrm{\Lambda }_k(\lambda )=\text{tr}_{V[\nu ]}\lambda ^{(k)}`$, $`ϵ_{k,l}=\text{tr}_{V[\nu ]}\mathrm{\Omega }^{(k,l)}`$, $`\gamma _k=_{l,lk}\epsilon _{k,l}`$. Set
(13) $`D_{V[\nu ]}(z_1,\mathrm{},z_n,\lambda )={\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{\frac{\mathrm{\Lambda }_k(\lambda )}{\kappa }\frac{\gamma _k}{2\kappa }}{\displaystyle \underset{1k<ln}{}}(z_kz_l)^{\frac{ϵ_{k,l}}{\kappa }}{\displaystyle \underset{\alpha \mathrm{\Sigma }_+}{}}X_{\alpha ,V[\nu ]}(\lambda ).`$
Let $`𝔤=sl_N`$, $`V=V_1\mathrm{}V_n`$ a tensor product of finite dimensional $`sl_N`$-modules. Fix a basis $`v_1,\mathrm{},v_d`$ in a weight subspace $`V[\nu ]`$. Suppose that $`u_i(z_1,\mathrm{},z_n,\lambda )=_{j=1}^du_{i,j}v_j`$, $`i=1,\mathrm{},d`$, is a set of $`V[\nu ]`$-valued solutions of the combined system of KZ equations (9) and dynamical equations (10).
###### Corollary 19.
$$\text{det}(u_{i,j})_{1i,jd}=C_{V[\nu ]}(\lambda )D_{V[\nu ]}(z_1,\mathrm{},z_n,\lambda )$$
where $`C_{V[\nu ]}(\lambda )`$ is a function of $`\lambda `$ (depending also on $`V_1,\mathrm{},V_n`$ and $`\nu `$) such that
$$C_{V[\nu ]}(\lambda +\kappa \omega )=C_{V[\nu ]}(\lambda )$$
for all $`\omega P^{}`$.
###### Proof.
The Corollary follows from the following simple Lemma.
###### Lemma 20.
For $`i=1,\mathrm{},N1`$, the operator $`𝔹_{w_{[i]},V}(\lambda )`$ is the product in a suitable order of all operators $`𝔹_V^\alpha (\lambda )`$ with $`\alpha \mathrm{\Sigma }_+`$ and $`(\omega _i^{},\alpha )>0`$.
Notice that Lemma 20 in particular implies that operators $`𝔹_{w_{[i]},V}(\lambda )`$ and the dynamical equations are well defined in the tensor product of any highest weight $`sl_N`$-modules.
## 4. Dynamical Difference Equations
In this section we introduce dynamical difference equations for arbitrary simple Lie algebra. The compatibility of the dynamical equations follows from \[Ch1\]. We prove the compatibility of dynamical and KZ equations.
### 4.1. Affine Root Systems, \[Ch1, Ch2\]
Let $`𝔤`$ be a simple Lie algebra. The vectors $`\stackrel{~}{\alpha }=[\alpha ,j]𝔥\times `$ for $`\alpha \mathrm{\Sigma },j`$ form the affine root system $`\mathrm{\Sigma }^a`$ corresponding to the root system $`\mathrm{\Sigma }𝔥`$. We view $`\mathrm{\Sigma }`$ as a subset in $`\mathrm{\Sigma }^a`$ identifying $`\alpha 𝔥`$ with $`[\alpha ,0]`$. The simple roots of $`\mathrm{\Sigma }^a`$ are $`\alpha _1,\mathrm{},\alpha _r\mathrm{\Sigma }`$ and $`\alpha _0=[\theta ,1]`$ where $`\theta \mathrm{\Sigma }`$ is the maximal root. The positive roots are $`\mathrm{\Sigma }_+^a=\{[\alpha ,j]\mathrm{\Sigma }^a|\alpha \mathrm{\Sigma },j>0\text{or}\alpha \mathrm{\Sigma }_+,j=0\}`$. The Dynkin diagram and its affine completion with $`\{\alpha _i\}_{0in}`$ as vertices are denoted $`\mathrm{\Gamma }`$ and $`\mathrm{\Gamma }^a`$, respectively. The set of the indices of the images of $`\alpha _0`$ with respect to all authomorphisms of $`\mathrm{\Gamma }^a`$ is denoted $`O`$ ($`O=\{0\}`$ for $`E_8,F_4,G_2`$ ). Let $`O^{}=\{iO|i0\}`$. For $`i=1,\mathrm{},r`$, the dual fundamental weight $`\omega _i^{}`$ is minuscule if and only if $`iO^{}`$.
Given $`\stackrel{~}{\alpha }=[\alpha ,j]\mathrm{\Sigma }^a`$ and $`\omega P^{}`$, set
$$s_{\stackrel{~}{\alpha }}(\stackrel{~}{z})=\stackrel{~}{z}(z,\alpha ^{})\stackrel{~}{\alpha },t_\omega (\stackrel{~}{z})=[z,\xi (z,\omega )]$$
for $`\stackrel{~}{z}=[z,\xi ]`$.
The affine Weyl group $`𝕎^a`$ is the group generated by reflections $`s_{\stackrel{~}{\alpha }},\stackrel{~}{\alpha }\mathrm{\Sigma }_+^a`$. One defines the length of elements of $`𝕎^a`$ taking the simple reflections $`s_i=s_{\alpha _i},i=0,\mathrm{},r`$, as generators of $`𝕎^a`$. The group $`𝕎^a`$ is the semidirect product $`𝕎Q_t^{}`$ of its subgroups $`𝕎=s_\alpha |\alpha \mathrm{\Sigma }_+`$ and $`Q_t^{}=\{t_\omega |\omega Q^{}\}`$, where for $`\alpha \mathrm{\Sigma }`$ we have $`t_\alpha ^{}=s_\alpha s_{[\alpha ,1]}=s_{[\alpha ,1]}s_\alpha `$.
Consider the group $`P_t^{}=\{t_\omega |\omega P^{}\}`$. The extended affine Weyl group $`𝕎^b`$ is the group of transformations of $`𝔥\times `$ generated by $`𝕎`$ and $`P_t^{}`$. $`𝕎^b`$ is isomorphic to $`𝕎P_t^{}`$ with action $`(w,\omega )([z,\xi ])=[w(z),\xi (z,\omega )]`$.
Notice that for any $`w𝕎^b`$ and $`\stackrel{~}{\alpha }\mathrm{\Sigma }^a`$, we have $`w(\stackrel{~}{\alpha })\mathrm{\Sigma }^a`$.
The extended affine Weyl group has a remarkable subgroup $`\mathrm{\Pi }=\{\pi _i|iO\}`$, where $`\pi _0\mathrm{\Pi }`$ is the identity element in $`𝕎^b`$ and for $`iO^{}`$ we have $`\pi _i=t_{\omega _i^{}}w_{[i]}^1`$. The group $`\mathrm{\Pi }`$ is isomorphic to $`P^{}/Q^{}`$ with the isomorphism sending $`\pi _i`$ to the minuscle weight $`\omega _i^{}`$. For $`iO^{}`$, the element $`w_{[i]}`$ preserves the set $`\{\theta ,\alpha _1,\mathrm{},\alpha _r\}`$ and $`\pi _i(\alpha _0)=\alpha _i=w_{[i]}^1(\theta )`$. We have
$$𝕎^b=\mathrm{\Pi }𝕎^a,\text{where}\pi _is_l\pi _i^1=s_k\text{if}\pi _i(\alpha _l)=\alpha _k\text{and}0kr.$$
We extend the notion of length to $`𝕎^b`$. For $`iO^{},w𝕎^a`$, we set the length of $`\pi _iw`$ to be equal to the length of $`w`$ in $`𝕎^a`$.
### 4.2. Affine R-matrices, \[Ch1, Ch2\]
Fix a $``$-algebra $`F`$. A set $`G=\{G^\alpha F|\alpha \mathrm{\Sigma }\}`$ is called a closed R-matrix if
$`G^\alpha G^\beta `$ $`=`$ $`G^\beta G^\alpha ,`$
$`G^\alpha G^{\alpha +\beta }G^\beta `$ $`=`$ $`G^\beta G^{\alpha +\beta }G^\alpha ,`$
$`G^\alpha G^{\alpha +\beta }G^{\alpha +2\beta }G^\beta `$ $`=`$ $`G^\beta G^{\alpha +2\beta }G^{\alpha +\beta }G^\alpha ,`$
$`G^\alpha G^{3\alpha +\beta }G^{2\alpha +\beta }G^{3\alpha +2\beta }G^{\alpha +\beta }G^\beta `$ $`=`$ $`G^\beta G^{\alpha +\beta }G^{3\alpha +2\beta }G^{2\alpha +\beta }G^{3\alpha +\beta }G^\alpha `$
under the assumption that $`\alpha ,\beta \mathrm{\Sigma }`$ and $`\alpha ,\beta =\{\pm \gamma \}`$ where $`\gamma `$ runs over all indices in the corresponding identity.
A set $`G^a=\{\stackrel{~}{G}^{\stackrel{~}{\alpha }}F|\stackrel{~}{\alpha }\mathrm{\Sigma }^a\}`$ is called a closed affine R-matrix if $`\stackrel{~}{G}^{\stackrel{~}{\alpha }}`$ satisfy the same relations where $`\alpha ,\beta `$ are replaced with $`\stackrel{~}{\alpha },\stackrel{~}{\beta }`$.
If $`G^a`$ is an affine R-matrix, for any $`w𝕎^b`$ define an element $`\stackrel{~}{G}_wF`$ as follows. Given a reduced presentation $`w=\pi _is_{j_l}\mathrm{}s_{j_1}`$, $`iO`$, $`0j_1,\mathrm{},j_lr`$, set $`\stackrel{~}{G}_w=\stackrel{~}{G}^{\stackrel{~}{\alpha }^l}\mathrm{}\stackrel{~}{G}^{\stackrel{~}{\alpha }^1}`$ where $`\stackrel{~}{\alpha }^1=\alpha _{j_1},\stackrel{~}{\alpha }^2=s_{j_1}(\alpha _{j_2}),\stackrel{~}{\alpha }^3=s_{j_1}s_{j_2}(\alpha _{j_3})`$,… The element $`\stackrel{~}{G}_w`$ does not depend on the reduced presentation of $`w`$. We set $`\stackrel{~}{G}_{\text{id}}=1`$.
The unordered set $`\{\stackrel{~}{\alpha }^1,\mathrm{},\stackrel{~}{\alpha }^l\}`$ is denoted $`\stackrel{~}{A}(w)`$. There is a useful formula valid for any (not necessarily minuscule) dual fundamental weight $`\omega _i^{}`$, $`i=1,\mathrm{},r`$,
(14) $`\stackrel{~}{A}(t_{\omega _i^{}})=\{[\alpha ,j]|\alpha \mathrm{\Sigma }_+,\text{and}(\omega _i^{},\alpha )>j0\},`$
Prop. 1.4 \[Ch2\].
Introduce the following formal notation: for $`w𝕎^b`$, $`\stackrel{~}{\alpha },\stackrel{~}{\beta }\mathrm{\Sigma }^a`$, set $`{}_{}{}^{w}(\stackrel{~}{G}^{\stackrel{~}{\alpha }})=G^{w(\stackrel{~}{\alpha })},{}_{}{}^{w}(\stackrel{~}{G}^{\stackrel{~}{\alpha }}\stackrel{~}{G}^{\stackrel{~}{\beta }})=G^{w(\stackrel{~}{\alpha })}G^{w(\stackrel{~}{\beta })}`$,… Then the elements $`\{\stackrel{~}{G}_w|w𝕎^b\}`$ form a 1-cocycle:
$$\stackrel{~}{G}_{xy}={}_{}{}^{y^1}\stackrel{~}{G}_{x}^{}\stackrel{~}{G}_y$$
for all $`x,y𝕎^b`$ such that $`l(xy)=l(x)+l(y)`$.
There is a way to construct a closed affine R-matrix if a closed nonaffine R-matrix $`G=\{G^\alpha F|\alpha \mathrm{\Sigma }\}`$ is given. Namely, assume that the group $`P_t^{}`$ acts on the algebra $`F`$ so that $`{}_{}{}^{t_\omega }(G^\alpha )=G^\alpha `$ whenever $`(\omega ,\alpha )=0`$, $`\omega P^{}`$, $`\alpha \mathrm{\Sigma }`$. Then for $`\stackrel{~}{\alpha }=[\alpha ,j]\mathrm{\Sigma }^a`$, choose $`\omega P^{}`$ so that $`(\omega ,\alpha )=j`$ and set $`\stackrel{~}{G}^{\stackrel{~}{\alpha }}={}_{}{}^{t_\omega }(G^\alpha )`$. The set $`G^a=\{\stackrel{~}{G}^{\stackrel{~}{\alpha }}F|\stackrel{~}{\alpha }\mathrm{\Sigma }^a\}`$ is well defined and forms a closed affine R-matrix called the affine completion of the R-matrix $`G`$.
Assume that a closed affine R-matrix $`G^a`$ is the affine completion of a closed nonaffine R-matrix $`G`$. Consider the system of equations for an element $`\mathrm{\Phi }F`$:
(15) $`{}_{}{}^{t_{\omega _i^{}}}(\mathrm{\Phi })=\stackrel{~}{G}_{t_{\omega _i^{}}}\mathrm{\Phi },i=1,\mathrm{},r,`$
where $`\omega _1^{},\mathrm{},\omega _r^{}`$ are the dual fundamental weights.
###### Theorem 21.
\[Ch1\] The system of equations (15) is compatible,
$`{}_{}{}^{t_{\omega _i^{}}}(\stackrel{~}{G}_{t_{\omega _j^{}}})\stackrel{~}{G}_{t_{\omega _i^{}}}={}_{}{}^{t_{\omega _j^{}}}(\stackrel{~}{G}_{t_{\omega _i^{}}})\stackrel{~}{G}_{t_{\omega _j^{}}}`$
for $`1i<jr`$.
Example, \[Ch1\]. Let $`\alpha =\alpha _1,\beta =\alpha _2,a=\omega _1^{},b=\omega _2^{}`$. Then the system for $`A_2`$ is
$`{}_{}{}^{t_a}(\mathrm{\Phi })=\stackrel{~}{G}^{\alpha +\beta }\stackrel{~}{G}^\alpha \mathrm{\Phi },{}_{}{}^{t_b}(\mathrm{\Phi })=\stackrel{~}{G}^{\alpha +\beta }\stackrel{~}{G}^\beta \mathrm{\Phi }.`$
The system for $`B_2`$ is
$`{}_{}{}^{t_a}(\mathrm{\Phi })=\stackrel{~}{G}^{\alpha +2\beta }\stackrel{~}{G}^{\alpha +\beta }\stackrel{~}{G}^\alpha \mathrm{\Phi },{}_{}{}^{t_b}(\mathrm{\Phi })=\stackrel{~}{G}^{[\alpha +2\beta ,1]}\stackrel{~}{G}^{\alpha +\beta }\stackrel{~}{G}^{\alpha +2\beta }\stackrel{~}{G}^\beta \mathrm{\Phi }.`$
The system for $`G_2`$ is
$`{}_{}{}^{t_a}(\mathrm{\Phi })`$ $`=`$ $`\stackrel{~}{G}^{[3\alpha +2\beta ,2]}\stackrel{~}{G}^{[3\alpha +\beta ,2]}\stackrel{~}{G}^{[2\alpha +\beta ,1]}\stackrel{~}{G}^{[3\alpha +2\beta ,1]}\stackrel{~}{G}^{[3\alpha +\beta ,1]}\times `$
$`\stackrel{~}{G}^{\alpha +\beta }\stackrel{~}{G}^{3\alpha +2\beta }\stackrel{~}{G}^{2\alpha +\beta }\stackrel{~}{G}^{3\alpha +\beta }\stackrel{~}{G}^\alpha \mathrm{\Phi },`$
$`{}_{}{}^{t_b}(\mathrm{\Phi })`$ $`=`$ $`\stackrel{~}{G}^{[3\alpha +2\beta ,1]}\stackrel{~}{G}^{3\alpha +\beta }\stackrel{~}{G}^{2\alpha +\beta }\stackrel{~}{G}^{3\alpha +2\beta }\stackrel{~}{G}^{\alpha +\beta }\stackrel{~}{G}^\beta \mathrm{\Phi }.`$
### 4.3. Affine R-matrix for Dynamical Equations
Fix $`\kappa `$ and a natural number $`n`$. Let $`F`$ be the algebra of meromorphic functions of $`z_1,\mathrm{},z_n`$ and $`\lambda 𝔥`$ with values in $`U𝔤_0^n`$. Define an action of $`𝕎`$ on $`F`$ by
$${}_{}{}^{w}f(z_1,\mathrm{},z_n,\lambda )=w(f(z_1,\mathrm{},z_n,w^1(\lambda )))$$
and an action of $`P_t^{}`$ on $`F`$ by
$${}_{}{}^{t_\omega }f(z_1,\mathrm{},z_n,\lambda )=\underset{k=1}{\overset{n}{}}z_k^{\omega ^{(k)}}f(z_1,\mathrm{},z_n,\lambda \kappa \omega )\underset{k=1}{\overset{n}{}}z_k^{\omega ^{(k)}}$$
where $`w𝕎`$, $`\omega P^{}`$, $`fF`$.
###### Lemma 22.
Those actions extend to an action of $`𝕎^b=𝕎P_t^{}`$ on $`F`$, i.e. $`{}_{}{}^{w}({}_{}{}^{t_\omega }f)={}_{}{}^{t_{w(\omega )}}({}_{}{}^{w}f)`$ for $`w𝕎`$, $`\omega P^{}`$, $`fF`$. $`\mathrm{}`$
Define a closed nonaffine $`F`$-valued R-matrix $`G_F=\{G_F^\alpha |\alpha \mathrm{\Sigma }\}`$ by
$$G_F^\alpha (\lambda )=\mathrm{\Delta }^{(n)}(p((\lambda ,\alpha ^{})1;H_\alpha ,E_\alpha ,F_\alpha )).$$
Properties of operators $`𝔹_V^\alpha `$ described in Section 2.7 ensure that $`G_F`$ is a closed R-matrix. The action of $`P_t^{}`$ on $`F`$ defined above clearly has the property: $`{}_{}{}^{t_\omega }(G_F^\alpha )=G_F^\alpha `$ whenever $`(\omega ,\alpha )=0`$, $`\omega P^{}`$, $`\alpha \mathrm{\Sigma }`$. This allows us to define a closed affine R-matrix $`G_F^a=\{\stackrel{~}{G}_F^{\stackrel{~}{\alpha }}F|\stackrel{~}{\alpha }\mathrm{\Sigma }^a\}`$ as the affine completion of the R-matrix $`G_F`$. Namely, for $`\stackrel{~}{\alpha }=[\alpha ,j]\mathrm{\Sigma }^a`$, we choose $`\omega P^{}`$ so that $`(\omega ,\alpha )=j`$ and set
$$\stackrel{~}{G}_F^{[\alpha ,j]}(z_1,\mathrm{},z_n,\lambda )={}_{}{}^{t_\omega }(G_F^\alpha )=\underset{k=1}{\overset{n}{}}z_k^{\omega ^{(k)}}G_F^\alpha (\lambda \kappa \omega )\underset{k=1}{\overset{n}{}}z_k^{\omega ^{(k)}}.$$
Let $`V=V_1\mathrm{}V_n`$ be a tensor product of finite dimensional $`𝔤`$-modules. Let $`F_V`$ be the algebra of meromorphic functions of $`z_1,\mathrm{},z_n`$ and $`\lambda 𝔥`$ with values in
End $`(V)`$. The closed affine R-matrix $`G_F^a`$ induces a closed affine R-matrix $`G_V^a=\{\stackrel{~}{G}_V^{\stackrel{~}{\alpha }}\}`$ where
$$\stackrel{~}{G}_V^{\stackrel{~}{\alpha }}(z_1,\mathrm{},z_n,\lambda )=\stackrel{~}{G}_F^{\stackrel{~}{\alpha }}(z_1,\mathrm{},z_n,\lambda +\frac{1}{2}\underset{k=1}{\overset{n}{}}h^{(k)})|_V.$$
In other words,
$$\stackrel{~}{G}_V^{[\alpha ,j]}(z_1,\mathrm{},z_n,\lambda )=\underset{k=1}{\overset{n}{}}z_k^{\omega ^{(k)}}𝔹_V^\alpha (\lambda \kappa \omega )\underset{k=1}{\overset{n}{}}z_k^{\omega ^{(k)}}$$
where $`(\omega ,\alpha )=j`$ and the operators $`𝔹_V^\alpha `$ are defined in Section 2.7. For any $`w𝕎^b`$ and $`\stackrel{~}{\alpha }\mathrm{\Sigma }^a`$, we have $`{}_{}{}^{w}(\stackrel{~}{G}_V^{\stackrel{~}{\alpha }})=\stackrel{~}{G}_V^{w(\stackrel{~}{\alpha })}`$.
Let $`\{\stackrel{~}{G}_w^VF_V|w𝕎^b\}`$ be the 1-cocycle associated with the affine R-matrix $`G_V^a`$. Consider the system
$$\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}\mathrm{\Phi }(z_1,\mathrm{},z_n,\lambda +\kappa \omega _i^{})\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}=\stackrel{~}{G}_{t_{\omega _i^{}}}^V(z_1,\mathrm{},z_n,\lambda )\mathrm{\Phi }(z_1,\mathrm{},z_n,\lambda ),$$
$`i=1,\mathrm{},r`$, of equations (15) associated with the affine R-matrix $`G_V^a`$. By Theorem 21 this system is compatible.
Example. For $`𝔤=sl_N`$, this system of equations for an element $`\mathrm{\Phi }F_V`$ has the form
$$\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}\mathrm{\Phi }(z_1,\mathrm{},z_n,\lambda +\kappa \omega _i^{})\underset{k=1}{\overset{n}{}}z_k^{(\omega _i^{})^{(k)}}=𝔹_{w_{[i]},V}(\lambda )\mathrm{\Phi }(z_1,\mathrm{},z_n,\lambda ),$$
$`i=1,\mathrm{},N1`$, cf. (10).
Introduce the dynamical difference equations on a $`V`$-valued function $`u(z_1,\mathrm{},z_n,\lambda )`$ as
(16)
$`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _i^{})^{(k)}}u(z_1,\mathrm{},z_n,\lambda +\kappa \omega _i^{})=\stackrel{~}{G}_{t_{\omega _i^{}}}^V(z_1,\mathrm{},z_n,\lambda )u(z_1,\mathrm{},z_n,\lambda ),`$
$`i=1,\mathrm{},r`$. Notice that the operators $`\stackrel{~}{G}_{t_{\omega _i^{}}}^V`$ preserve the weight decomposition of $`V`$. Notice also that the operators $`\stackrel{~}{G}_{t_{\omega _i^{}}}^V`$ are well defined on the tensor product of any highest weight $`𝔤`$-modules according to formula (14).
An easy corollary of the compatibility of system (15) is
###### Lemma 23.
The dynamical difference equations (16) form a compatible system of equations for a $`V`$-valued function $`u(z_1,\mathrm{},z_n,\lambda )`$.
In particular, for $`𝔤=sl_N`$, the Lemma says that the system (10) is compatible.
###### Theorem 24.
Assume that the Lie algebra $`𝔤`$ has a minuscle dual fundamental weight, i.e. $`𝔤`$ is not of type $`E_8,F_4,G_2`$. Then the dynamical equations (16) together with the KZ equations (1) form a compatible system of equations.
The Theorem is proved in Section 4.4.
We conjecture that the statement of the Theorem holds for any simple Lie algebra.
Let $`𝔤`$ be a simple Lie algebra for which the KZ and dynamical equations are compatible. Let $`V=V_1\mathrm{}V_n`$ be a tensor product of finite dimensional $`𝔤`$-modules. Fix a basis $`v_1,\mathrm{},v_d`$ in a weight subspace $`V[\nu ]`$. Suppose that $`u_i(z_1,\mathrm{},z_n,\lambda )=_{j=1}^du_{i,j}v_j`$, $`i=1,\mathrm{},d`$, is a set of $`V[\nu ]`$-valued solutions of the combined system of KZ equations (1) and dynamical equations (16).
###### Corollary 25.
$$\text{det}(u_{i,j})_{1i,jd}=C_{V[\nu ]}(\lambda )D_{V[\nu ]}(z_1,\mathrm{},z_n,\lambda )$$
where $`C_{V[\nu ]}(\lambda )`$ is a function of $`\lambda `$ (depending also on $`V_1,\mathrm{},V_n`$ and $`\nu `$) such that
$$C_{V[\nu ]}(\lambda +\kappa \omega )=C_{V[\nu ]}(\lambda )$$
for all $`\omega P^{}`$ and $`D_{V[\nu ]}(z_1,\mathrm{},z_n,\lambda )`$ is defined in (13).
The Corollary follows from formula (14).
### 4.4. Proof of Theorem 24
Introduce an action of $`𝕎^b`$ on the KZ operators $`_j(\lambda ,\kappa ),j=1,\mathrm{},n`$. Namely, for any $`w𝕎`$, set
$${}_{}{}^{w}_{j}^{}(\lambda ,\kappa )=w(_j(w^1(\lambda ),\kappa ))=\kappa z_j\frac{}{z_j}\underset{l,lj}{}w(r(z_j/z_l))^{(j,l)}\lambda ^{(j)}$$
and for any $`\omega P_t^{}`$ set
$`{}_{}{}^{t_\omega }_{j}^{}(\lambda ,\kappa )`$ $`=`$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{\omega ^{(k)}}_j(\lambda \kappa \omega ,\kappa ){\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{\omega ^{(k)}}=`$
$`\kappa z_j{\displaystyle \frac{}{z_j}}`$ $``$ $`{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{\omega _i^{(k)}}\left({\displaystyle \underset{l,lj}{}}r(z_j/z_l)^{(j,l)}\right){\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{\omega _i^{(k)}}\lambda ^{(j)}.`$
The compatibility conditions of the dynamical and KZ equations take the form
$$\stackrel{~}{G}_{t_{\omega _i^{}}}^V(z_1,\mathrm{},z_n,\lambda )_j(\lambda ,\kappa )={}_{}{}^{t_{\omega _i^{}}}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_{t_{\omega _i^{}}}^V(z_1,\mathrm{},z_n,\lambda )$$
for $`i=1,\mathrm{},r`$, $`j=1,\mathrm{},n`$.
The compatibility conditions follow from a more general statement.
###### Theorem 26.
Assume that the Lie algebra $`𝔤`$ has a minuscle dual fundamental weight, i.e. $`𝔤`$ is not of type $`E_8,F_4,G_2`$. Then for any $`j=1,\mathrm{},n`$ and any $`w𝕎^b`$ we have
$`\stackrel{~}{G}_w^V(z_1,\mathrm{},z_n,\lambda )_j(\lambda ,\kappa )={}_{}{}^{w^1}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_w^V(z_1,\mathrm{},z_n,\lambda ).`$
We conjecture that the statement of the Theorem holds for any simple Lie algebra.
The Theorem follows from the next four Lemmas.
###### Lemma 27.
Let $`j=1,\mathrm{},n`$. Assume that
$$\stackrel{~}{G}_{s_l}^V_j(\lambda ,\kappa )={}_{}{}^{s_l}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_{s_l}^V,{}_{}{}^{\pi _i}_{j}^{}(\lambda ,\kappa )=_j(\lambda ,\kappa )$$
for $`l=0,\mathrm{},r`$ and $`iO^{}`$. Then
$`\stackrel{~}{G}_w^V(z_1,\mathrm{},z_n,\lambda )_j(\lambda ,\kappa )={}_{}{}^{w^1}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_w^V(z_1,\mathrm{},z_n,\lambda )`$
for all $`w𝕎^b`$.
###### Proof.
If $`w=\pi _is_{m_l}\mathrm{}s_{m_1}`$ is a reduced presentation, then
$`\stackrel{~}{G}_w^V={}_{}{}^{s_{m_1}\mathrm{}s_{m_{l1}}}(\stackrel{~}{G}_{s_{m_l}}^V)\mathrm{}{}_{}{}^{s_{m_1}}(\stackrel{~}{G}_{s_{m_2}}^V)\stackrel{~}{G}_{s_{m_1}}^V`$ and
$`\stackrel{~}{G}_w^V_j(\lambda ,\kappa )={}_{}{}^{s_{m_1}\mathrm{}s_{m_{l1}}}(\stackrel{~}{G}_{s_{m_l}}^V)\mathrm{}{}_{}{}^{s_{m_1}}(\stackrel{~}{G}_{s_{m_2}}^V)\stackrel{~}{G}_{s_{m_1}}^V_j(\lambda ,\kappa )=`$
$`{}_{}{}^{s_{m_1}\mathrm{}s_{m_{l1}}}(\stackrel{~}{G}_{s_{m_l}}^V)\mathrm{}{}_{}{}^{s_{m_1}}(\stackrel{~}{G}_{s_{m_2}}^V){}_{}{}^{s_{m_1}}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_{s_{m_1}}^V=`$
$`{}_{}{}^{s_{m_1}\mathrm{}s_{m_{l1}}}(\stackrel{~}{G}_{s_{m_l}}^V)\mathrm{}{}_{}{}^{s_{m_1}s_{m_2}}_{j}^{}(\lambda ,\kappa ){}_{}{}^{s_{m_1}}(\stackrel{~}{G}_{s_{m_2}}^V)\stackrel{~}{G}_{s_{m_1}}^V=`$
$`{}_{}{}^{s_{m_1}s_{m_2}\mathrm{}s_{m_l}}_{j}^{}(\lambda ,\kappa ){}_{}{}^{s_{m_1}\mathrm{}s_{m_{l1}}}(\stackrel{~}{G}_{s_{m_l}}^V)\mathrm{}{}_{}{}^{s_{m_1}}(\stackrel{~}{G}_{s_{m_2}}^V)\stackrel{~}{G}_{s_{m_1}}^V=`$
$`{}_{}{}^{w^1}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_w^V.`$
###### Lemma 28.
Let $`j=1,\mathrm{},n`$ and $`w𝕎`$. Then
$$\stackrel{~}{G}_w^V_j(\lambda ,\kappa )={}_{}{}^{w^1}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_w^V.$$
###### Proof.
For $`w𝕎`$ we have $`\stackrel{~}{G}_w^V(z_1,\mathrm{},z_n\lambda )=𝔹_{w,V}(\lambda )`$, and Lemma 28 is equivalent to Lemma 18. ∎
###### Lemma 29.
Let $`j=1,\mathrm{},n`$ and $`iO^{}`$. Then
$${}_{}{}^{\pi _i}_{j}^{}(\lambda ,\kappa )=_j(\lambda ,\kappa ).$$
###### Proof.
We have $`\pi _i=t_{\omega _i^{}}w_{[i]}^1`$. Hence
$`{}_{}{}^{\pi _i}_{j}^{}(\lambda ,\kappa )={}_{}{}^{t_{\omega _i^{}}}({}_{}{}^{w_{[i]}^1}_{j}^{}(\lambda ,\kappa ))={}_{}{}^{t_{\omega _i^{}}}(\kappa z_j{\displaystyle \frac{}{z_j}}{\displaystyle \underset{l,lj}{}}w_{[i]}^1(r(z_j/z_l))^{(j,l)}\lambda ^{(j)})=`$
$`\kappa z_j{\displaystyle \frac{}{z_j}}{\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _i^{})^{(k)}}\left({\displaystyle \underset{l,lj}{}}w_{[i]}^1(r(z_j/z_l))^{(j,l)}\right){\displaystyle \underset{k=1}{\overset{n}{}}}z_k^{(\omega _i^{})^{(k)}}\lambda ^{(j)}=_j(\lambda ,\kappa ).`$
The last equality follows from Lemma 3. ∎
###### Lemma 30.
Let $`j=1,\mathrm{},n`$. Assume that the Lie algebra $`𝔤`$ has a minuscle dual fundamental weight. Then
$$\stackrel{~}{G}_{s_0}^V_j(\lambda ,\kappa )={}_{}{}^{s_0}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_{s_0}^V.$$
###### Proof.
Let $`\omega _i^{}`$ be a minuscle dual fundamental weight. We have $`s_0=\pi _i^1s_i\pi _i`$ and $`\stackrel{~}{G}_{s_0}^V={}_{}{}^{\pi _i^1}(\stackrel{~}{G}_{s_i}^V)`$ according to the 1-cocycle property. Now
$`{}_{}{}^{s_0}_{j}^{}(\lambda ,\kappa )\stackrel{~}{G}_{s_0}^V={}_{}{}^{\pi _i^1s_i\pi _i}_{j}^{}(\lambda ,\kappa ){}_{}{}^{\pi _i^1}(\stackrel{~}{G}_{s_i}^V)={}_{}{}^{\pi _i^1}({}_{}{}^{s_i}({}_{}{}^{\pi _i}_{j}^{}(\lambda ,\kappa ))\stackrel{~}{G}_{s_i}^V)=`$
$`{}_{}{}^{\pi _i^1}({}_{}{}^{s_i}(_j(\lambda ,\kappa ))\stackrel{~}{G}_{s_i}^V)={}_{}{}^{\pi _i^1}(\stackrel{~}{G}_{s_i}^V_j(\lambda ,\kappa ))={}_{}{}^{\pi _i^1}(\stackrel{~}{G}_{s_i}^V){}_{}{}^{\pi _i^1}(_j(\lambda ,\kappa ))=\stackrel{~}{G}_{s_0}^V_j(\lambda ,\kappa ).`$
Theorems 24 and 26 are proved. |
warning/0002/astro-ph0002298.html | ar5iv | text | # X–ray variability and prediction of TeV emission in the HBL 1ES 1101–232
## 1 Introduction
BL Lac objects form a minority class of active nuclei (see e.g. Urry & Padovani, 1995), but nevertheless their high luminosities and extreme variability in all bands make them an interesting subclass to study intrinsic properties of nuclear emission. Furthermore, it is matter of discussion the position of BL Lacs with respect to other active nuclei in the framework of Unification Models. The current picture claims that BL Lacs are the fraction of Fanaroff-Riley I galaxies that point their jet towards us, but many details still need to be worked out. Often, BL Lacs are studied together with other classes of flat radio spectrum sources, to form the class of blazars.
Recently a sequence has been proposed for the subclass of gamma ray bright blazars, and possibly valid for all blazars, based on their bolometric luminosity (Ghisellini et al. 1998, Fossati et al. 1998). The overall spectral energy distribution has two peaks (in $`\nu F_\nu `$ representation), the one at lower energies due to synchrotron radiation and the higher energy one to Inverse Compton scattering. It is proposed that the value of the peak frequency is the result of the balance between radiative cooling and acceleration of the corresponding electrons (see e.g. Ghisellini, 1999): in less luminous objects the radiative cooling is less efficient, allowing the accelerated electrons to reach higher energies. As a consequence, both the synchrotron and the inverse Compton peaks shift to higher frequencies as the bolometric intrinsic luminosity decreases. In this view, the class of BL Lacertae objects known as HBL (High frequency peak BL Lacs) has a smaller bolometric luminosity and a higher synchrotron peak frequency than the class of LBL (Low frequency peak BL Lacs). HBL are mostly found in X–ray selected samples, since they are expected to produce most of their synchrotron emission in the X–ray band. If the emission peak (as in the LBL objects) is at frequencies smaller than the observed X–ray band, the X–ray spectrum is steep (i.e. the steepening part above the peak), while HBL with a flat X–ray spectrum should have their peak in the observed X–ray band.
As part of a program aimed at a spectral survey of soft X–ray selected BL Lacs with BeppoSAX (Wolter et al. 1998), we have studied also 1ES 1101–232 (z=0.186), a bright BL Lac selected from the Slew Survey (Perlman et al. 1996) that shows an extreme behavior with a very flat X–ray spectrum. The object was detected, besides in the 0.1-10 keV energy band, also up to $``$ 100 keV in only $`6000`$ sec. The best fit of the source in the 0.1-100 keV band is given by a broken power law, with Galactic low energy absorption, that has a break energy E<sub>0</sub>=1.36 (1.11–1.65) keV. The spectral slope (energy index) derived from the PDS is consistent with the one derived from the MECS above the break energy E<sub>0</sub>: $`\alpha _x`$ = 1.03 (0.99–1.08) (Wolter et al. 1998).
We have constructed also the Spectral Energy Distribution (SED) for this object, using flux measurements collected from the literature, from radio to X–rays. We have fitted a cubic polynomial to the distribution in order to find the peak of the SED, that indeed falls in the BeppoSAX band (log $`\nu _{peak}`$ 17.48, corresponding to $`1.3`$ keV). For this object, therefore, the break energy derived from the spectral fit in the X–ray band is consistent, albeit within its large indetermination, with the position of the synchrotron peak as derived from the overall distribution (SED).
The SED of 1ES 1101–232 is similar to that of the flaring state of Mkn 501 (Pian et al. 1998) and 1ES 2344+514 (Catanese et al. 1998; Giommi et al. 2000), and therefore we could expect a strong TeV emission. A quasi–simultaneous X–ray and TeV observation has been therefore scheduled, to confirm the X–ray spectrum of the source, and its overall shape, to detect possible variations in flux, that could constrain the physical parameters of the source, and to monitor the TeV emission to search for possible detection.
The BeppoSAX data are presented here. The TeV observation, conducted in non–optimal weather condition, did not yield a detection, but we will use the upper limit (Chadwick et al. 1999a) to derive useful information on the physical mechanisms at work in this source. The plan of the paper is as follows: in Section 2 we describe the observational data obtained with BeppoSAX in the two epochs, in Section 3 we summarize the TeV predictions and observations, that help constraining the parameters of the source, by using the SED and theoretical models of emission as explained in Section 4. Section 5 presents our results and conclusions.
Throughout the paper a Hubble constant H<sub>0</sub>=50 $`kms^1Mpc^1`$ and a Friedman universe with a deceleration parameter $`q_0`$=0 are assumed.
## 2 BeppoSAX data
The X–ray astronomy satellite BeppoSAX is a project of the Italian Space Agency (ASI) with a participation of the Netherlands Agency for Aerospace Programs (NIVR). The scientific payload comprises four Narrow Field Instruments \[NFI: Low Energy Concentrator Spectrometer (LECS), Medium Energy Concentrator Spectrometer (MECS), High Pressure Gas Scintillation Proportional Counter (HPGSPC), and Phoswich Detector System (PDS)\], all pointing in the same direction, and two Wide Field Cameras (WFC), pointing in opposite directions perpendicular to the NFI common axis. A detailed description of the entire BeppoSAX mission can be found in Butler & Scarsi (1990) and Boella et al. (1997a).
The MECS consists of three equal units, each composed of a grazing incidence mirror unit and of a position sensitive gas scintillation proportional counter, with a field of view of 56 arcmin diameter, working range 1.3–10 keV, energy resolution $`8\%`$ and angular resolution $`0.7`$ arcmin (FWHM) at 6 keV. The effective area at 6 keV is 155 cm<sup>2</sup> (Boella et al. 1997b)
The LECS is a unit similar to the MECS, with a thinner window that grants a lower energy cut-off (sensitive in the energy range 0.1-10.0 keV) but also reduces the FOV to 37 arcmin diameter (Parmar et al. 1997). The LECS energy resolution is a factor $`2.4`$ better than that of the ROSAT PSPC ($`32\%`$ at 0.28 keV), while the effective area is smaller: 22 cm<sup>2</sup> at 0.28 and and 50 cm<sup>2</sup> at 6 keV.
The PDS is a system of four crystals, sensitive in the 13–200 keV band and mounted on a couple of rocking collimators, which points two units on the targets and two units $`3.5^{}`$ aside respectively, to monitor the background. The position of the collimators flips every 96 seconds. Thanks to the stability of the instrumental background, the PDS has shown an unprecedented sensitivity in its energy range, allowing 3$`\sigma `$ detection of $`\alpha 1`$ sources as faint as 10 m Crab with 10 ks of effective exposure time (Guainazzi & Matteuzzi, 1997).
The source is not detected by the HPGSPC, so we will not discuss this instrument.
### 2.1 Observation of June 1998
The object has been observed in AO2 on 19 June 1998 for a total of 8958 sec (LECS – 3167 net counts); 24895 sec (MECS(2+3) – 10612 net counts) and 10792 sec (PDS on source – 1320 net counts). The AO2 observation has been performed with only 2 MECS (MECS2 and MECS3) since MECS1 was no longer active. The source has been observed almost in the same period (May 1998) by the Mark6 telescope (working in the GeV-TeV range).
The extraction of the BeppoSAX data has been performed with FTOOLS v4.0 and the spectral analysis with XSPEC v9.0, using the most recently available matrices (September 1997 release). The data analysis has been performed on the same guidelines as outlined by the BeppoSAX Cookbook (`http://www.sdc.asi.it/software/cookbook/`) and described e.g. in Wolter et al. (1998). We summarize here that counts are extracted in a circular region of 8.5’/4’ (LECS/MECS) radius, and the background is taken from the blank sky images distributed by the BeppoSAX Data Center, in a region corresponding to the one used to extract source counts. Counts are binned so as to have at least 30 total counts in each bin to ensure applicability of the $`\chi ^2`$ statistics. Fit to LECS data are performed only up to 4 keV, as the response matrix of LECS is not well calibrated above this energy (see Orr et al. 1998). All confidence levels are computed using $`\mathrm{\Delta }\chi ^2`$ = 2.7 (corresponding to 90% for 1 interesting parameter), unless otherwise stated.
For the LECS+MECS combined data, a single power law is not acceptable ($`>3\sigma `$) with $`N_\mathrm{H}`$ fixed at the Galactic value. If $`N_\mathrm{H}`$ is left free, the fit with a single power law is acceptable only at about $`3\sigma `$ (Prob $`1.5\%`$). The residuals are however skewed, showing that the fit is not good. The fit is significantly improved (F–test at $`>`$ 99.99 % probability) by using a broken power law shape. Results for the three models for LECS and MECS data are listed in Table 1.
For each observation, in the first fit (single power-law model with Galactic $`N_\mathrm{H}`$) we left the LECS normalization free with respect to the MECS normalization to account for the residual errors in intensity cross-calibration (see Cusumano, Mineo, Guainazzi et al. in preparation). The fitted value of 0.696 (1997) and 0.701 (1998) fall in the range expected given the current knowledge of the cross-calibration (F. Fiore, private communication; see also `http://www.sdc.asi.it/software/cookbook/cross_cal.html`). Since the ratio of the two normalizations depends on the position of the source in the detector, and not on the model chosen, we fix the LECS/MECS normalization to 0.7 also for the other subsequent models.
The PDS exposure time is not even twice than in AO1 (10.8 ks vs. 6.4 ks). Since the spectrum is steeper and the source is fainter than in the AO1 observation, the PDS detection is not more significant than the AO1 detection ($`2.6\sigma `$). In order to fit the PDS data alone we rebin them to 6 data points. We fix the absorbing column to the Galactic value (N$`{}_{H}{}^{}=5.76\times 10^{20}`$ cm<sup>-2</sup>) while the normalization with respect to the MECS and index are left free. The best fit slope is $`\alpha _{PDS}`$ = 1.02 \[$`<`$ 2.65\], and the PDS/MECS ratio is 0.90, consistent with what expected on the basis of cross-calibration of the instruments, with a $`\chi ^2`$ = 2.2(4 dof), for a probability of 70%, therefore statistically acceptable. However, given the low statistical significance of the PDS data, the uncertainty on the slope is high. We therefore fit the total spectrum, from 0.1 keV to $``$50 keV, using LECS, MECS and PDS data together: the result of the broken power law fit is shown in Figure 1. It yields an unabsorbed flux in the \[2–10 keV\] band of $`F_X=2.54\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, and a corresponding luminosity in the same band of $`L_X=4.7\times 10^{45}`$ erg s<sup>-1</sup>. The PDS data are consistent with the LECS+MECS extrapolation.
### 2.2 Comparison of the 1997 and 1998 observations
For ease of comparison, we report in Table 2 the AO1 observation results of 1ES 1101–232, from Wolter et al. (1998). The LECS exposure was 5195 sec (2484 net counts), the MECS exposure was 13830 (9509 net counts) and the PDS on-source exposure was 6410 sec (1996 net counts). The best fit models of the LECS+MECS spectra for various spectral shapes are listed in Table 2, while in Figure 2 the LECS+MECS+PDS spectrum is plotted, with the best fit of model (3).
In order to compare the two observations, we first check directly the count rates in the two epochs. The best representation is the ratio of the two observed spectra, that does not depend on the choice of models and parameters. We therefore bin the two spectra, the relative background and response matrices in 32 channels, in order to avoid having empty bins after background subtraction. We then divide the two background-subtracted count rates and plot the results (after having flagged out the energy ranges that are not well calibrated in the matrices) in Fig. 3. The ratio is consistent with being flat up to $`2`$ keV, and steepening after it, showing that a change in the spectrum occurred above $`2`$ keV. The two instruments, LECS and MECS, give a consistent ratio in the overlapping energy range. A fit with a constant in the interval 0.5–10 keV is not statistically acceptable.
Also comparing the fit results for the two observations of Jan 1997 and Jun 1998, we see that the high energy slope ($`\alpha _2`$) is steeper in the second one. The flux, with all the three best fit models reported in the tables, is a factor of $``$ 32% lower in the second than in the first one. The low energy slope ($`\alpha _1`$) and the break energy ($`E_0`$) are instead consistent within the errors between the two observations.
We can make therefore the hypothesis that the different fluxes and spectra between the two epochs are explained by a change in the spectral slope above the peak of the synchrotron emission. Hence we fit all the data (1997+1998) simultaneously, keeping the low energy slope and break energy tied (equal one to each other, but free to vary) between the two observations. On the contrary, the high energy slopes in the two observations are left independent. The $`N_\mathrm{H}`$ is fixed to the galactic value, that fits well both observations. The PDS data are not used for the comparison, being of low statistical significance; the LECS and MECS data are re-binned to 100 total counts for each bin, to improve the significance of the individual data points, since no small range feature is present. The LECS/MECS normalization is again fixed to 0.7.
The resulting values of $`\alpha _1`$ and $`E_0`$ are consistent with those of both single observations: $`\alpha _1`$=0.72 \[0.44, 0.88\]; $`E_0=1.17[0.93,1.43]`$ keV. The high energy slope is $`\alpha _2^{1998}`$ = 1.26 \[1.19, 1.32\], vs. $`\alpha _2^{1997}`$ = 1.01 \[0.95, 1.07\], confirming that the slope indeed steepened significantly. The $`\chi ^2`$ of the fit is 184.4 with 178 dof, corresponding to a probability of 36%. The errors quoted here are 90% confidence for 4 parameters of interest ($`\mathrm{\Delta }\chi ^2=7.76`$). The combined spectra (LECS and MECS) for both observations with the best fit model are shown in Figure 4. We report in Table 3 the unabsorbed fluxes and rest-frame luminosities in the 0.1–2 and 2–10 keV bands for the two observations using the combined fit results.
The fluxes derived using the combined model are consistent with the fluxes derived from the two independent fits to the observations. We can therefore attribute the observed flux variation entirely to a steepening of the high energy ($`>`$2 keV) portion of the spectrum. There is no evidence, within the statistical uncertainties, of a change in $`\alpha _1`$ or $`E_0`$, although in other well known sources (Mkn501, Mkn421, PKS2155–304) an increase in flux seems to be linked to an increase in $`E_0`$ and/or the high energy slope.
## 3 The TeV band
Blazars as a class have been shown to emit a large fraction of their power at high energies, in the MeV–TeV band. The current models assume that the high energy emission is produced by Inverse Compton scattering (e.g. the SSC model or the EC model, see Ghisellini et al. (1998) for the relevance of the two mechanisms), and the location of the peak of the Compton component depends mainly on the lower energy peak due to the Synchrotron component.
In particular, objects of the LBL kind, that have the synchrotron peak at soft energies, show their second peak at energies in the MeV–GeV range, and are in fact detected by EGRET on board CGRO (see Mukherjee et al. 1997). Objects of the HBL kind, that have their synchrotron peak at UV–X–ray frequencies, instead, should show the Compton peak in an even higher energy band. In fact, up to now, the only sources of VHE gamma rays (in the TeV band, by using $`\stackrel{ˇ}{\mathrm{C}}`$erenkov detectors) are HBL: Mkn 421 (Punch et al. 1992), Mkn 501 (Quinn et al. 1996), 1ES 2344+514 (Catanese et al. 1998), PKS 2155–304 (Chadwick et al. 1999b).
The new $`\stackrel{ˇ}{\mathrm{C}}`$erenkov arrays allow the detection of bright sources in relatively short exposure times, and therefore it has been possible to monitor their variability: some of these objects have in fact shown periods of flaring activity on time-scales as short as 15 minutes (Gaidos et al. 1996, Aharonian et al. 1999).
Another point of debate is the amount of absorption of VHE photons due to the cosmic infra–red background. The aforementioned detected objects in fact are all nearby. Detection of sources that lie further away can therefore help in constraining the amount of the IR background (e.g. Stecker and De Jager, 1998).
The simultaneous detection of X–ray and TeV emission allows us to estimate a number of physical parameters of the source, such as the magnetic field and the Doppler factor (see e.g. Tavecchio et al. 1998). We will show in Section 4 that a number of interesting constraints can be derived also by using the upper limit derived in the TeV band, together with the X–ray band information.
### 3.1 Predictions of TeV emission in 1ES 1101–232 from phenomenological constraints.
A very simple prediction of the TeV emission can be made by following the scheme presented e.g. in Fossati et al. (1998). In this model a) the ratio of the frequencies of the high (Compton) and low (synchrotron) energy peaks is constant and equal to $`5\times 10^8`$ and b) the high energy peak and the radio luminosity have a fixed ratio, $`\frac{\nu _\gamma L_{peak,\gamma }}{\nu _{5GHz}L_{5GHz}}`$ =$`3\times 10^3`$. This model represents the average SED observed for the various classes, while single sources can deviate from this average phenomenological parameterization. This very simple relationship, however, allows us to make order of magnitude prediction even not knowing the physical conditions at the source. From the two expressions above we can derive both the expected frequency of the second peak and its intensity. The peak of the synchrotron emission is measured (this paper and Wolter et al. 1998) at $`\nu _s=3\times 10^{17}`$ Hz: the Compton peak is therefore expected at $`\nu _C=1.5\times 10^{26}`$ Hz which corresponds to $``$0.6 TeV. The measured radio flux at 5 GHz is $`5\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> (see references in Table 4). The expected flux at the Compton peak (0.6 TeV) is therefore $`8.\times 10^{12}`$ photon cm<sup>-2</sup> s<sup>-1</sup>.
Another prediction can be made from the observed X–ray flux, by using the recipe of Stecker, De Jager & Salamon (1996). Within an SSC scenario they use simple scaling arguments to predict the TeV fluxes for HBL, based on the X–ray flux and the assumption that the properties of the emission are similar to those observed for Mkn 421. Their argument is partially supported by the actual detection of (a few) other sources for which they predicted possible detectability. The factor of increase between the synchrotron and Compton component is $`10^9`$, and, assuming $`L_C/L_s1`$, they derive $`\nu _{TeV}F_{TeV}\nu _XF_X`$. Therefore, from the observed X–ray properties for 1ES 1101–232 (peak around 1 keV and flux of 2.5–3.8 $`\times 10^{11}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$), we infer that the expected Compton peak is around 1 TeV with a flux of $`1.5`$$`2.5\times 10^{11}`$ photon cm<sup>-2</sup> s<sup>-1</sup>.
These estimates make 1ES 1101–232 a good TeV candidate, at the border of current sensitivities, and possibly well above it during flaring activities.
### 3.2 TeV Observations
1ES 1101–232 has been observed on the nights of 19–27 May 1998 with the Durham University Mark 6 atmospheric $`\stackrel{ˇ}{\mathrm{C}}`$erenkov telescope. We summarize here the data analysis and the results presented in Chadwick et al. (1999a). The telescope uses the imaging technique to separate VHE gamma rays from the cosmic ray background, together with a robust noise-free trigger (Armstrong et al. 1999). Data are taken in 15 minutes segments, alternating ON-source with an equal number of OFF-source observations. After removal of cloud–affected data, there are a total of 10.5 hours ON-source data and the same amount of OFF-source data. Data are screened for “good” events by the selection criteria listed in Table 2 of Chadwick et al. (1999a). The source was not detected and an upper limit of $`F_{TeV}`$ \[$`>300`$ GeV\] = 3.7$`\times 10^{11}`$ photons cm<sup>-2</sup> s<sup>-1</sup> has been derived.
The data have been investigated for time variability on time-scales of days (at a flux limit of $`1\times 10^{10}`$ cm<sup>-2</sup> s<sup>-1</sup>) and 15 min intervals. There is no evidence for any bursting behavior (Chadwick, private communication).
## 4 SED construction and TeV predictions from theoretical models.
We have constructed the overall Spectral Energy Distribution (SED), using both data from literature and the BeppoSAX spectra from the two observations. We construct two different SED, one for each BeppoSAX observation, that are clearly modeled by a different synchrotron state. The two SED are presented in Fig. 5.
We have reproduced the SED in both X–ray states with a simple homogeneous SSC model. The source is modeled as a spherical region with size $`R`$, uniform and tangled magnetic field $`B`$, in motion toward the observer with a bulk Lorentz factor $`\mathrm{\Gamma }`$. The region is filled by a population of relativistic electrons with a distribution of Lorentz factors given by: $`N(\gamma )=K\gamma ^{n_1}(1+\gamma /\gamma _b)^{n_1n_2}`$, where the asymptotic slopes are $`n_1`$ and $`n_2`$, the break point is $`\gamma _b`$ and $`K`$ is a normalization factor. The self–Compton emission is derived taking into account the full Klein–Nishina (KN) cross section, computed using the relations reported in Jones (1968; see also Blumenthal & Gould 1970). We do not take into account absorption of IR photons by the infrared background, whose emission level is still uncertain. This implies that in principle the derived curves are upper limits to the detectable VHE emission, also because the redshift of 1ES 1101–232 is only slightly smaller than the limit ($`z=0.2`$) chosen to monitor blazar emission in the TeV band by the HEGRA experiment (Rhode & Meyer, 1997). On the other hand, a detection of the source in this band would provide also a measure of the density of IR background photons.
Although the complete determination of the set of physical parameters for the SSC model requires the knowledge of the positions of both the synchrotron peak and the Inverse Compton peak (as discussed in Tavecchio et al. 1998), we can put strong constraints to the parameters using the informations provided by the X–ray spectrum and the TeV upper limit suggesting that the condition $`L_C/L_s1`$ applies.
We note here that the analytical relations discussed in Tavecchio et al. (1998) in the KN regime are obtained with a step approximation for the KN cross–section. In the extreme KN regime the numerical values given by this approximation might differ from the results of the numerical model derived with the full KN treatment, but we can use the analytical discussion as a guideline for the numerical model.
Typical variability time-scales observed in HBLs ($`t_{var}\mathrm{\hspace{0.17em}10}^{34}`$ s, see e.g Zhang et al. 2000, Giommi et al. 1999, 2000) suggest Doppler factors $`\delta `$ in the range 10–20 and sizes of $`R10^{16}`$ cm. We fixed the radius to $`R=10^{16}`$ cm. This choice directly puts a lower limit to the value of the magnetic field:
$$B\delta ^{2+\alpha _1}>A\times \left(1+z\right)^{\alpha _1}\left[\frac{g}{\nu _c\nu _s}\right]^{\left(1\alpha _1\right)/2}\left(\frac{L_s}{t_{var}L_C^{1/2}}\right)$$
(1)
where $`g(\alpha _1,\alpha _2)`$ is a constant related to the spectral indices $`\alpha _1`$ and $`\alpha _2`$ and A is the appropriate constant (see Tavecchio et al. 1998 for details on (1)). From this relation we can infer that an upper limit on $`L_C`$ and an estimate of $`t_{var}`$ directly puts a lower limit for $`B`$. The physical reason for this lower limit is related to the fact that the synchrotron peak Luminosity is proportional to the product $`N_eB^2`$ (where $`N_e`$ is the number of emitting electrons), while the ratio between the Compton peak Luminosity and the synchrotron peak Luminosity is proportional to $`N_e`$: therefore an upper limit to $`L_C/L_s`$ gives an upper limit to $`N_e`$ and this, together with the synchrotron Luminosity, provides the lower limit for the magnetic field $`B`$. Another way of describing the effect, assuming that the scattering is in the Thomson regime, is that in this case requiring $`L_C/L_s<1`$ corresponds to require that the ratio between the synchrotron and the magnetic energy densities $`U_{syn}/U_B<1`$. For a given source size and Doppler factor the synchrotron radiation energy density is fixed, and the above relation then corresponds to a lower limit on the value of the magnetic field.
In Fig. 5 we plot the spectrum of both states computed for two different values of $`\delta `$ ($`\delta =10`$ for the dashed lines, $`\delta =20`$ for the solid lines), and for two different values of the magnetic field $`B`$, as listed in Table 5. The lowest value of $`B`$ has been determined by requiring not to over-produce the high energy (TeV) emission, since, for a given synchrotron luminosity, size and Doppler factor, the ratio between the self Compton and the synchrotron powers depends only on the magnetic field.
We take this value, listed in Table 5, for the models 1 and 3 shown in Fig. 5. For models 2 and 4 we have doubled the B value, and decreased the relativistic electron density and $`\gamma _b`$ accordingly, in order to produce the same amount of synchrotron flux and about the same synchrotron peak frequency. In this case the self–Compton flux decreases, due to the decreased electron density.
The transition from the high to the low state is consistent with a change of the second slope $`n_2`$ and with a decrease of $`\gamma _b`$ by a factor of 1.5–2. This behavior is similar to what observed in the other well known TeV BL Lac, such as PKS 2155–304 (see e.g. Chiappetti et al. 1999), Mkn 501 (e.g. Pian et al 1998) and Mkn 421 (Maraschi et al. 1999), where high X-ray states are interpreted as states with either higher $`\gamma _b`$ and/or higher magnetic field.
It is evident from Fig. 5 that a small change in the magnetic field, while still consistent with the X–ray (BeppoSAX) observations, produces a dramatically different amount of TeV photons. Assuming that the size of the source and the Doppler factor do not vary substantially with time, variations of the synchrotron flux can be attributed to changes of the density of electrons and/or the magnetic field. If this is the case, we expect that the TeV emission can be easily detected, either for X–ray fluxes slightly brighter than what observed up to now, or by longer TeV exposure times.
## 5 Results and Conclusions
The X–ray spectrum of 1ES 1101–232 as measured by BeppoSAX is fitted only by a broken power law (a single power law or an absorbed power law are not statistically acceptable) with a break at 1.3 - 1.9 keV. From the first to the second observation, the spectrum varied at high energies, becoming softer (steeper). The flux decrease, by about 32%, has occurred in the 2–10 keV band. The PDS observation are not of statistical significance sufficient to put a real constraint on the spectrum.
Even if the variation in the X–ray band is not dramatic, we can clearly distinguish between the two states, that are modeled by different parameters of the synchrotron component.
By using the TeV upper limit and the two BeppoSAX observations we model also the higher energy portion of the spectrum as a self–Compton component, by using the model described e.g. in Tavecchio et al. (1998) that assumes a simple homogeneous SSC model in the KN regime, in which the relativistic electrons have a broken power law energy distribution.
The two X-ray states of the source are described by varying this distribution, assuming that the other relevant parameters ($`R`$ and $`\delta `$) are nearly constant.
We can compare these results with what found for the few TeV detected sources. The choice of Doppler factor of 10 and 20 made here is in the interval of the values of $`\delta `$ found by other authors for Mkn 421, PKS2155–304 and Mkn501, that range from $`\delta 5`$ (e.g. Takahashi, 1999; Mkn 421, Catanese et al. 1998: Mkn 501) to $`\delta 30`$ (e.g. Bednarek & Protheroe 1999: Mkn501; Kataoka et al 2000: PKS2155–304; Bednarek & Protheroe 1997: Mkn 421). At the same time, values of $`B`$ in excess of 0.03 G and up to 1 G (Chiappetti et al. 1999) are found for the same sources.
Of the three above mentioned objects, the most similar to 1ES 1101–232 is Mkn 501, for which different measures have been produced: e.g. $`\delta 15`$ & $`B=0.8`$G (Pian et al. 1998); $`\delta =15`$ & $`B=0.2`$G (Kataoka et al. 1999); $`\delta =30`$ & $`B=0.7`$G (Bednarek & Protheroe 1999), in good agreement with the values of Table 5.
Even if based, besides the accurate X-ray spectral determination up to $`50`$ keV, only on a TeV upper limit we can infer that the physical conditions in 1ES 1101–232 are similar to the brightest TeV sources, making it a very promising candidate for TeV observations, and a testbed for the SSC model. Furthermore, since the redshift of 1ES 1101–232 is intermediate, a detection of this source in the VHE range would pose constraints on the density of the IR background photons that is still at the moment very uncertain.
###### Acknowledgements.
This work has received partial financial support from the Italian Space Agency. We would like to thank Paula Chadwick and S.J. McQueen for informing us about their VHE results in advance of publication, Paolo Giommi and Roberto Della Ceca for helpful discussion and comments, and an anonymous referee for useful suggestions that improved the readability of the paper. |
warning/0002/gr-qc0002026.html | ar5iv | text | # Some integrals ocurring in a topology change problem
## I Introduction
A few years ago De Lorenci et al. presented a model of quantum cosmology which allowed space topology changes, having as main idea the use of the “conditional probability interpretation” to establish selection rules for the possible changes of topology; the wavefunctions involved in the process were of the type
$$\mathrm{\Psi }=\mathrm{\Psi }(\alpha ,\beta ,\xi ,\varphi )=A_k(a,\varphi )e^{\xi F_k},$$
(1)
where $`\alpha `$ and $`\beta `$ are appropriated canonical variables built upon the more common set of spherical coordinates $`(\chi ,\theta ,\phi )`$, the scale factor $`a`$ and the curvature $`k`$; $`\xi `$ and $`\varphi `$ are, respectively, a dust field describing a “distribution of irrotational dust particles” and a scalar field, both representing the matter content of the model; and $`F_k`$ is basically a numerical coeficient obtained by integration of certain functions constructed upon the ‘value’ $`\chi _0(\theta ,\phi ;V^3)`$ of the radial coordinate of the fundamental polyhedron’s boundary of the 3-dimensional manifold $`V^3`$, of curvature $`k`$, considered, written explicitly as
$$F_k=\frac{a}{2\pi \mathrm{}m}_{V^3}\frac{\mathrm{sin}2\sqrt{k}\chi _0(\theta ,\phi ;V^3)}{2\sqrt{k}}𝑑\theta \mathrm{sin}\theta d\phi .$$
(2)
The topology changes would occur at some value $`\xi `$ of the the dust field, when $`a=\overline{a}`$ and $`\varphi =\overline{\varphi }`$, such that the conditional probability of having $`k=1`$, $`0`$ or $`+1`$ would be
$`P_c\left(k|\overline{a},\overline{\varphi }\right)`$ $`=`$ $`{\displaystyle \frac{\left|\mathrm{\Psi }(k,\overline{a},\overline{\varphi })\right|^2}{_{k=0,\pm 1}\left|\mathrm{\Psi }(k,\overline{a},\overline{\varphi })\right|^2}}`$ (3)
$`=`$ $`{\displaystyle \frac{A_k^2(\overline{a},\overline{\varphi })e^{2\xi F_k}}{_{k^{}=0,\pm 1}A_k^{}^2(\overline{a},\overline{\varphi })e^{2\xi F_k^{}}}}.`$ (4)
So, when $`\xi \pm \mathrm{}`$ one has one of the $`P_c\left(k|\overline{a},\overline{\varphi }\right)`$ equal to one and the other two null, depending upon the value of $`F_k`$.
In the values of the functions $`F_k`$ were only estimated for two different compact manifolds, the Poincaré dodecahedral space $`D^3`$, of positive curvature, and the hyperbolic icosaedral space $`I^3`$ (also known as Best space), of negative curvature; since there the authors claimed that “it is not possible to calculate the $`F_i`$’s exactly” for these manifolds, the importance of the present work is in the exact calculation of the functions $`F_k`$ for several compact manifolds of constant negative curvature, including the cited $`I^3`$.
## II Some calculus in compact manifolds
The functions $`F_k`$, such as presented in equation (2), are probably uncomputable since the specific form of the functions $`\chi _0(\theta ,\phi ;V^3)`$ are difficult, if not impossible, to determine; however, one can simply establish the following limits for the $`F_k`$’s:
$$4\pi \frac{\mathrm{sin}2\sqrt{k}\chi _{\mathrm{min}}}{2\sqrt{k}}\frac{F_k}{\left(a/2\pi \mathrm{}m\right)}4\pi \frac{\mathrm{sin}2\sqrt{k}\chi _{\mathrm{max}}}{2\sqrt{k}},$$
(5)
where $`\chi _{\mathrm{min}}`$ and $`\chi _{\mathrm{max}}`$ are, respectively, the radii of the inscribed and circumscribed circunference of the the fundamental cell of the manifold in consideration. In the functions $`\chi _0`$ appear after performing an “integration with respect to the variable $`\chi `$”, using as interval of integration $`[0,\chi _0(\theta ,\phi ;V^3)]`$; so, in order to obtain a numerical value for the functions $`F_k`$, it is easy to see that one can start with the integral
$$F_k=\frac{a}{2\pi \mathrm{}m}_{V^3}\left[\mathrm{sin}^2\sqrt{k}\chi \mathrm{cos}^2\sqrt{k}\chi \right]𝑑\chi 𝑑\theta \mathrm{sin}\theta d\phi .$$
(6)
Noticing now that
$$dV=\frac{\mathrm{sin}^2\sqrt{k}\chi }{k}d\chi d\theta \mathrm{sin}\theta d\phi $$
(7)
is simply the element of volume for the spatial part of a Friedmann-Robertson-Walker metric, written in spherical coordinates, there are two possible ways to follow, one plainer and the other a little more sophisticated; in both, however, one needs to redefine the coordinates and limits of integration used. So, the next step consists in the use of cylindrical coordinates $`(\rho ,\phi ,z)`$, related to the spherical coordinates $`(\chi ,\theta ,\phi )`$ by means of the relations
$$\{\begin{array}{c}\mathrm{cos}\sqrt{k}\chi =\mathrm{cos}\sqrt{k}\rho \mathrm{cos}\sqrt{k}z\\ \mathrm{sin}\sqrt{k}\chi \mathrm{sin}\theta =\mathrm{sin}\sqrt{k}\rho \end{array}$$
(8)
or
$`d\chi ^2`$ $`+`$ $`{\displaystyle \frac{\mathrm{sin}^2\sqrt{k}\chi }{k}}\left[d\theta ^2+\mathrm{sin}^2\theta d\phi ^2\right]=`$ (10)
$`d\rho ^2+\mathrm{cos}^2\sqrt{k}\rho dz^2+{\displaystyle \frac{\mathrm{sin}^2\sqrt{k}\rho }{k}}d\phi ^2.`$
Now, one has the interval of integration $`[0,\rho _0(z,\phi ;V^3)]`$ for the coordinate $`\rho `$; the expression for $`\rho _0(z,\phi ;V^3)`$ is easily obtainable, since is only a matter of using trigonometrical identities in the plane, i.e., in the triangles that compose the faces of each tetrahedron in which the fundamental polyhedron can be divided<sup>*</sup><sup>*</sup>*For more information on trigonometric identities in non-euclidean spaces one can see references to ; is a classical book of cosmology with one section on spherical trigonometry. using the following procedure:
* for each face draw a geodesic line perpendicular to it, connecting it to the center $`A`$ of the polyhedron, and crossing it or its plane in a point $`B`$ (this line $`AB`$ gives the height $`z`$ of the tetraedron);
* for each edge draw a geodesic line perpendicular to it and connecting it to the point $`B`$ of the face to which the edge belongs, crossing the edge or its extension in a point $`C`$;
* complete the tetrahedron with one of the two vertices of the edge, naming it as $`D`$.
These steps will create some ‘negative’ tetrahedra, covering also regions outside the polyhedron, and some ‘positive’, covering only regions of the polyhedron, each one of them having four right-angled triangles, one of which (named here $`BCD`$) is the base of the tetrahedron; integration on the compact manifold represented by the polyhedron is the difference between the sums of the integrations on all of the positive tetrahedra and the integrations on all of the negative tetrahedra.
The easiest path of integration consists of simply making
$`F_k`$ $`=`$ $`{\displaystyle \frac{a}{2\pi \mathrm{}m}}{\displaystyle _{V^3}}\left[2\mathrm{sin}^2\sqrt{k}\chi 1\right]𝑑\chi 𝑑\theta \mathrm{sin}\theta d\phi `$ (11)
$`=`$ $`{\displaystyle \frac{a}{2\pi \mathrm{}m}}\left[2kv{\displaystyle _{V^3}}{\displaystyle \frac{kdV}{\mathrm{sin}^2\sqrt{k}\chi }}\right],`$ (12)
where $`v`$ is the volume of the compact manifold where the integration is being performed. The remaining integral in the right hand side of the last equality must be done in the new set of cylindrical coordinates, where the limits of integration for the particular case of negative curvature ($`k=1`$) are, in each tetrahedronThe trigonometric identities that lead to such result are showed in an appendix at the end of this work.,
$$0zd_{AB},\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}\phi \stackrel{}{CBD},$$
(13)
and
$$0\rho \rho _0(z,\phi )=\mathrm{arctanh}\left[\mathrm{tan}\stackrel{}{BAC}\frac{\mathrm{sinh}z}{\mathrm{cos}\phi }\right].$$
(14)
The integration in the coordinate $`\rho `$ is easily done and gives finally
$`{\displaystyle \frac{F_1}{a/2\pi \mathrm{}m}}=2v+`$ (15)
$`{\displaystyle _0^{\stackrel{}{CBD}}}𝑑\phi {\displaystyle _0^{d_{AB}}}{\displaystyle \frac{dz}{\mathrm{cosh}^2z}}\mathrm{ln}\sqrt{{\displaystyle \frac{\mathrm{cos}^2\phi +\mathrm{tan}^2\stackrel{}{BAC}}{\mathrm{cos}^2\phi \mathrm{tan}^2\stackrel{}{BAC}\mathrm{sinh}^2z}}}`$ (16)
from where numerical results can be obtained by plain numerical integration. Notice that the same procedure can yield a formula for the volume of the manifold.
Alternatively, one can start doing
$`F_k`$ $`=`$ $`{\displaystyle \frac{ak}{2\pi \mathrm{}m}}{\displaystyle _{V^3}}\left[1\mathrm{cot}^2\sqrt{k}\chi \right]𝑑V`$ (17)
$`=`$ $`{\displaystyle \frac{ak}{2\pi \mathrm{}m}}\left[v+{\displaystyle _{V^3}}_\mu V^\mu dV\right]`$ (18)
where $`V^\mu `$ is a vector satisfying the differential equation
$$_\mu V^\mu =\left(_\mu +\mathrm{\Gamma }_{\mu \rho }^\rho \right)V^\mu =k\mathrm{cot}^2\sqrt{k}\chi ,$$
(19)
whose solution in spherical coordinates isHere is used the identity
$$\frac{\pi ^2}{4m^2}\mathrm{csc}^2\frac{\pi }{m}+\frac{\pi }{4m}\mathrm{cot}\frac{\pi }{m}\frac{1}{2}=\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{\left(1k^2m^2\right)^2}$$
found as equation $`1.423`$ of reference .
$`V^\chi `$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\sqrt{k}\mathrm{cot}\sqrt{k}\chi +k\chi \mathrm{csc}^2\sqrt{k}\chi \right]`$ (20)
$`=`$ $`{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{\mathrm{}=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{k^2\chi ^3}{\left(k\chi ^2\pi ^2\mathrm{}^2\right)^2}},`$ (21)
where the last equality was put to show clearly the behavior of the solution when $`k=0`$. This result permits the use of Stokes’s theorem to make
$$_{V^3}k\mathrm{cot}^2\sqrt{k}\chi dV=_{S=V^3}g_{\mu \nu }V^\mu n^\nu 𝑑A,$$
(22)
where $`n^\nu `$ is a vector normal to the boundary $`S`$ of the fundamental cell of the compact manifold $`V^3`$, obbeying the constraint $`n^\mu n_\mu =1`$.
In the procedure presented here the faces of the fundamental polyhedron that represents a compact manifold appear, by construction, as surfaces of constant $`z`$, allowing to use as element of area
$$dA=\frac{\mathrm{sin}\sqrt{k}\rho }{\sqrt{k}}d\rho d\phi .$$
(23)
Finally, to carry out the integration the vector $`V^\chi `$ must be written in cylindrical coordinates; only the component $`V^z=V^\chi _\chi z`$, normal to the base of the tetrahedron, is important. This procedure can be used also to give the volume of each tetrahedron, what allows to write, in the case of negative curvature,
$`{\displaystyle \frac{F_1}{a/2\pi \mathrm{}m}}=`$ (24)
$`{\displaystyle _0^{\stackrel{}{CBD}}}{\displaystyle \frac{d\phi }{\mathrm{coth}z}}\mathrm{ln}\sqrt{{\displaystyle \frac{\mathrm{cos}^2\phi +\mathrm{tan}^2\stackrel{}{BAC}}{\mathrm{cos}^2\phi \mathrm{tan}^2\stackrel{}{BAC}\mathrm{sinh}^2z}}},`$ (25)
where $`z=d_{AB}`$.
## III Numerical results
To obtain numerical results the data – volumes and coordinates of all vertices for several hyperbolic compact manifolds – contained in the literature were used (see, for instance, e ) together with those of the software SnapPea<sup>§</sup><sup>§</sup>§SnapPea is an electronic catalog of thousands of hyperbolic compact manifolds, each one of them identified by volume and a code such as $`m036(3,2)`$. ; part of the data used are presented in Table I. The manifolds choosed present in some way a degree of symmetry which simplified the calculus, but, in principle, the approach followed can be used to any compact manifold. All results are presented in the Table II where they are compared with estimates done as in ; the result obtained for the Weeks manifold was used in . TABLE I.: Data for each manifold studied. Manifold volume $`\chi _{min}`$ $`\chi _{max}`$ Weeks $`0.942707`$ $`0.519162`$ $`0.752470`$ Thurston $`0.981369`$ $`0.535437`$ $`0.748538`$ $`m036(3,2)`$ $`2.029883`$ $`0.675646`$ $`1.014814`$ $`m016(4,3)`$ $`2.343017`$ $`0.691286`$ $`0.895576`$ $`m036(2,3)`$ $`2.568971`$ $`0.726205`$ $`0.895576`$ Best $`4.686034`$ $`0.868298`$ $`1.382571`$ $`v3469(+3,1)`$ $`5.137941`$ $`0.808931`$ $`1.45241`$ TABLE II.: Summary of the results. Manifold $`F_1/\left(a/2\pi \mathrm{}m\right)`$ $`2\pi \mathrm{sinh}2\chi _{min}`$ $`2\pi \mathrm{sinh}2\chi _{max}`$ Weeks $`9.28474`$ $`7.76109`$ $`13.4518`$ Thurston $`9.48385`$ $`8.09029`$ $`13.3355`$ $`m036(3,2)`$ $`13.4897`$ $`11.3208`$ $`23.4987`$ $`m016(4,3)`$ $`14.5526`$ $`11.7314`$ $`18.3142`$ $`m036(2,3)`$ $`15.3167`$ $`12.6901`$ $`20.6181`$ Best $`21.4948`$ $`17.2847`$ $`49.6976`$ $`v3469(+3,1)`$ $`22.5418`$ $`15.2178`$ $`57.1996`$
## IV Conclusion
There are several formulations of quantum cosmology and the intention of this work is to put some new light over a particular one, showing that the wavefunctions built by the procedure of present a dependence on the volume of the compact manifold in consideration; aside that, such wavefunctions have an additional dependence on the shape of the fundamental cell of the manifold, due to a surface term that does not appear in several other models.
To finish, it is also interesting to notice that the results presented here, though of specific relevance for a particular model of quantum cosmology, can be seen in a more generalized context, since this work presents a method that allows to easily calculate the volume of the fundamental polyhedron of a compact manifold. Explicitly, for the particular case of negative curvature, the volume of each tetrahedron in which the fundamental polyhedron can be divided is
$`v={\displaystyle _0^{\stackrel{}{CBD}}}{\displaystyle \frac{d\phi }{2}}\times `$ (26)
$`\left\{{\displaystyle \frac{\mathrm{arctanh}\left[\mathrm{tanh}z\mathrm{sec}\phi \sqrt{\mathrm{cos}^2\phi +\mathrm{tan}^2\stackrel{}{BAC}}\right]}{\mathrm{sec}\phi \sqrt{\mathrm{cos}^2\phi +\mathrm{tan}^2\stackrel{}{BAC}}}}z\right\}`$ (27)
where again $`z=d_{AB}`$; alternatively,
$`v={\displaystyle _0^{d_{AB}}}{\displaystyle \frac{dz}{2}}\times `$ (28)
$`\left\{{\displaystyle \frac{arctanh\left[\mathrm{tan}\stackrel{}{CBD}\left(\mathrm{cot}^2\stackrel{}{BAC}csch^2z1\right)^{1/2}\right]}{\sqrt{\mathrm{cot}^2\stackrel{}{BAC}csch^2z1}}}\right\}.`$ (29)
These results must be compared with the more traditional ones given in , and .
## Acknowledgments
The author wants to thank the anonymous people of the Brazilian state of São Paulo, who gave him financial support through grant n. 96/0052-3 of the Fundação de Amparo à Pesquisa do Estado de São Paulo (FAPESP); the author is also deeply grateful for the help of R.G. Teixeira and Prof. Helio V. Fagundes.
## A Trigonometric identities
In the non-euclidean geometry the trigonometric identities valid for a triangle $`XYZ`$, with right angle $`Z`$, of sides $`x`$, $`y`$ and hypothenuse $`z`$ are
$$\mathrm{sin}Y=\frac{\mathrm{sin}\sqrt{k}y}{\mathrm{sin}\sqrt{k}z};\mathrm{cos}Y=\frac{\mathrm{tan}\sqrt{k}x}{\mathrm{tan}\sqrt{k}z};\mathrm{tan}Y=\frac{\mathrm{tan}\sqrt{k}y}{\mathrm{sin}\sqrt{k}x}.$$
(A1)
Using the second identity in the right-angled triangle $`BCD`$, of right angle $`C`$, and the third one in the right-angled triangle $`ABC`$, of right angle $`B`$, one can write, for the tetrahedron $`ABCD`$ built as in the section 2,
$$\mathrm{cos}\stackrel{}{CBD}=\mathrm{cos}\phi =\frac{\mathrm{tan}\sqrt{k}d_{BC}}{\mathrm{tan}\sqrt{k}\rho }=\frac{\mathrm{tan}\stackrel{}{BAC}\mathrm{sin}\sqrt{k}d_{AB}}{\mathrm{tan}\sqrt{k}\rho }$$
(A2)
from where one obtains equation (14), after identification of $`d_{AB}`$ with $`z`$. |
warning/0002/astro-ph0002364.html | ar5iv | text | # Nonlinear Velocity-Density Coupling: Analysis by Second-Order Perturbation Theory
## 1 Introduction
The peculiar velocity field is one of the most fundamental quantities to analyze the large-scale structure in the universe (e.g. Peebles 1980). It is considered to reflect dynamical nature of density fluctuations of gravitational matter. The peculiar velocity field is usually observed using astrophysical objects (e.g. galaxies), as determination of distances is crucial for measuring peculiar velocities (Dekel 1994, Strauss & Willick 1995). There is a possibility that statistical aspects of the velocity field traced by these objects and that traced by dark matter particles might be different. This difference is generally called “velocity bias” and its elucidation becomes highly important in observational cosmology (e.g. Cen & Ostriker 1992, Narayanan, Berlin & Weinberg 1998, Kaufmann et al. 1999).
Velocity bias is often discussed numerically with making “galaxy particles” in some effective manners. But here we discussed a more basic phenomenon. It is known that statistics of the peculiar velocity field depend largely on local density contrast. For example, both the single particle and pairwise velocity dispersions of dark matter particles are known to be increasing function of local density (Kepner, Summers, & Strauss 1997, Strauss, Cen & Ostriker 1998, Narayanan et al. 1998). Analysis of pairwise velocity statistics is interesting from theoretical point of views, and also very important in observational cosmology (Peebles 1976, Davis & Peebles 1983, Zurek et al. 1994, Fisher et al. 1994, Sheth 1996, Diaferio & Geller 1996, Suto & Jing 1997, Seto & Yokoyama 1998a, 1999, Jing & B$`\ddot{\mathrm{o}}`$rner 1998, Juszkiewicz, Fisher & Szapudi 1998, Seto 1999b, Juszkiewicz, Springel & Durrer 1999). But we do not discuss it here and concentrate on velocity field characterized by single point which is simpler to analyze theoretically.
Linear perturbation theory predicts that the peculiar velocity $`𝑽(𝒙)`$ and the density contrast $`\delta (𝒙)`$ at a given point $`𝒙`$ is statistically independent, as long as the initial density fluctuations are random Gaussian distributed. Namely, the joint probability distribution function $`P(𝑽,\delta )`$ can be written in a form as $`P_1(𝑽)P_2(\delta )`$.
It is not surprising that the peculiar velocity of each particle is largely affected by nonlinear gravitational effects and shows local density dependence described above. But what can we expect for the smoothed (bulk) velocity that is field coarse grained at some spatial scale $`R`$? Due to nonlinear mode couplings, the relation $`P_1(𝑽)P_2(\delta )`$ valid for linear theory must be modified and bulk velocity dispersion must also depend on local density contrast defined at the same smoothing scale $`R`$ (see Bernardeau 1992, Chodorowski & Łokas 1997, Bernardeau et al. 1999 for the velocity divergence field).
However, Kepner, Summers & Strauss (1997) showed from cold-dark-matter (CDM) and hot-dark-matter (HDM) N-body simulations that at nonlinear scales ($`0.77h^1\mathrm{Mpc}R4.88h^1\mathrm{Mpc}`$), such a local density dependence was not observed (see Figs.2(a) and 3(a) of their paper). This is an interesting contrast to the behavior of velocity field traced by each particle, as described before (Kepner et al. 1997, Narayanan et al. 1998).
In this article, we investigate local density dependence of smoothed (bulk) velocity dispersion using framework of second-order perturbation theory. We calculate the first-order nonlinear correction of the constrained velocity dispersion. Our target is weakly nonlinear scale and somewhat larger than scale analyzed by Kepner et al. (1997). Since current survey depth of the cosmic velocity field is highly limited, our constrained statistics might not be useful in observational cosmology at present (e.g. Seto & Yokoyama 1998b, see also Seto 1999a). Our interest in this article is theoretically motivated one about nonlinear gravitational dynamics.
As the peculiar velocity field is more weighted to large-scale fluctuations (smaller wave number $`k`$) than the density field, perturbative treatment of smoothed velocity field would be reasonable at weakly nonlinear scale. Actually, Bahcall, Gramann & Cen (1994) showed that smoothed unconstrained velocity dispersion in N-body simulations are well predicted by linear theory even at smoothing scale $`R=3h^1\mathrm{Mpc}`$ (see their Table 1). Second-order analysis by Makino, Sasaki & Suto (1992) also gives consistent results to their simulations.
## 2 Formulation
First we define the (unsmoothed) density contrast field $`\delta (𝒙)`$ in terms of the mean density of the universe $`\overline{\rho }`$ and the local density field $`\rho (𝒙)`$ as
$$\delta (𝒙)=\frac{\rho (𝒙)\overline{\rho }}{\overline{\rho }}.$$
(1)
Many theoretical predictions of the large-scale structure are based on continuous fields, but observations as well as numerical experiments (such as, N-body simulations) are usually sampled by points where point-like galaxies (or mass elements) exist. In comparison of theoretical predictions with actual observations or numerical experiments, smoothing operation becomes sometimes crucially important to remove sparseness of particles’ system. This operation is also important to reduce strong nonlinear effects which are difficult to handle theoretically. Thus it is favorable to make theoretical predictions of the large-scale structure including smoothing operation. We can express the smoothed density contrast field $`\delta _R(𝒙)`$ and the smoothed velocity field $`𝑽_R(𝒙)`$ with (spatially isotropic) filter $`W(x,R)`$ as
$$\delta _R(𝒙)𝑑𝒙^3\delta (𝒙^{})W(|𝒙𝒙^{}|,R),𝑽_R(𝒙)𝑑𝒙^3𝑽(𝒙^{})W(|𝒙𝒙^{}|,R).$$
(2)
As we discuss only the smoothed fields in this article, we hereafter omit the suffix $`R`$ which indicates smoothing radius.
The velocity dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ for points $`𝒙`$ with a given overdensity $`\delta (𝒙)=\delta `$ is formally written as
$$\mathrm{\Sigma }_V^2(\delta )=\frac{𝑽(𝒙)^2\delta _{Drc}[\delta (𝒙)\delta ]}{\delta _{Drc}[\delta (𝒙)\delta ]},$$
(3)
where $`\delta _{Drc}()`$ is Dirac’s delta function and brackets $``$ represent to take ensemble average.
We assume that the initial (linear) density fluctuations are isotropic random Gaussian. At the linear-order we have $`𝑽(𝒙)\mathrm{\Delta }^1\delta (𝒙)`$ and $`𝑽(𝒙)\delta (𝒙)=0`$ due to isotropy of matter fluctuations. This means that $`\delta `$ and $`𝑽`$ at a same point are statistically independent quantities, as a multivariate probability distribution function (hereafter PDF) of Gaussian variables is completely decided by their covariance matrix (e.g. Bardeen et al. 1986). Thus the constrained velocity dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ does not depend on the density contrast $`\delta `$ at linear order. However, nonlinear mode couplings would change the situation. Let us examine weakly non-Gaussian effects on $`\mathrm{\Sigma }_V^2(\delta )`$. We can express the first nonlinear correction of $`\mathrm{\Sigma }_V^2(\delta )`$, using framework of the Edgeworth expansion method (Cramer 1946, Matsubara 1994, 1995, Juszkiewicz et al. 1995, Bernardeau & Kofman 1995). This method is an excellent tool to explore weakly nonlinear effects of the large-scale structure induced by gravity.
When a field $`F`$ is defined by weakly non-Gaussian variables $`\{A_\mu (𝒙)\}`$ with vanishing means, we can expand the expectation value $`F`$ as (see appendix A)
$$F(A_1,\mathrm{},A_m)=F_G+\frac{1}{6}\underset{\mu ,\nu ,\lambda }{}A_\mu A_\nu A_\lambda _c\frac{^3F}{A_\mu A_\nu A_\lambda }_G+O(\sigma ^2F),$$
(4)
where $`_G`$ is the expectation value under the assumption that variables $`\{A_\mu (𝒙)\}`$ are multivariate Gaussian distributed, characterized by their covariance matrix $`A_\mu A_\nu `$. The quantity $`A_\mu A_\nu A_\lambda _c`$ is the third-order connected moment of variables $`\{A_\mu \}`$, and we have $`A_\mu A_\nu A_\lambda _c=A_\mu A_\nu A_\lambda `$ at third-order. The variance $`\sigma ^2=O(A_i^2)`$ is the order parameter of perturbative expansion around the Gaussian distribution, and we can regard $`\sigma ^2=\delta ^2`$ in this article. The denominator of $`\mathrm{\Sigma }_V^2(\delta )`$ in equation (3) is nothing but the one point PDF of density contrast $`\delta `$. From equation (4) we obtain the famous perturbative formula as follows ($`\nu \delta /\sigma `$)
$$\delta _{Drc}[\delta (𝒙)\delta ]=\frac{e^{\nu ^2/2}}{\sqrt{2\pi \sigma ^2}}\left(1+\frac{S\sigma H_3(\nu )}{6}+O(\sigma ^2)\right),$$
(5)
(e.g. Juszkiewicz et al. 1995, Bernardeau & Kofman 1995) and this is the most simplified version of the Edgeworth expansion. Here the function $`H_n(\nu )(1)^ne^{\nu ^2/2}(d/d\nu )^ne^{\nu ^2/2}`$ is $`n`$-th order Hermite polynomial, and $`S`$ is a parameter of order unity and called skewness (Peebles 1980, Fry 1984, Goroff et al. 1986, see also Seto 1999c),
$$S\frac{\delta ^3}{\sigma ^4}.$$
(6)
Due to the nonlinear correction term proportional to $`S\sigma `$, points with high-$`\sigma `$ overdensity are more abundant than the linear prediction by a Gaussian distribution.
Next the numerator of $`\mathrm{\Sigma }_V^2(\delta )`$ is given by
$$𝑽(𝒙)^2\delta _{Drc}[\delta (𝒙)\delta ]=\sigma _V^2\frac{e^{\nu ^2/2}}{\sqrt{2\pi \sigma ^2}}\left(1+\frac{S\sigma H_3(\nu )}{6}+C\sigma H_1(\nu )+O(\sigma ^2)\right),$$
(7)
where $`\sigma _V^2𝑽^2`$ is the unconstrained velocity dispersion. The parameter $`C=O(1)`$ is defined by
$$C=\frac{𝑽(𝒙)^2\delta (𝒙)}{\sigma ^2\sigma _V^2}.$$
(8)
In the studies of the large-scale structure, the Edgeworth expansion or the third-order moments have been mainly discussed for scalar fields, such as density field $`\delta (𝒙)`$ or velocity divergence field $`𝑽(𝒙)`$ (e.g. Chodorowski & Łokas 1997). Here we present analytical study for couplings of $`\delta (𝒙)`$ and $`𝑽(𝒙)`$, but numerical investigation of our method is also important as well as interesting. From equations (5) and (7) we obtain the constrained velocity dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ up to the first-order nonlinear correction as
$`\mathrm{\Sigma }_V^2(\delta )`$ $`=`$ $`{\displaystyle \frac{𝑽(𝒙)^2\delta _{Drc}[\delta (𝒙)\delta ]}{\delta _{Drc}[\delta (𝒙)\delta ]}}`$ (9)
$`=`$ $`{\displaystyle \frac{\sigma _V^2{\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi \sigma ^2}}}\left(1+{\displaystyle \frac{S\sigma H_3(\nu )}{6}}+C\sigma H_1(\nu )+O(\sigma ^2)\right)}{{\displaystyle \frac{e^{\nu ^2/2}}{\sqrt{2\pi \sigma ^2}}}\left(1+{\displaystyle \frac{S\sigma H_3(\nu )}{6}}+O(\sigma ^2)\right)}}`$
$`=`$ $`\sigma _V^2\left(1+{\displaystyle \frac{S\sigma H_3(\nu )}{6}}+C\delta {\displaystyle \frac{S\sigma H_3(\nu )}{6}}+O(\sigma ^2)\right)`$
$`=`$ $`\sigma _V^2(1+C\delta +O(\sigma ^2)).`$
Note that our result $`\mathrm{\Sigma }_V^2(\delta )`$ does not depend on the skewness parameter $`S`$. Nonlinear effects appear through the quantity $`C`$.
Next let us evaluate non-Gausssianity induced by gravity, using higher-order perturbation theory. We perturbatively expand the density and velocity fields as
$`\delta (𝒙)`$ $`=`$ $`\delta _1(𝒙)+\delta _2(𝒙)+\delta _3(𝒙)+\mathrm{},`$ (10)
$`𝑽(𝒙)`$ $`=`$ $`𝑽_1(𝒙)+𝑽_2(𝒙)+𝑽_3(𝒙)+\mathrm{},`$ (11)
where $`\delta _1(𝒙)`$ and $`𝑽_1(𝒙)`$ are the linear modes, $`\delta _2(𝒙)`$ and $`𝑽_2(𝒙)`$ are the second-order modes, and so on. We solve the following three basic equations (continuity, Euler and Poisson equations) order by order (Peebles 1980)
$`{\displaystyle \frac{}{t}}\delta (𝒙)+{\displaystyle \frac{1}{a}}[𝑽(𝒙)\{1+\delta (𝒙)\}]`$ $`=`$ $`0,`$ (12)
$`{\displaystyle \frac{}{t}}𝑽(𝒙)+{\displaystyle \frac{1}{a}}[𝑽(𝒙)]𝑽(𝒙)+{\displaystyle \frac{_ta}{a}}𝑽(𝒙)+{\displaystyle \frac{1}{a}}\varphi (𝒙)`$ $`=`$ $`0,`$ (13)
$`^2\varphi (𝒙)4\pi a^2\rho (t)\delta (𝒙)`$ $`=`$ $`0,`$ (14)
where $`a`$ represents the scale factor. In these equations we have omitted explicit time dependence of fields for notational simplicities. We only discuss quantities at a specific epoch and there would be no confusion.
Fourier space representation is convenient to analyze the nonlinear mode couplings. We denote the unsmoothed linear Fourier mode by $`\delta _{lin}(𝒌)`$. Then $`\delta _1(𝒙)`$ and $`𝑽_1(𝒙)`$ are written in terms of $`\delta _{lin}(𝒌)`$ and $`W(kR)`$, the Fourier transform of the filter function $`W(|𝒙|,R)`$, as
$$\delta _1(𝒙)=\frac{d𝒌}{(2\pi )^3}\mathrm{exp}(i𝒌𝒙)\delta _{lin}(𝒌)W(kR),𝑽_1(𝒙)=Hf\frac{d𝒌}{(2\pi )^3}\frac{i𝒌}{k^2}\mathrm{exp}(i𝒌𝒙)\delta _{lin}(𝒌)W(kR),$$
(15)
where $`H(d\mathrm{ln}a/dt`$) is the Hubble parameter and $`f(d\mathrm{ln}D/d\mathrm{ln}a`$, $`D`$: linear growth rate of density fluctuation) is a function of cosmological parameters $`\mathrm{\Omega }`$ and $`\lambda `$, and well fitted by
$$f\mathrm{\Omega }^{0.6}+\lambda /30,$$
(16)
in the ranges $`0.05\mathrm{\Omega }1.5`$ and $`0\lambda 1.5`$ (Martel 1991).
We define the linear matter power spectrum $`P(k)`$ by
$$\delta _{lin}(𝒌)\delta _{lin}(𝒍)=(2\pi )^3\delta _{Drc}^3(𝒌+𝒍)P(k).$$
(17)
Then the dispersions $`\sigma ^2`$ and $`\sigma _V^2`$ are given by the following simple integrals of $`P(k)`$ up to required order to evaluate the first nonlinear effects of $`\mathrm{\Sigma }_V^2(\delta )`$,
$`\sigma ^2`$ $`=`$ $`{\displaystyle \frac{d𝒌}{(2\pi )^3}P(k)W(kR)^2}+O(\sigma ^4),`$ (18)
$`\sigma _V^2`$ $`=`$ $`H^2f^2{\displaystyle \frac{d𝒌}{(2\pi )^3k^2}P(k)W(kR)^2}+O(\sigma ^4),`$ (19)
In this article we only use the Gaussian filter defined by $`W(kR)=\mathrm{exp}[(kR)^2/2]`$.
As shown in equation (8), the first-order nonlinear correction of the constrained dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ is characterized by the factor $`C`$. We need the second-order modes $`\delta _2(𝒙)`$ and $`𝑽_2(𝒙)`$ to calculate the first nonvanishing contributions of $`𝑽(𝒙)^2\delta (𝒙)`$. These second-order modes are given with linear mode $`\delta _{lin}(𝒌)`$ as (Fry 1984, Goroff 1986)
$`\delta _2(𝒙)`$ $`=`$ $`{\displaystyle \frac{d𝒌d𝒍}{(2\pi )^6}\mathrm{exp}[i(𝒌+𝒍)𝒙]\delta _{lin}(𝒌)\delta _{lin}(𝒍)_{2\delta }(𝒌,𝒍)W(R|𝒌+𝒍|)},`$ (20)
$`𝑽_2(𝒙)`$ $`=`$ $`Hf{\displaystyle \frac{d𝒌d𝒍}{(2\pi )^6}\frac{i(𝒌+𝒍)}{|𝒌+𝒍|^2}\mathrm{exp}[i(𝒌+𝒍)𝒙]\delta _{lin}(𝒌)\delta _{lin}(𝒍)_{2V}(𝒌,𝒍)W(R|𝒌+𝒍|)},`$ (21)
where kernels $`_{2\delta }`$ and $`_{2V}`$ are defined as follows
$`_{2\delta }(𝒌,𝒍)`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+K)+{\displaystyle \frac{𝒌𝒍}{2kl}}\left({\displaystyle \frac{k}{l}}+{\displaystyle \frac{l}{k}}\right)+{\displaystyle \frac{1}{2}}(1K)\left({\displaystyle \frac{𝒌𝒍}{kl}}\right)^2,`$ (22)
$`_{2V}(𝒌,𝒍)`$ $`=`$ $`L+{\displaystyle \frac{𝒌𝒍}{2kl}}\left({\displaystyle \frac{k}{l}}+{\displaystyle \frac{l}{k}}\right)+(1L)\left({\displaystyle \frac{𝒌𝒍}{kl}}\right)^2.`$ (23)
The factors $`K`$ and $`L`$ depend very weakly on cosmological parameters $`\mathrm{\Omega }`$ and $`\lambda `$, and are fitted as (Matsubara 1995)
$`K(\mathrm{\Omega },\lambda )`$ $``$ $`{\displaystyle \frac{3}{7}}\mathrm{\Omega }^{1/30}{\displaystyle \frac{\lambda }{80}}\left(1{\displaystyle \frac{3}{2}}\lambda \mathrm{log}_{10}\mathrm{\Omega }\right),`$ (24)
$`L(\mathrm{\Omega },\lambda )`$ $``$ $`{\displaystyle \frac{3}{7}}\mathrm{\Omega }^{11/200}{\displaystyle \frac{\lambda }{70}}\left(1{\displaystyle \frac{7}{3}}\lambda \mathrm{log}_{10}\mathrm{\Omega }\right),`$ (25)
in the ranges $`0.1\mathrm{\Omega }1`$ and $`0\lambda 1`$. In the followings we neglect these weak dependence and simply put
$$K=L=\frac{3}{7}.$$
(26)
Thus we can write down the third-order moment $`𝑽𝑽\delta `$ in the following form
$`𝑽𝑽\delta `$ $`=`$ $`𝑽_1𝑽_1\delta _2+2𝑽_1𝑽_2\delta _1+O(\sigma ^6)`$ (28)
$`=`$ $`2H^2f^2{\displaystyle \frac{d𝒌d𝒍}{(2\pi )^6}P(k)P(l)\left[\frac{𝒌𝒍}{k^2l^2}_{2\delta }(𝒌,𝒍)+2\frac{𝒌(𝒌+𝒍)}{k^2|𝒌+𝒍|^2}_{2V}(𝒌,𝒍)\right]}`$
$`\times W(kR)W(lR)W(|𝒌+𝒍|R)+O(\sigma ^6).`$
Due to the rotational symmetry around the origin, we can simplify the six dimensional integral $`d𝒌d𝒍`$ to three dimensional integral $`dkdldu`$. Here, $`1u1`$ is the cosine between two vectors $`𝒌`$ and $`𝒍`$ and given by $`u=𝒌𝒍/kl`$. Then we obtain the first nonvanishing order of $`𝑽𝑽\delta `$ as
$`𝑽𝑽\delta `$ $`=`$ $`2H^2f^2{\displaystyle _1^1}𝑑u{\displaystyle \frac{k^2l^2dkdl}{8\pi ^4}P(k)P(l)\mathrm{exp}[k^2l^2klu]}`$ (29)
$`\times [{\displaystyle \frac{u}{kl}}\{{\displaystyle \frac{5}{7}}+{\displaystyle \frac{u}{2}}({\displaystyle \frac{k}{l}}+{\displaystyle \frac{l}{k}})+{\displaystyle \frac{2}{7}}u^2\}`$
$`+2{\displaystyle \frac{k+lu}{k(k^2+l^2+2klu)}}\{{\displaystyle \frac{3}{7}}+{\displaystyle \frac{u}{2}}({\displaystyle \frac{k}{l}}+{\displaystyle \frac{l}{k}})+{\displaystyle \frac{4}{7}}u^2\}].`$
Note that the parameter $`C`$ does not depend on the normalization of power spectrum (see eqs. and ). Furthermore, the factors $`Hf`$ cancel out between $`\sigma _V^2`$ and $`𝑽𝑽\delta `$ and cosmological parameters are irrelevant for the factor $`C`$ in our treatment ($`K=L=3/7`$). Finally, we comment that even though the constrained dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ changes by $`\sigma _V^2C\delta `$ from the unconstrained value $`\sigma _V^2`$, the shape of the one-point PDF of velocity field with a given $`\delta `$ keeps Gaussian distribution at the same order of nonlinearity. We can easily confirm this by calculating the ratio
$$\frac{\delta _{Drc}^3(𝑽(𝒙)𝑽)\delta _{Drc}(\delta (𝒙)\delta )}{\delta _{Drc}(\delta (𝒙)\delta )}=\frac{1}{(23^1\pi \mathrm{\Sigma }_V^2(\delta ))^{3/2}}\left[\mathrm{exp}\left(\frac{𝑽^2}{23^1\mathrm{\Sigma }_V^2(\delta )}\right)+O(\sigma ^2)\right].$$
(30)
The factor $`3^1`$ in the right-hand side arises from the dimensionality of the velocity vector. First-order correction is completely absorbed to the velocity dispersion $`\mathrm{\Sigma }_V^2(\delta )`$.
## 3 Results
In this section we numerically evaluate the parameter $`C`$ for various power spectra. We first examine pure power-law spectra $`P(k)`$ given by $`(n>1)`$
$$P(k)=Ak^n.$$
(31)
In this case $`C`$ does not depend on the smoothing radius $`R`$, and we can simply put $`R=1`$. Then the dispersions $`\sigma ^2`$ and $`\sigma _V^2`$ are given as
$$\sigma ^2=\frac{1}{(2\pi )^2}\mathrm{\Gamma }\left(\frac{3+n}{2}\right),\sigma _V^2=\frac{(Hf)^2}{(2\pi )^2}\mathrm{\Gamma }\left(\frac{1+n}{2}\right),$$
(32)
where $`\mathrm{\Gamma }(n)`$ is the Gamma function. As for the nonlinear coupling $`𝑽𝑽\delta =𝑽_1𝑽_1\delta _2+2𝑽_1𝑽_2\delta _1`$, we can write down the first contribution $`𝑽_1𝑽_1\delta _2`$ explicitly in terms of Hypergeometric functions as in the case of skewness parameter $`S`$ (Matsubara 1994, Łokas et al 1995). However, using mathematica (Wolfram 1996), we confirm that the second term $`2𝑽_1𝑽_2\delta _1`$ cannot be expressed in a closed form and numerical integration is required. These two terms diverge in the limit $`n1`$ where velocity dispersion $`\sigma _V^2`$ also diverges, but the factor $`C`$ approaches $`0`$ in this limit.
In Fig.1 we plot $`C`$ as a function of spectral index $`n`$ in the range $`1<n<2`$. The correction $`C`$ is a positive and increasing function of $`n`$. This means that the velocity dispersion of high density regions are larger than that of low density regions. We have $`C=0.314`$ at the scale-invariant spectrum $`n=1`$.
Next we examine $`C`$ for a more realistic power spectrum $`P(k)`$. We use CDM transfer function given in Bardeen, Bond, Kaiser & Szalay (1986) and assume that the primordial spectral index is equal to 1. Then $`P(k)`$ can be written as
$$P(k)=Ak\left[\frac{\mathrm{ln}(1+2.34q)}{2.34q}\right]^2[1+3.89q+(16.1q)^2+(5.46q)^3+(6.71q)^4]^{1/2},$$
(33)
where $`qk/[(\mathrm{\Gamma }h)\mathrm{Mpc}^1]`$. $`\mathrm{\Gamma }`$ is the shape parameter of the CDM transfer function and recent observational analyses of galaxy clusterings support $`\mathrm{\Gamma }=0.20.3`$ (e.g. Tadros et al. 1999, Dodelson & Gazta$`\stackrel{~}{\mathrm{n}}`$aga 1999). In Fig.2 we plot $`C`$ as a function of smoothing radius $`R`$ in units of $`[(\mathrm{\Gamma }h)^1\mathrm{Mpc}]`$. For this model the factor $`C`$ depends weakly on the smoothing radius $`R`$ and we have $`C0.30`$ at a weakly nonlinear regime $`R10h^1\mathrm{Mpc}`$. In the limit $`R\mathrm{}`$, $`C`$ converges to $`0.314`$ which is the same value of $`C`$ for the power-law model with $`n=1`$ presented in Fig.1. This is reasonable as we have
$$\underset{k0}{lim}\frac{P(k)}{k}=const,$$
(34)
for CDM models analyzed here.
Our results obtained so far are the velocity dispersion for points constrained by the matter density contrast $`\delta `$. One might have interest in the velocity dispersion constrained by the galaxy density contrast $`\delta _g`$. Here, let us assume deterministic but nonlinear biasing relation for the smoothed galaxy distribution $`\delta _g(𝒙)`$ and the matter distribution $`\delta (𝒙)`$ as
$$\delta _g(𝒙)=b_1\delta (𝒙)+b_2(\delta (𝒙)^2\sigma ^2)+O(\sigma ^3),$$
(35)
where $`b_1`$ and $`b_2`$ are some constants (e.g. Fry & Gazta$`\stackrel{~}{\mathrm{n}}`$aga 1993). In this case we can easily show that the velocity dispersion $`\mathrm{\Sigma }_V(\delta _g)^2`$ for points $`𝒙`$ with $`\delta _g(𝒙)=\delta _g`$ is given by
$$\mathrm{\Sigma }_V^2(\delta _g)=\sigma _V^2(1+C\delta _g/b_1+O(\sigma ^2)),$$
(36)
where the factors $`C`$ and $`\sigma _V`$ are same as those appeared in $`\mathrm{\Sigma }_V^2(\delta )`$ (eq.). Thus $`\mathrm{\Sigma }_V^2(\delta _g)`$ does not depend on the nonlinear coefficient $`b_2`$. This is also apparent when we write down $`\delta (𝒙)`$ using $`\delta _g(𝒙)`$ and then insert this solution to equation (9). The factor proportional to $`b_2`$ is higher effects than analyzed here. Note that in equation (36), the linear bias parameter $`b_1`$ appears by itself not in the usual form $`\beta \mathrm{\Omega }^{0.6}/b_1`$, and the overdensity $`\delta _g`$ dependence becomes smaller for larger $`b_1`$.
Kepner et al. (1997) numerically investigated the mean magnitude $`|𝑽(𝒙)|`$ of smoothed bulk velocity for points with given overdensity $`\delta `$. Following the fact commented in the last paragraph of section 2, we can easily calculate this magnitude $`\mu _V(\delta )`$ and obtain following result (see appendix B)
$$\mu _V(\delta )=\sqrt{\frac{8}{3\pi }\mathrm{\Sigma }_V^2(\delta )(1+O(\sigma ^2))}=\sqrt{\frac{8}{3\pi }}\sigma _V\left(1+\frac{C\delta }{2}+O(\sigma ^2)\right).$$
(37)
For a typical CDM model, our analytical result predicts that the magnitude $`\mu _V(\delta )`$ is expected to change $`0.15\delta `$, according to the local density contrast $`\delta `$. If we constrain points using overdensity of galaxies instead of that of the gravitating matter, the combination $`C\delta `$ is replaced by $`C\delta _g/b_1`$ in the above equation.
Numerical results of Kepner et al. (1997) were given for CDM and HDM models with $`\sigma _8=0.67`$ normalization. Here $`\sigma _8`$ is the linear rms density fluctuation in a sphere of $`8h^1\mathrm{Mpc}`$ radius. They calculated the smoothed density and velocity fields with smoothing radius $`R`$ at $`0.77h^1\mathrm{Mpc}R4.88h^1\mathrm{Mpc}`$. Thus their results are quantities at nonlinear regimes. It is true that simple application of our perturbative formula to their results would not be valid. However, surprisingly enough, the quantity $`\mu _V(\delta )`$ shows almost no $`\delta `$ dependence in the range $`0<\delta \stackrel{<}{}\text{ }30`$. <sup>1</sup><sup>1</sup>1It seems that the function $`\mu _V(\delta )`$ in their figures shows extremely weak dependence of $`\delta `$ around $`\delta 0`$.
If $`\mathrm{\Sigma }_V^2(\delta )`$ (and thus $`\mu _V(\delta )`$) shows no $`\delta `$-dependence at nonlinear scale and our second-order analysis is valid for weakly nonlinear regime, we confront an interesting possibility. Namely, with parameterization of overdensity by normalized value $`\nu \delta /\sigma `$, velocity dispersion does not depend on $`\nu `$ at linear and nonlinear regime, but depends on it at (intermediate) weakly nonlinear regime. Furthermore, we should notice that the velocity dispersion of dark-matter particles (without coarse graining) depends largely on $`\delta `$ (Kepner et al. 1997, Narayanan et al. 1998).<sup>2</sup><sup>2</sup>2Definitions of velocity dispersion in these two papers are not identical.
To make clear understanding of these transitions, we need to numerically investigate the constrained dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ in detailed manner with various smoothing length $`R`$, from linear to nonlinear scales. Performance of second-order perturbation theory for the velocity vector is also worth studying.
## 4 Summary
It is commonly accepted that the large-scale structure observed today is formed by gravitational instability from small primordial density fluctuations (Peebles 1980). In this picture, the peculiar velocity and the density contrast are fundamental quantities to characterize inhomogeneities in the universe. Linear analysis of cosmological perturbation theory predicts that, as long as initial fluctuations are random Gaussian distributed, the one-point PDF of the velocity field $`𝑽(𝒙)`$ is statistically independent of the local density contrast $`\delta (𝒙)`$. This is an important aspect of cosmological perturbation.
However nonlinear gravitational evolution changes the situation. Due to nonlinear mode-couplings, the peculiar velocity field is no longer statistically independent of the local density field. Here we have investigated bulk velocity dispersion ($`\mathrm{\Sigma }_V^2(\delta )`$) as a function of the local density contrast and calculated its first nonlinear correction using framework of second-order perturbation theory. Our target has been set at weakly nonlinear regimes where perturbative treatment must be reasonable. At present, survey depth of velocity field is highly limited and our constrained statistics might not be directly useful for observational cosmology. However, we believe that our theoretical study is important to understand one interesting aspect of the cosmic velocity field peculiar to its nonlinear evolution.
We have shown that the first nonlinear correction of $`\mathrm{\Sigma }_V^2(\delta )`$ is proportional to the local density and strongly depends on the matter power spectra, but weakly on the cosmological parameters $`\mathrm{\Omega }`$ and $`\lambda `$. For typical CDM model with primordially scale-invariant fluctuations, this first-order correction is about $`0.3\delta `$. If we use overdensity of galaxies $`\delta _g`$, this correction term becomes $`0.3\delta _g/b_1`$ (the factor $`b_1`$ is defined in eq.). This dependence might be used to constrain the linear bias parameter $`b_1`$ itself (not in the usual form $`\mathrm{\Omega }^{0.6}/b_1`$) in future peculiar velocity surveys. We have also shown that the constrained one-point PDF of velocity field keeps the Gaussian shape up to second-order of perturbation. First nonlinear effects are completely absorbed to the velocity dispersion $`\mathrm{\Sigma }_V^2(\delta )`$.
Numerical results by Kepner et al. (1997) have been compared with our analytical results. Their results show almost no $`\delta `$-dependence, contrary to ours. However spatial scale of their analysis (strongly nonlinear regime) is largely different from ours (weakly nonlinear regime). In the forthcoming paper, detailed numerical analysis will be presented for various smoothing lengths $`R`$ and power spectra (see also Seto 2000). Validity of second-order analysis as well as the Edgeworth expansion method for velocity vector would be also investigated numerically.
The author would like to thank J. Yokoyama for useful discussion and the referee R. Juszkiewicz for helpful comments. This work was partially supported by the Japanese Grant in Aid for Science Research Fund of the Ministry of Education, Science, Sports and Culture No. 3161.
## Appendix A Weakly Non-Gaussian Averages
In this appendix, we derive expression (4) for weakly non-Gaussian variables $`\{A_\mu \}(\mu =1,\mathrm{},n)`$ with $`A_\mu =0`$ (e.g. Matsubara 1994). The partition function $`Z(J_\mu )`$ for a multivariate probability distribution function $`P(A_\mu )`$ is defined as
$$Z(J_\mu )_{\mathrm{}}^{\mathrm{}}d^nP(A_\mu )\mathrm{exp}\left(iJ_\nu A_\nu \right).$$
(A1)
According to the cumulant expansion theorem (Ma 1985), the function $`\mathrm{ln}Z(J_\mu )`$ is a generating function of connected moments $`A_{\mu _1}\mathrm{}A_{\mu _N}_c`$. Therefore, taking the inverse Fourier transform of equation (A1), the probability distribution function $`P(A_\mu )`$ is written as
$$P(A_\mu )=\mathrm{exp}\left(\underset{N=3}{\overset{\mathrm{}}{}}\frac{(1)^N}{N!}\underset{\mu _1,\mathrm{},\mu _N}{}A_{\mu _1}\mathrm{}A_{\mu _N}_c\frac{^N}{A_{\mu _1}A\mathrm{}A_{\mu _N}}\right)P_G(A_\mu ),$$
(A2)
where the function $`P_G(A_\mu )`$ is the multivariate Gaussian probability distribution function determined by a ($`n\times n`$) correlation matrix $`M_{\mu \nu }A_\mu A_\nu `$ as
$$P_G(A_\mu )=\frac{1}{\sqrt{(2\pi )^n\mathrm{det}M}}\mathrm{exp}\left(\frac{1}{2}\underset{\mu ,\nu }{}A_\mu (M^1)_{\mu \nu }A_\nu \right).$$
(A3)
If we have relations $`A_{\mu _1}\mathrm{}A_{\mu _N}_c=O(\sigma ^{2N2})`$ as predicted by higher-order perturbation theory, equation (A2) is perturbatively expanded as
$$P(A_\mu )=P_G(A_\mu )\frac{1}{6}\underset{\mu ,\nu ,\lambda }{}A_\mu A_\nu A_\lambda _c\frac{^3}{A_\mu A_\nu A_\lambda }P_G(A_\mu )+O(\delta ^2).$$
(A4)
Evaluating the ensemble average of a field $`F(A_\mu )`$ with this perturbative expression and taking partial integrals, we obtain expansion (4).
One might consider that formula (3) with Dirac’s delta function is somewhat indirect. We can obtain same results using formula for probability distribution function $`P(𝑽|\delta )=P(𝑽,\delta )/P(\delta )`$ and evaluating $`P(𝑽,\delta )`$ and $`P(\delta )`$ with expression (A4).
## Appendix B Derivation of $`\mu _V(\delta )`$
Here we derive expression (37). As in the case of the constrained velocity dispersion $`\mathrm{\Sigma }_V^2(\delta )`$ (eq.), the quantity $`\mu _V(\delta )`$ is formally defined as
$$\mu _V(\delta )\frac{|𝑽(𝒙)|\delta _{Drc}[\delta (𝒙)\delta ]}{\delta _{Drc}[\delta (𝒙)\delta ]}.$$
(B1)
The r.h.s. of this equation is written as
$$\mu _V(\delta )=_{\mathrm{}}^{\mathrm{}}𝑑𝑽\frac{|𝑽|\delta _{Drc}^3(𝑽(𝒙)𝑽)\delta _{Drc}(\delta (𝒙)\delta )}{\delta _{Drc}(\delta (𝒙)\delta )}.$$
(B2)
With the perturbative expansion given in equation (30) the r.h.s. of this equation is evaluated as
$$\frac{1}{(23^1\pi \mathrm{\Sigma }_V^2(\delta ))^{3/2}}_{\mathrm{}}^{\mathrm{}}𝑑𝑽\left[\mathrm{exp}\left(\frac{𝑽^2}{23^1\mathrm{\Sigma }_V^2(\delta )}\right)+O(\sigma ^2)\right]|𝑽|.$$
(B3)
It is straightforward to obtain expression (37) given in the main text as follows
$$\mu _V(\delta )=\sqrt{\frac{8}{3\pi }\mathrm{\Sigma }_V^2(\delta )(1+O(\sigma ^2))}=\sqrt{\frac{8}{3\pi }}\sigma _V\left(1+\frac{C\delta }{2}+O(\sigma ^2)\right).$$
(B4) |
warning/0002/hep-th0002167.html | ar5iv | text | # Vector supersymmetry in topological field theories.
## 1 Introduction
It is well known that there exist two classes of topological quantum field theories (TQFT’s) . The so-called Witten-type models, e.g. topological Yang-Mills (YM) theory on a four-dimensional space-time manifold - and its higher-dimensional generalizations and the Schwarz-type models, e.g. the Chern-Simons theory in any odd space-time dimension and the BF model in arbitrary dimensions. These models are not only invariant under BRST-transformations, but also (at least in certain gauges) under the so-called vector supersymmetry (VSUSY). The anticommutator of this transformation with the BRST-operator yields space-time translations (either off-shell or on-shell) and thus generates a superalgebra of Wess-Zumino type . The VSUSY plays a central role for proving the perturbative finiteness of Schwarz-type models (e.g. see ) and it is most helpful for discussing the algebraic renormalization of Witten-type models .
In reference , a general procedure for obtaining the explicit form of VSUSY-transformations was presented in the case of three-dimensional Chern-Simons theory. However, this derivation becomes rather involved for more complex models. Another approach to VSUSY consists of writing the most general Ansatz for these transformations which is compatible with dimensional and covariance constraints and subsequently eliminating terms by imposing the aforementioned anticommutation relations and the invariance of the action. In practice, this also turns out to be a laborious undertaking.
The aim of the present work is to give a short derivation of the two different forms of VSUSY’s appearing in topological models. Our results also allow us to cast these transformations into a compact form. The proposed procedure closely parallels the derivation of BRST-transformations in field theories with local symmetries from the so-called horizontality conditions or Russian formulae . The latter enclose all field variations in a single (or a few) simple equation(s). Though the horizontality conditions admit a clear geometric interpretation in the case of ordinary YM-theories , they seem to be a bit mysterious for more general field theoretic models. Therefore, we will first provide some insight into their working mechanism. In particular, we will emphasize that they not only encode all information concerning the explicit form of BRST-transformations, but also about their nilpotency.
Our paper is organized as follows. Section 2 is devoted to reviewing the horizontality conditions for some of the prototype models of TQFT mentioned above. In section 3, we discuss the VSUSY before presenting our general derivation of VSUSY-transformations in section 4. Some comments are gathered in a concluding section. We note that all of our considerations concern the classical theory (tree-level).
## 2 BRST-symmetry in topological field theories
As prototype models, we consider topological YM-theory on a Riemannian $`4`$-manifold $`_4`$, Chern-Simons theory on $`𝐑^3`$ and the BF model on a Riemannian manifold $`_{p+2}`$ of dimension $`p+2`$. All fields in these models are given by differential forms with values in a Lie algebra. The local symmetries will be described in the BRST-framework, i.e. infinitesimal symmetry parameters are turned into ghost fields and symmetry transformations are collected in a BRST-transformation. Thus, all fields are characterized by their form degree and ghost-number which we specify by means of lower and upper indices for the fields. All models under consideration involve a YM-connection one-form<sup>1</sup><sup>1</sup>1Here, the $`T_a`$ represent a basis of the Lie algebra and they are assumed to satisfy $`[T_a,T_b]=\mathrm{i}f_{ab}^cT_c`$ and $`\mathrm{tr}(T_aT_b)=\delta _{ab}`$. $`A=A_\mu ^aT_adx^\mu `$. The associated field strength is defined by
$$F=dA+\frac{1}{2}[A,A],$$
(1)
where $`[.,.]`$ denote graded brackets. From the nilpotency of the exterior derivative $`d`$, it follows that $`F`$ satisfies the Bianchi identity: $`DFdF+[A,F]=0`$.
It is convenient to combine the gauge field and the associated ghost $`c`$ as well as the exterior derivative and the BRST-operator $`s`$ by introducing the expressions
$`\stackrel{~}{A}`$ $`=`$ $`A+c`$ (2)
$`\stackrel{~}{d}`$ $`=`$ $`d+s.`$
Accordingly, one defines
$`\stackrel{~}{F}`$ $`=`$ $`\stackrel{~}{d}\stackrel{~}{A}+{\displaystyle \frac{1}{2}}[\stackrel{~}{A},\stackrel{~}{A}]`$ (3)
$`\stackrel{~}{D}`$ $`=`$ $`\stackrel{~}{d}+[\stackrel{~}{A},].`$
By definition, the $`s`$-operator anticommutes with $`d`$, and $`d`$ is nilpotent, henceforth
$`\stackrel{~}{d}^2=s^2\text{on all fields}.`$ (4)
By expanding $`\stackrel{~}{F}`$ with respect to the ghost-number, we find that it has an expression of the form
$`\stackrel{~}{F}=F_2^0+F_1^1+F_0^2,`$ (5)
where
$`F_2^0`$ $``$ $`F`$
$`F_1^1`$ $``$ $`sA+Dc`$ (6)
$`F_0^2`$ $``$ $`sc+{\displaystyle \frac{1}{2}}[c,c].`$
From equations (2)-(4), it follows that
$`\stackrel{~}{D}\stackrel{~}{F}=\stackrel{~}{d}^2\stackrel{~}{A}=s^2A+s^2c`$ (7)
and similarly equations (3)-(5) imply
$`\stackrel{~}{D}(\stackrel{~}{D}\stackrel{~}{F})=\stackrel{~}{d}^2\stackrel{~}{F}=s^2F+s^2F_1^1+s^2F_0^2.`$ (8)
So far, we have simply derived some equations involving $`s`$-variations of $`A`$ and $`c`$ without specifying the latter. According to equations (5),(6), these can be determined by imposing a horizontality condition, i.e. by prescribing $`F_1^1`$ and $`F_0^2`$ in equation (5), $`F_2^0`$ being necessarily equal to $`FdA+\frac{1}{2}[A,A]`$. Equations (7) and (8) then allow us to discuss the nilpotency of the resulting BRST-transformations and thereby to check the consistency of the imposed horizontality condition.
### 2.1 Witten-type models
#### Topological Yang-Mills theory
The classical action reads
$$\mathrm{\Sigma }_{inv}^W=__4\mathrm{tr}\left(FF\right),$$
(9)
where the wedge product symbol has been omitted.
Due to the shift- (or topological $`Q`$-) symmetry $`\delta A=\psi _1^1`$ which is present in this type of model -, the connection $`A`$ is associated with ghost fields $`\psi _1^1`$ and $`\phi _0^2\phi ^2`$. Henceforth, one imposes the horizontality condition
$`\stackrel{~}{F}=F+\psi _1^1+\phi ^2.`$ (10)
Substitution of equations (5),(6) into this relation yields the BRST-transformations
$`sA`$ $`=`$ $`\psi _1^1Dc`$
$`sc`$ $`=`$ $`\phi ^2{\displaystyle \frac{1}{2}}[c,c].`$ (11)
Since $`sA`$ involves an inhomogeneous term $`\psi _1^1`$, the requirement $`s^2A=0`$ determines $`s\psi _1^1`$ in terms of the other fields and analogously the condition $`s^2c=0`$ determines $`s\phi ^2`$. In order to obtain the explicit form of these variations, as well as the one of $`F`$, we note that substitution of $`s^2A=0=s^2c`$ in equation (7) yields the generalized Bianchi identity $`\stackrel{~}{D}\stackrel{~}{F}=0`$. By expanding the latter with respect to the ghost-number, one readily obtains
$`sF`$ $`=`$ $`D\psi _1^1[c,F]`$ (12)
$`s\psi _1^1`$ $`=`$ $`D\phi ^2[c,\psi _1^1]`$
$`s\phi ^2`$ $`=`$ $`[c,\phi ^2].`$
Finally, substitution of $`\stackrel{~}{D}\stackrel{~}{F}=0`$ into equation (8) allows us to conclude that $`s^2(F,\psi _1^1,\phi ^2)=0`$. Last, but not least, it can be verified explicitly that the transformation (12) of $`F`$ leaves the classical action (9) invariant.
### 2.2 Schwarz-type models
#### Chern-Simons theory
The classical action of this model is given by
$$\mathrm{\Sigma }_{inv}^{CS}=\frac{1}{2}_{𝐑^3}\mathrm{tr}(AdA+\frac{1}{3}A[A,A]).$$
(13)
Due to the absence of the shift-symmetry in this model, one imposes the horizontality condition
$$\stackrel{~}{F}=F,$$
(14)
which implies
$`sA`$ $`=`$ $`Dc`$
$`sc`$ $`=`$ $`{\displaystyle \frac{1}{2}}[c,c].`$ (15)
From equations (11) and (12), we see that the truncation (14) of (5) is consistent and that it leads to nilpotent $`s`$-variations.
#### BF model
Apart from a YM $`1`$-form $`A`$ and its associated ghost $`c`$, this model involves a $`p`$-form $`BB_p^0`$ transforming with the adjoint representation of the gauge group. Its field strength $`HDB`$ automatically satisfies the second Bianchi identity $`DH=[F,B]`$. The model is characterized by the classical action
$$\mathrm{\Sigma }_{inv}^{BF}=_{_{p+2}}\mathrm{tr}\left(BF\right),$$
(16)
which is not only invariant under ordinary gauge transformations, but also under the reducible local symmetry $`\delta B=DB_{p1}^1`$. Henceforth, the field $`B`$ is associated with a series of ghosts $`B_{p1}^1,B_{p2}^2,\mathrm{},B_0^p`$ which can be combined in a generalized field $`\stackrel{~}{B}`$, by analogy to the definition of $`\stackrel{~}{A}`$ in the YM-sector:
$$\stackrel{~}{B}=B+B_{p1}^1+\mathrm{}+B_1^{p1}+B_0^p.$$
(17)
Then, the generalized field strength
$$\stackrel{~}{H}\stackrel{~}{D}\stackrel{~}{B}$$
admits the expansion
$$\stackrel{~}{H}=H_{p+1}^0+H_p^1+\mathrm{}+H_0^{p+1},$$
(18)
where
$`H_{p+1}^0`$ $`=`$ $`DB_p^0`$
$`H_p^1`$ $`=`$ $`DB_{p1}^1+sB_p^0+[c,B_p^0]`$
$`\mathrm{}`$
$`H_0^{p+1}`$ $`=`$ $`sB_0^p+[c,B_0^p].`$
Due to the absence of shift-symmetries, we proceed by analogy with the Chern-Simons model and truncate the expansions of the field strengths, i.e. we impose the horizontality conditions
$`\stackrel{~}{F}`$ $`=`$ $`F`$
$`\stackrel{~}{H}`$ $`=`$ $`H.`$ (20)
From equations (2.2), we then obtain the BRST-transformations
$`sB_{pq}^q`$ $`=`$ $`DB_{pq1}^{q+1}[c,B_{pq}^q]\mathrm{for}q=0,1,\mathrm{},p1`$
$`sB_0^p`$ $`=`$ $`[c,B_0^p],`$ (21)
the transformations of the connection $`A`$ and ghost field $`c`$ being given by equations (15).
In order to check the nilpotency of the $`s`$-variations (21), we note that
$$s^2B+s^2B_{p1}^1+\mathrm{}+s^2B_0^p=\stackrel{~}{d}^2\stackrel{~}{B}=\stackrel{~}{D}\stackrel{~}{H}[\stackrel{~}{F},\stackrel{~}{B}],$$
which implies (by matching the ghost-numbers on the left and right hand sides)
$`s^2B_{pq}^q`$ $`=`$ $`[F,B_{pq2}^{q+2}]\mathrm{for}q=0,1,\mathrm{},p2`$ (22)
$`s^2B_{pq}^q`$ $`=`$ $`0\mathrm{for}q=p1,p.`$
Here, the right hand side vanishes, if we use the equation of motion $`F=0`$ following from the classical action (16). Thus, the $`s`$-variations (21) are only nilpotent on-shell . The origin of this result can be drawn back to the fact that the used horizontality conditions (which were motivated by the absence of shift-symmetries) enforced a truncation of the ghost-expansion of $`\stackrel{~}{H}`$, which is not consistent in the sense that it leads to an on-shell algebra. This kind of phenomenon is familiar from supersymmetric field theories where the elimination of auxiliary fields from a superfield expansion leads to a supersymmetry algebra which only closes on-shell. In the BRST-framework, the on-shell closure of the symmetry algebra is reflected by the fact that the $`s`$-variations of the ghost fields are only nilpotent on-shell. From our discussion, it followed that this information can be directly extracted from the horizontality conditions.
## 3 Vector supersymmetry
In flat space-time, infinitesimal VSUSY-transformations are parametrized by a constant vector field. On curved space-time manifolds, one has to consider a covariantly constant vector field . Although this can be done at the expense of technical complications, we will limit our discussion of VSUSY to flat space-time for the sake of simplicity.
The total action $`\mathrm{\Sigma }=\mathrm{\Sigma }_{inv}+\mathrm{\Sigma }_{gf}`$ of a topological model not only involves classical and ghost fields, but also anti-ghost and Lagrange multiplier fields<sup>2</sup><sup>2</sup>2Here, “classical” fields are not opposed to quantum fields, but simply refer to the fields appearing in the classical action.. Let us denote all these fields collectively by $`(\mathrm{\Phi }_i)_{i=1,2,\mathrm{}}`$. Their infinitesimal VSUSY-variations $`\delta _\tau \mathrm{\Phi }_i\tau ^\mu \delta _\mu \mathrm{\Phi }_i`$ are parametrized by a constant, $`s`$-invariant vector field $`\tau =\tau ^\mu _\mu `$ of ghost-number zero. The operator $`\delta _\tau `$ acts as an antiderivation which lowers the ghost-number by one unit and which anticommutes with $`d`$.
The existence and explicit form of VSUSY-transformations for a topological model described by the action $`\mathrm{\Sigma }`$ depends, in general, in a sensitive way on the choice of the gauge-fixing condition. In order to study the existence of this symmetry and to determine the explicit form of VSUSY-transformations, one can apply a general procedure presented in reference . This method is based on the facts that the gauge-fixing term is BRST-exact and that it is metric dependent (while the classical action is metric independent): thus, the energy-momentum tensor is a BRST-exact expression,
$$T_{\mu \nu }=s\mathrm{\Lambda }_{\mu \nu },$$
(23)
which is a typical feature of topological models . After determining $`\mathrm{\Lambda }_{\mu \nu }`$, one expresses $`^\mu \mathrm{\Lambda }_{\mu \nu }`$ in terms of the functional derivatives $`\delta \mathrm{\Sigma }/\delta \mathrm{\Phi }_i`$, thereby producing contact terms, i.e. expressions which vanish when the equations of motion are used. If
$$^\mu \mathrm{\Lambda }_{\mu \nu }=\text{contact terms}+^\mu \epsilon _{\mu \nu },\mathrm{where}s(^\mu \epsilon _{\mu \nu })=0,$$
(24)
then the quantities
$`\widehat{\mathrm{\Lambda }}_{\mu \nu }`$ $`=`$ $`\mathrm{\Lambda }_{\mu \nu }\epsilon _{\mu \nu }`$ (25)
$`\widehat{T}_{\mu \nu }`$ $`=`$ $`T_{\mu \nu }s\epsilon _{\mu \nu }`$
are conserved up to equations of motion. They are still related by $`\widehat{T}_{\mu \nu }=s\widehat{\mathrm{\Lambda }}_{\mu \nu }`$ and $`\widehat{T}_{\mu \nu }`$ can be viewed as an improvement of the energy-momentum tensor. More explicitly, one finds
$$^\nu \widehat{\mathrm{\Lambda }}_{\nu \mu }=\mathrm{tr}(V_{i\mu }\frac{\delta \mathrm{\Sigma }}{\delta \mathrm{\Phi }_i}),$$
where $`V_{i\mu }`$ are polynomials in the fields $`\mathrm{\Phi }_i`$ and their derivatives. Integration of the last equation over the $`n`$-dimensional space-time on which the topological model is defined, yields
$$0=__nd^nx\mathrm{tr}(V_{i\mu }\frac{\delta \mathrm{\Sigma }}{\delta \mathrm{\Phi }_i}).$$
(26)
This relation expresses the invariance of $`\mathrm{\Sigma }`$ under the VSUSY-transformations $`\delta _\mu \mathrm{\Phi }_i:=V_{i\mu }`$. A nice feature of this approach is that, by construction, the obtained variations of the fields represent a symmetry of the theory. Yet, for a given TQFT, it may be quite tedious to carry out the calculations.
For all known models, the VSUSY- and BRST-operators satisfy the anticommutation relations
$$[s,\delta _\tau ]\mathrm{\Phi }_i=_\tau \mathrm{\Phi }_i+\text{contact terms}.$$
(27)
Here, $`_\tau =[i_\tau ,d]`$ represents the Lie derivative along the vector field $`\tau `$ and $`i_\tau `$ denotes the interior product with $`\tau `$. Since the algebra closes on space-time translations, it describes a superalgebra of Wess-Zumino type and, for brevity, we will refer to (27) as the SUSY-algebra. More precisely, this algebra closes off-shell for Witten-type models and on-shell for Schwarz-type models. The lack of off-shell closure for the latter theories can be explained by the fact that the shift-symmetry and thereby the associated ghosts are “missing”.
In the next section, we will use differential forms to present the known, as well as some new, results for our prototype models. Before doing so, we summarize some useful algebraic relations.
The graded commutation relations between the basic operators read:
$`0`$ $`=`$ $`[s,d]=[s,i_\tau ]=[s,_\tau ]`$ (28)
$`0`$ $`=`$ $`[\delta _\tau ,d]=[\delta _\tau ,i_\tau ]=[\delta _\tau ,_\tau ].`$
As usual, the Hodge dual of a Lie algebra-valued differential form $`\mathrm{\Omega }`$ will be denoted by $`\mathrm{\Omega }`$. On a $`n`$-dimensional space-time manifold $`_n`$, the star operator can be used to define a scalar product of Lie algebra-valued $`p`$-forms $`\mathrm{\Omega }_p^q`$ which, in addition, have some ghost-number $`q`$:
$$\mathrm{\Omega }_p^q,\mathrm{\Lambda }_p^r__n\mathrm{tr}(\mathrm{\Omega }_p^q\mathrm{\Lambda }_p^r).$$
This product has the graded symmetry
$`\mathrm{\Omega }_p^q,\mathrm{\Lambda }_p^r=(1)^{(p+n)(q+r)+qr}\mathrm{\Lambda }_p^r,\mathrm{\Omega }_p^q.`$ (29)
If the space-time dimension is odd, the star operation represents a mapping between forms of even and odd degree, henceforth it anticommutes with the antiderivation $`s`$.
By using the metric tensor $`(g_{\mu \nu })`$ of the space-time manifold, one can associate the $`1`$-form $`g(\tau )\tau ^\mu g_{\mu \nu }dx^\nu `$ to the vector field $`\tau =\tau ^\mu _\mu `$. (In the mathematical literature, this mapping and its inverse are known as the “musical isomorphisms”, which are usually denoted by $`\mathrm{}`$ and $`\mathrm{}`$, respectively .) The Hodge operator intertwines between the interior product $`i_\tau `$ and the exterior multiplication with $`g(\tau )`$:
$`g(\tau )\mathrm{\Omega }_p^q`$ $`=`$ $`(1)^pi_\tau \mathrm{\Omega }_p^q`$ (30)
$`g(\tau )\mathrm{\Omega }_p^q`$ $`=`$ $`(1)^{p+1}i_\tau \mathrm{\Omega }_p^q.`$
### 3.1 Witten-type models
In the literature, two different types of gauge-fixings have been considered for topological YM-theory.
* The first choice consists of a linear gauge condition for both the shift-symmetry and the ordinary gauge symmetry :
$$\mathrm{\Sigma }_1^W=__4\mathrm{tr}\left(FF\right)+s__4\mathrm{tr}\left\{\chi _2^1F^++\overline{\varphi }^2d\psi _1^1+\overline{c}dA\right\}.$$
(31)
Here, the fields $`F^+\frac{1}{2}(F+F),\chi _2^1`$ and $`s\chi _2^1B_2`$ are self-dual and the BRST-variations are defined by (11),(12) and
$$\begin{array}{ccccccc}\hfill s\chi _2^1& =& B_2\hfill & ,& \hfill sB_2& =& 0\hfill \\ \hfill s\overline{\varphi }^2& =& \eta ^1\hfill & ,& \hfill s\eta ^1& =& 0\hfill \\ \hfill s\overline{c}& =& b\hfill & ,& \hfill sb& =& 0.\hfill \end{array}$$
The action (31) is also invariant under the following VSUSY-variations :
$`\delta _\tau A`$ $`=`$ $`0,\delta _\tau \psi _1^1=_\tau A`$
$`\delta _\tau c`$ $`=`$ $`0,\delta _\tau \phi ^2=_\tau c`$ (32)
and
$$\begin{array}{ccccccc}\hfill \delta _\tau \overline{c}& =& _\tau \overline{\varphi }^2\hfill & ,& \hfill \delta _\tau b& =& _\tau \overline{c}+_\tau \eta ^1\hfill \\ \hfill \delta _\tau \chi _2^1& =& 0\hfill & ,& \hfill \delta _\tau B_2& =& _\tau \chi _2^1\hfill \\ \hfill \delta _\tau \overline{\varphi }^2& =& 0\hfill & ,& \hfill \delta _\tau \eta ^1& =& _\tau \overline{\varphi }^2.\hfill \end{array}$$
* The second choice is as follows . For the shift-symmetry, one considers a covariant gauge condition and for the ordinary gauge symmetry, one chooses either (a) a linear, (b) a covariant or (c) no gauge condition at all:
$$\mathrm{\Sigma }_{2\alpha }^W=__4\mathrm{tr}\left(FF\right)+s__4\mathrm{tr}\left\{\chi _2^1F^++\overline{\varphi }^2D\psi _1^1\right\}+s\mathrm{\Psi }_\alpha \mathrm{with}\alpha \{a,b,c\},$$
(33)
where
$`\mathrm{\Psi }_a={\displaystyle __4}\mathrm{tr}\left\{\overline{c}dA\right\},\mathrm{\Psi }_b={\displaystyle __4}\mathrm{tr}\left\{\overline{c}DA\right\},\mathrm{\Psi }_c=0.`$ (34)
The BRST-transformations are again given by (11),(12) and (3.1) (or by adding a gauge symmetry contribution $`[c,\mathrm{\Phi }_i]`$ to each variation $`s\mathrm{\Phi }_i`$ in (3.1) - such a contribution does not matter for our considerations).
For the action (33), we have found the VSUSY-variations
$`\delta _\tau A`$ $`=`$ $`0,\delta _\tau \psi _1^1=i_\tau F=i_\tau dA+[i_\tau A,A]`$
$`\delta _\tau c`$ $`=`$ $`i_\tau A,\delta _\tau \phi ^2=i_\tau \psi _1^1`$ (35)
and
$$\begin{array}{ccccccc}\hfill \delta _\tau \overline{c}& =& 0\hfill & ,& \hfill \delta _\tau b& =& _\tau \overline{c}\hfill \\ \hfill \delta _\tau \chi _2^1& =& 2\left(i_\tau D\overline{\varphi }^2\right)^+\hfill & ,& \hfill \delta _\tau B_2& =& _\tau \chi _2^1+2\left(i_\tau ([\psi _1^1Dc,\overline{\varphi }^2]D\eta ^1)\right)^+\hfill \\ \hfill \delta _\tau \overline{\varphi }^2& =& 0\hfill & ,& \hfill \delta _\tau \eta ^1& =& _\tau \overline{\varphi }^2.\hfill \end{array}$$
By contrast to the transformations (32),(3.1), these variations are not linear in the basic fields. Since $`A`$ and $`\overline{c}`$ do not transform under $`\delta _\tau `$, the gauge-fixing term $`s\mathrm{\Psi }_\alpha `$ for ordinary gauge symmetry is, taken by itself, invariant under VSUSY. In section 4.1 below, we will explain how we determined the given VSUSY-transformations and why they represent a symmetry of the action.
Both sets of VSUSY-transformations have several features in common, which can thereby be considered as characteristic for topological models of Witten-type. First, both of them fulfill the SUSY-algebra (27) off-shell. Second, both of them leave the classical field $`A`$ inert and therefore they do not act on the classical action. Thus, they only represent a non-trivial symmetry of the gauge-fixing part of the action. This should explain the fact that the VSUSY is not restrictive enough for topological YM-theory to make the model perturbatively finite, though its existence considerably improves the algebraic renormalization procedure, leading to an anomaly-free quantized theory . Finally, we remark that the classical and ghost fields, i.e. the fields which belong to the geometric part of the BRST-algebra, transform among themselves under VSUSY: none of the anti-ghosts or Lagrange multipliers involved in the gauge-fixing action appears in the variations (32) and (35).
### 3.2 Schwarz-type models
#### Chern-Simons theory
In the Landau gauge, the total Chern-Simons action is given by
$`\mathrm{\Sigma }^{CS}`$ $`=`$ $`\mathrm{\Sigma }_{inv}^{CS}+\mathrm{\Sigma }_{gf}^{CS}`$ (36)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{𝐑^3}}\mathrm{tr}(AdA+{\displaystyle \frac{2}{3}}A^3)+{\displaystyle _{𝐑^3}}\mathrm{tr}\left\{bdA+\overline{c}dDc\right\},`$
where $`\overline{c}`$ and $`b`$ are, respectively, the anti-ghost and the corresponding Lagrange multiplier, both forming a BRST-doublet: $`s\overline{c}=b,sb=0`$.
Substitution of the functional derivatives of $`\mathrm{\Sigma }^{CS}`$ with respect to $`A,b`$ and $`c`$, e.g.
$`{\displaystyle \frac{\delta \mathrm{\Sigma }^{CS}}{\delta A}}`$ $`=`$ $`Fdb[c,d\overline{c}],`$ (37)
into expression (26), leads to the (linear) VSUSY-variations
$`\delta _\tau A`$ $`=`$ $`i_\tau d\overline{c},\delta _\tau \overline{c}=\mathrm{\hspace{0.33em}0}`$
$`\delta _\tau c`$ $`=`$ $`i_\tau A,\delta _\tau b=_\tau \overline{c}.`$ (38)
The SUSY-algebra (27) now closes off-shell for $`c,\overline{c}`$ and $`b`$, but not for the classical field $`A`$:
$$[s,\delta _\tau ]A=_\tau Ai_\tau \frac{\delta \mathrm{\Sigma }^{CS}}{\delta A}.$$
(39)
From these results (and similar results for the BF model discussed below), we conclude that the VSUSY-transformations in Schwarz-type models differ substantially from those in Witten-type models: the classical fields are not invariant, but transform into anti-ghost fields, and the SUSY-algebra does not close off-shell for the classical fields. The fact that the transformations mix the classical and gauge-fixing parts of the total action renders the VSUSY highly non-trivial and constraining for the quantum theory: it is at the origin of the perturbative finiteness of these models.
We note that a (linear) VSUSY is also present if the axial gauge is chosen , but it does not exist for a covariant gauge
$$\frac{\delta \mathrm{\Sigma }}{\delta b}=dA+\alpha b(\alpha 𝐑^{}).$$
The fact that the presence of VSUSY implies a certain class of gauges is a feature that is reminiscent of the anti-BRST symmetry, whose presence has similar consequences (if considered in addition to the usual BRST-symmetry) . However, the VSUSY has a considerably richer structure which entails interesting results for the quantum theory, which is not the case for the anti-BRST symmetry.
#### BF model
The total BF action $`\mathrm{\Sigma }^{BF}=\mathrm{\Sigma }_{inv}^{BF}+\mathrm{\Sigma }_{gf}^{BF}`$ involves
$$\mathrm{\Sigma }_{gf}^{BF}=s_{_{p+2}}\mathrm{tr}\left\{\overline{c}dA+\overline{c}_{p1}dB+\mathrm{}\right\},$$
where $`\overline{c}`$ and $`\overline{c}_{p1}`$ are the anti-ghosts which fix the Landau gauge in the YM- and $`B`$-field sector, respectively, and where we only wrote out the terms which are relevant here. The derivation of VSUSY-transformations proceeds along the lines of the Chern-Simons theory, though one has to take into account the fact that the $`s`$-variations (21) of the BF model are only nilpotent by virtue of the classical equation of motion $`F=0`$ (cf. equations (22)): since we are now dealing with the complete, gauge-fixed action, these $`s`$-variations have to be extended in an appropriate way so as to relate to the complete equation of motion. This can be achieved by standard methods and the following VSUSY-transformations can be found :
$`\delta _\tau A`$ $`=`$ $`i_\tau d\overline{c}_{p1}`$
$`\delta _\tau B`$ $`=`$ $`(1)^pi_\tau d\overline{c}`$ (40)
$`\delta _\tau B_{pk}^k`$ $`=`$ $`i_\tau B_{pk+1}^{k1}\mathrm{for}k=1,\mathrm{},p.`$
For the classical fields, the SUSY-algebra (27) only closes on-shell:
$`[s,\delta _\tau ]A`$ $`=`$ $`_\tau Ai_\tau {\displaystyle \frac{\delta \mathrm{\Sigma }^{BF}}{\delta B}}`$
$`[s,\delta _\tau ]B`$ $`=`$ $`_\tau Bi_\tau {\displaystyle \frac{\delta \mathrm{\Sigma }^{BF}}{\delta A}}.`$ (41)
## 4 Derivation of vector supersymmetry transformations
In the sequel, we will repeatedly refer to the quantity
$$i_\tau \stackrel{~}{F}=i_\tau (\stackrel{~}{d}\stackrel{~}{A}+\frac{1}{2}[\stackrel{~}{A},\stackrel{~}{A}])=i_\tau \stackrel{~}{d}\stackrel{~}{A}[\stackrel{~}{A},i_\tau \stackrel{~}{A}].$$
By virtue of $`[i_\tau ,\stackrel{~}{d}]=[i_\tau ,d+s]=_\tau `$ and $`\stackrel{~}{D}=\stackrel{~}{d}+[\stackrel{~}{A},]`$, this expression takes the compact form
$$i_\tau \stackrel{~}{F}=_\tau \stackrel{~}{A}\stackrel{~}{D}i_\tau \stackrel{~}{A}.$$
(42)
We now present an alternative approach to the derivation of VSUSY-transformations for topological models. To stress the analogy with the method of horizontality conditions for the derivation of BRST-transformations, we briefly summarize the main steps which are followed for the latter derivation in the case of Chern-Simons (or ordinary YM-) theory. One assumes that the BRST-operator is nilpotent on the fields $`A`$ and $`c`$, i.e. that the graded algebra
$$[s,s]=0$$
(43)
holds for these fields. Then, one postulates the horizontality conditions involving the generalized field strength $`\stackrel{~}{F}\stackrel{~}{d}\stackrel{~}{A}+\frac{1}{2}[\stackrel{~}{A},\stackrel{~}{A}]`$, i.e. for Chern-Simons theory, one postulates $`\stackrel{~}{F}=F`$. By expanding this relation with respect to the ghost-number, one immediately obtains the BRST-transformations. Their off-shell nilpotency, i.e. the consistency of the final equations with the starting point (43), can either be checked explicitly or by resorting to the general arguments indicated in section 2. As a last step, the invariance of a given action is to be verified. (Eventually, we can also reverse the problem and look for action functionals admitting the derived BRST-variations as symmetries.)
Let us now proceed with VSUSY. First of all, we assume that the SUSY-algebra (27) is satisfied off-shell. Next, we evaluate the $`\delta _\tau `$-variation of $`\stackrel{~}{F}`$:
$`\delta _\tau \stackrel{~}{F}`$ $`=`$ $`\delta _\tau \stackrel{~}{d}\stackrel{~}{A}[\stackrel{~}{A},\delta _\tau \stackrel{~}{A}]`$ (44)
$`=`$ $`[\delta _\tau ,\stackrel{~}{d}]\stackrel{~}{A}\stackrel{~}{d}\delta _\tau \stackrel{~}{A}[\stackrel{~}{A},\delta _\tau \stackrel{~}{A}]`$
$`=`$ $`[\delta _\tau ,\stackrel{~}{d}]\stackrel{~}{A}\stackrel{~}{D}\delta _\tau \stackrel{~}{A}.`$
Substitution of the assumed off-shell algebra entails
$$\delta _\tau \stackrel{~}{F}=_\tau \stackrel{~}{A}\stackrel{~}{D}\delta _\tau \stackrel{~}{A}.$$
(45)
Comparison of the expressions (42) and (45) now motivates us to postulate either $`\mathrm{}`$-type symmetry conditions,
$`\delta _\tau \stackrel{~}{A}`$ $`=`$ $`i_\tau \stackrel{~}{A}`$ (46)
$`\delta _\tau \stackrel{~}{F}`$ $`=`$ $`i_\tau \stackrel{~}{F},`$
or $`0`$-type symmetry conditions,
$`\delta _\tau \stackrel{~}{A}`$ $`=`$ $`0`$
$`\delta _\tau \stackrel{~}{F}`$ $`=`$ $`_\tau \stackrel{~}{A}.`$ (47)
In both sets of equations, the second relation is a consequence of (or consistency condition for) the first one by virtue of equations (42),(45). The terminology $`\mathrm{}`$ versus $`0`$ simply expresses the fact that $`\delta _\tau \stackrel{~}{A}0`$ as opposed to $`\delta _\tau \stackrel{~}{A}=0`$.
For the $`B`$-field sector, we follow the same line of reasoning. From the definition $`\stackrel{~}{H}=\stackrel{~}{D}\stackrel{~}{B}`$, it follows that
$$i_\tau \stackrel{~}{H}=_\tau \stackrel{~}{B}\stackrel{~}{D}i_\tau \stackrel{~}{B}+[i_\tau \stackrel{~}{A},\stackrel{~}{B}],$$
(48)
while the assumption that $`\stackrel{~}{B}`$ satisfies the SUSY-algebra off-shell, i.e. $`[s,\delta _\tau ]\stackrel{~}{B}=_\tau \stackrel{~}{B}`$, leads to
$$\delta _\tau \stackrel{~}{H}=_\tau \stackrel{~}{B}\stackrel{~}{D}\delta _\tau \stackrel{~}{B}+[\delta _\tau \stackrel{~}{A},\stackrel{~}{B}].$$
(49)
Comparison of both relations then motivates us again to postulate either $`\mathrm{}`$-type symmetry conditions,
$`\delta _\tau \stackrel{~}{B}`$ $`=`$ $`i_\tau \stackrel{~}{B}`$ (50)
$`\delta _\tau \stackrel{~}{H}`$ $`=`$ $`i_\tau \stackrel{~}{H},`$
in conjunction with equations (46), or $`0`$-type symmetry conditions,
$`\delta _\tau \stackrel{~}{B}`$ $`=`$ $`0`$ (51)
$`\delta _\tau \stackrel{~}{H}`$ $`=`$ $`_\tau \stackrel{~}{B},`$
in conjunction with equations (47).
### 4.1 Witten-type models
Let us substitute $`\stackrel{~}{A}=A+c`$ and the horizontality condition for topological YM-theories, i.e. $`\stackrel{~}{F}=F+\psi _1^1+\phi ^2`$, into the $`0`$-type symmetry conditions (47). By decomposing with respect to the ghost-number, we immediately obtain the VSUSY-transformations (32). Similarly, from the $`\mathrm{}`$-type symmetry conditions (46), we reproduce the non-linear VSUSY-transformations (35). (Actually, this is how we found these variations!) Thus, the two representations of VSUSY defined, respectively, by the $`0`$-type and $`\mathrm{}`$-type symmetry conditions, manifest themselves in Witten-type models, the symmetry depending on the chosen gauge-fixing conditions.
Let us now come back to the VSUSY-variations (3.1) of the anti-ghosts and Lagrange multipliers. The transformations of the anti-ghosts can be found by assuming that the off-shell SUSY-algebra $`[s,\delta _\tau ]=_\tau `$ is valid, and by varying the gauge-fixing action $`\mathrm{\Sigma }_{gf}sL`$:
$$\delta _\tau \mathrm{\Sigma }_{gf}=\delta _\tau sL=_\tau Ls\delta _\tau L$$
By choosing the $`\delta _\tau `$-variations of the anti-ghosts $`(\overline{c},\chi _2^1,\overline{\varphi }^2)`$ in an appropriate way, the last term vanishes and thereby ensures the $`\delta _\tau `$-invariance of $`\mathrm{\Sigma }_{gf}`$. Finally, the transformations of the Lagrange multipliers $`(b,B_2,\eta ^1)`$ are also determined by imposing the VSUSY-algebra for them, e.g. from
$$\delta _\tau b=\delta _\tau (s\overline{c})=_\tau \overline{c}s(\delta _\tau \overline{c})$$
and the known $`\delta _\tau `$-variation of $`\overline{c}`$, one finds the one of $`b`$. Thus, it is by construction that the VSUSY-transformations (35),(3.1) represent a symmetry of the action (33). (The same arguments can be used to determine the $`\delta _\tau `$-variations (3.1) and to check the $`\delta _\tau `$-invariance of the action (31).)
We also applied our procedure to higher-dimensional TQFT’s of Witten-type, in particular to the six-dimensional model of reference . In this case as well, we could determine the corresponding VSUSY-transformations in a straightforward way , thereby confirming the usefulness of the approach to VSUSY outlined here.
### 4.2 Schwarz-type models
Before discussing the examples, we should note right away that the transformation laws that we will derive from the symmetry conditions in the present case, are to be considered with caution. In fact, our symmetry conditions are based on the assumption that the SUSY-algebra closes off-shell and this is not the case for the classical fields occurring in Schwarz-type models.
#### Chern-Simons theory
If we were to combine the horizontality conditions of the present model, i.e. $`\stackrel{~}{F}=F`$, with the $`0`$-type symmetry conditions (47), we would obtain the inadmissible result $`_\tau A=0=_\tau c`$. Henceforth, the $`0`$-type symmetry conditions (47) can only occur for Witten-type models where a shift-symmetry exists.
Thus, let us apply the $`\mathrm{}`$-type symmetry conditions (46): by decomposing the latter with respect to the ghost-number and by using $`\stackrel{~}{F}=F`$, we get
$$\delta _\tau A=0,\delta _\tau c=i_\tau A$$
(52)
and
$$\delta _\tau F=0,i_\tau F=0.$$
(53)
From the transformations (52) and the $`s`$-variations of $`A`$ and $`c`$, it follows that
$`[s,\delta _\tau ]A`$ $`=`$ $`_\tau Ai_\tau F`$ (54)
$`[s,\delta _\tau ]c`$ $`=`$ $`_\tau c,`$
where $`i_\tau F=0`$, if the classical equations of motion are used. Thus, we have obtained an on-shell algebra after having assumed the validity of an off-shell algebra as our starting point: this result is due to the truncation of the ghost-expansion $`\stackrel{~}{F}`$.
The algebra (54) can be interpreted as follows. If we only consider the classical action, the latter is invariant under the $`\delta _\tau `$-variations (52) which satisfy the SUSY-algebra on-shell. We will now try to promote this trivial symmetry of the classical action to a non-trivial symmetry of the total action (again allowing for an on-shell closure of the SUSY-algebra). To do so, we retain the non-trivial transformation law $`\delta _\tau c=i_\tau A`$ and we consider $`\delta _\tau A`$ to be unknown.
Let us evaluate the expression $`\delta _\tau sA`$ in terms of $`\delta _\tau A`$: by substituting the known expressions of $`sA`$ and $`\delta _\tau c`$, and by using $`di_\tau =_\tau di_\tau `$ as well as $`dA+\frac{1}{2}[A,A]=F`$, we obtain
$$\delta _\tau sA=_\tau Ai_\tau F+[c,\delta _\tau A].$$
(55)
By virtue of the complete equation of motion for $`A`$, as given by equation (37), the classical contact term $`i_\tau F`$ in equation (55) can be expressed in terms of the contact term $`i_\tau (\delta \mathrm{\Sigma }^{CS}/\delta A)`$. (In this way, the anti-ghosts enter our geometric framework which only involves classical and ghost fields.) Subsequent use of $`b=s\overline{c}`$ then entails
$$\delta _\tau sA=_\tau Ai_\tau \frac{\delta \mathrm{\Sigma }^{CS}}{\delta A}i_\tau (ds\overline{c})+[c,i_\tau (d\overline{c})]+[c,\delta _\tau A].$$
By adding the unknown quantity $`s\delta _\tau A`$ to both sides of this equation, we get the result
$$[s,\delta _\tau ]A=_\tau Ai_\tau \frac{\delta \mathrm{\Sigma }^{CS}}{\delta A}+s\{\delta _\tau A+i_\tau (d\overline{c})\}+[c,\delta _\tau A+i_\tau (d\overline{c})].$$
(56)
Obviously, the choice
$$\delta _\tau A=i_\tau (d\overline{c})$$
(57)
ensures the validity of the SUSY-algebra (27) and gives the known results (38),(39). The requirement of invariance of the Chern-Simons action under the determined $`\delta _\tau `$-variations fixes the transformation laws of $`\overline{c}`$ and $`b`$, again in agreement with equations (38). (We could also argue that $`\delta _\tau \overline{c}`$ has to vanish for dimensional reasons; then $`\delta _\tau b`$ follows again by imposing the SUSY-algebra on $`\overline{c}=sb`$.)
Let us summarize once more our procedure: by starting from the $`\mathrm{}`$-type symmetry conditions, we could derive the VSUSY-transformations for the Chern-Simons theory solely from the knowledge of the total action and BRST-transformations and by assuming that the SUSY-algebra is fulfilled up to contact terms.
As we have shown in the appendix, the $`\delta _\tau `$-transformations of $`A`$ and $`c`$ can also be obtained in a direct way by redoing our initial derivation (44)-(46) after having determined the contact terms in the SUSY-algebra by dimensional arguments.
#### BF model
One proceeds as for the Chern-Simons theory. If we only consider the classical action, the $`\mathrm{}`$-type symmetry conditions (46) and (50) lead to equations (52) and to the following variations of the $`B`$-fields:
$`\delta _\tau B`$ $`=`$ $`0`$
$`\delta _\tau B_{pk}^k`$ $`=`$ $`i_\tau B_{pk+1}^{k1}\mathrm{for}k=1,\mathrm{},p.`$ (58)
When extending these results to the complete gauge-fixed action, one has to take into account the fact that the SUSY-algebra is only valid on-shell and that it involves the complete equations of motion. By modifying the transformation laws $`\delta _\tau A=0`$ and $`\delta _\tau B=0`$ along the lines indicated above, one obtains the VSUSY-transformations (40) which fulfill the on-shell algebra (41).
## 5 Conclusion
From the previous considerations, we conclude that the VSUSY-transformations for Witten-type models follow straightforwardly from the $`0`$-type or $`\mathrm{}`$-type symmetry conditions (their presence depending on which gauge-fixing condition is chosen). This derivation seems to be quite efficient, in particular for higher-dimensional TQFT’s .
The VSUSY-transformations for Schwarz-type models follow from $`\delta _\tau \stackrel{~}{A}=i_\tau \stackrel{~}{A}`$ by checking the algebra and by taking into account the equations of motion of the model under consideration. An off-shell formulation for these theories can be obtained by considering the linearized Slavnov-Taylor operator which involves external sources. These sources are associated with the non-linear terms in the BRST-transformations and they also transform under the VSUSY which is now linearly broken . Since the Batalin-Vilkovisky formalism naturally incorporates sources under the disguise of anti-fields (e.g. see references for the application to topological models), it should represent a more convenient framework for discussing Schwarz-type models. This will be reported upon elsewhere .
Acknowledgments
F.G. wishes to thank R.Bertlmann for a stimulating discussion on the BF model and he expresses his gratitude to all the members of the Institut für Theoretische Physik of the Technical University of Vienna for the warm hospitality extended to him.
## Appendix A Another derivation of VSUSY-variations in Chern-Simons theory
In the following, we present a slightly different derivation of VSUSY-transformations for Schwarz-type models by using the Chern-Simons theory as an example. This approach is motivated by the results (52),(53) which entail the vanishing of $`F`$, i.e. the classical, rather than the complete equation of motion for Chern-Simons theory. This fact exhibits the inadequacy of our starting point, i.e. of the assumption that the SUSY-algebra is fulfilled off-shell. Hence, we simply review the derivation (44)-(46) after having determined the contact terms in the SUSY-algebra (27) for the commutators $`[s,\delta _\tau ]A`$ and $`[s,\delta _\tau ]c`$. These terms can be found without explicitly knowing the $`\delta _\tau `$-variations since their form is strongly constrained: for the commutator $`[s,\delta _\tau ]A`$, this term has to be a function of the functional derivatives of $`\mathrm{\Sigma }^{CS}`$ with respect to the fields of the model, and this function must be linear in $`\tau `$ and of the same dimension and ghost-number as $`A`$ (similarly for the contact term in the commutator $`[s,\delta _\tau ]c`$). From these arguments, we can deduce that the algebra can only have the form
$`[s,\delta _\tau ]A`$ $`=`$ $`_\tau A+\xi i_\tau {\displaystyle \frac{\delta \mathrm{\Sigma }^{CS}}{\delta A}}`$
$`[s,\delta _\tau ]c`$ $`=`$ $`_\tau c,`$
where $`\xi `$ is a real factor. We now collect these two equations into a single one involving $`\stackrel{~}{A}=A+c`$ and we substitute the known expression (37) for the functional derivative:
$`[s,\delta _\tau ]\stackrel{~}{A}=_\tau \stackrel{~}{A}+\xi i_\tau \{F+s(d\overline{c})[c,d\overline{c}]\}.`$ (59)
This relation represents the correct form of the SUSY-algebra for the present model. Thus, we substitute it in the $`\delta _\tau `$-variation (44) of $`\stackrel{~}{F}`$:
$$\delta _\tau \stackrel{~}{F}=_\tau \stackrel{~}{A}+\xi i_\tau F+\xi \{i_\tau s(d\overline{c})[c,i_\tau (d\overline{c})]\}\stackrel{~}{D}(\delta _\tau \stackrel{~}{A}).$$
Next, we substitute the horizontality condition $`\stackrel{~}{F}=F`$ in $`i_\tau F`$ and eliminate $`_\tau \stackrel{~}{A}`$ by means of the general relation (42): if $`\xi =1`$, the $`i_\tau F`$-term drops out from the last equation and we are left with
$$\delta _\tau \stackrel{~}{F}=\stackrel{~}{D}(i_\tau \stackrel{~}{A})+si_\tau (d\overline{c})[c,i_\tau (d\overline{c})]\stackrel{~}{D}(\delta _\tau \stackrel{~}{A}).$$
Thanks to $`\stackrel{~}{D}\stackrel{~}{d}+[\stackrel{~}{A},]=D+s+[c,]`$, this result can be rewritten as
$$\delta _\tau \stackrel{~}{F}=\stackrel{~}{D}[\delta _\tau \stackrel{~}{A}i_\tau \stackrel{~}{A}+i_\tau (d\overline{c})]+Di_\tau (d\overline{c}).$$
(60)
Henceforth, the postulated conditions (47), which did not take into account the equations of motion, should be modified according to
$`\delta _\tau \stackrel{~}{A}`$ $`=`$ $`i_\tau \stackrel{~}{A}i_\tau (d\overline{c})`$ (61)
$`\delta _\tau \stackrel{~}{F}`$ $`=`$ $`Di_\tau (d\overline{c}).`$
Expansion with respect to the ghost-number now yields the known results (38).
In summary, the key point of this model-dependent derivation was to determine the general form of contact terms in the SUSY-algebra for the considered model. |
warning/0002/math0002142.html | ar5iv | text | # Asymptotic homomorphisms into the Calkin algebra 11footnote 1This research was partially supported by RFBR (grant No 99-01-01201).
## 1 Introduction
Let $`A`$, $`B`$ be $`C^{}`$-algebras. Remind that a collection of maps
$$\phi =(\phi _t)_{t[1,\mathrm{})}:AB$$
is called an asymptotic homomorphism if for every $`aA`$ the map $`t\phi _t(a)`$ is continuous and if for any $`a,bA`$, $`\lambda 𝐂`$ one has
$$\underset{t\mathrm{}}{lim}\phi _t(ab)\phi _t(a)\phi _t(b)=0;$$
$$\underset{t\mathrm{}}{lim}\phi _t(a+\lambda b)\phi _t(a)\lambda \phi _t(b)=0;$$
$$\underset{t\mathrm{}}{lim}\phi _t(a^{})\phi _t(a)^{}=0.$$
Two asymptotic homomorphisms $`\phi ^{(0)}`$ and $`\phi ^{(1)}`$ are homotopic if there exists an asymptotic homomorphism $`\mathrm{\Phi }`$ from $`A`$ to $`BC[0,1]`$ such that its compositions with the evaluation maps at $`0`$ and at $`1`$ coincide with $`\phi ^{(0)}`$ and $`\phi ^{(1)}`$ respectively. The set of homotopy classes of asymptotic homomorphisms from $`A`$ to $`B`$ is denoted by $`[[A,B]]`$ .
Throughout this paper we always assume that $`A`$ is separable and that $`B`$ has a strictly positive element and is stable, $`BB𝒦`$, where $`𝒦`$ denotes the $`C^{}`$-algebra of compact operators. We will sometimes write $`B=B_1𝒦`$, where $`B_1=B`$, to distinguish $`B`$ from $`B𝒦`$ when necessary.
By $`Ext(A,B)`$ we denote the set of homotopy classes of extensions of $`A`$ by $`B`$. We identify extensions with homomorphisms into the Calkin algebra $`Q(B)=M(B)/B`$ by the Busby invariant . Two extensions $`f_0,f_1:AQ(B)`$ are homotopic if there exists an extension $`F:AQ(BC[0,1])`$ such that its composition with the evaluation maps at $`0`$ and at $`1`$ coincide with $`f_0`$ and $`f_1`$ respectively.
Similarly we denote by $`Ext^{as}(A,B)`$ the set of homotopy classes of asymptotic homomorphisms from $`A`$ to $`Q(B)`$. Two asymptotic homomorphisms
$$\phi ^{(i)}=(\phi _t^{(i)})_{t[1,\mathrm{})}:AQ(B),i=0,1,$$
are homotopic if there exists an asymptotic homomorphism $`\mathrm{\Phi }=(\mathrm{\Phi }_t)_{t[1,\mathrm{})}:AQ(BC[0,1])`$ such that its compositions with the evaluation maps at $`0`$ and at $`1`$ coincide with $`\phi ^{(0)}`$ and $`\phi ^{(1)}`$ respectively. Asymptotic homomorphisms into $`Q(B)`$ are sometimes called asymptotic extensions.
All these sets are equipped with a natural group structure when $`A`$ is a suspension, i.e. $`A=SD=C_0(𝐑)D`$ for some $`C^{}`$-algebra $`D`$.
As every genuine homomorphism can be viewed as an asymptotic one, so we have a natural map
$$i:Ext(A,B)Ext^{as}(A,B).$$
(1)
It is well-known that usually there is much more asymptotic homomorphisms than genuine ones, e.g. for $`A=C_0(𝐑^2)`$ all genuine homomorphisms of $`A`$ into $`𝒦`$ are homotopy trivial though the group $`[[C_0(𝐑^2),𝒦]]`$ coincides with $`K_0(C_0(𝐑^2))=𝐙`$ via the Bott isomorphism.
The main purpose of the paper is to prove epimorphity of the map (1) when $`A`$ is a suspension. This makes a contrast with the case of mappings into the compacts. As a by-product we get another description of the $`E`$-theory in terms of asymptotic extensions.
The main tool in this paper is the Connes–Higson map
$$CH:Ext(A,B)[[SA,B]],$$
which plays an important role in $`E`$-theory. Remind that for $`fExt(A,B)`$ this map is defined by $`CH(f)=(\phi _t)_{t[1,\mathrm{})}`$, where $`\phi `$ is given by
$$\phi _t:\alpha a\alpha (u_t)f^{}(a),aA,\alpha C_0(0,1).$$
Here $`f^{}:AM(B)`$ is a set-theoretic lifting for $`f:AQ(B)`$ and $`u_tB`$ is a quasicentral approximate unit for $`f^{}(A)`$. We are going to show that by fine tuning of this quasicentral approximate unit one can define also a map
$$\stackrel{~}{CH}:Ext^{as}(A,B)[[SA,B]]$$
extending $`CH`$ and completing the commutative triangle diagram
$$\begin{array}{ccc}Ext(A,B)& \stackrel{i}{}& Ext^{as}(A,B)\\ & CH& \stackrel{~}{CH}\\ & & [[SA,B]].\end{array}$$
We will show that the map $`\stackrel{~}{CH}`$ is a isomorphism when $`A`$ is a suspension.
I am grateful to K. Thomsen for his hospitality during my visit to Århus University in 1999 when the present paper was conceived.
## 2 An extension of the Connes–Higson map
A useful tool for working with asymptotic homomorphisms is the possibility of discretization suggested in . Let $`Ext_{discr}^{as}(A,B)`$ denote the set of homotopy classes of discrete asymptotic homomorphisms $`\phi =(\phi _n)_{n𝐍}:AQ(B)`$ with the additional crucial property suggested by Mishchenko: for every $`aA`$ one has
$$\underset{n\mathrm{}}{lim}\phi _{n+1}(a)\phi _n(a)=0.$$
(2)
In a similar way we define a set $`[[A,B]]_{discr}`$ as a set of homotopy classes of discrete asymptotic homomorphisms with the property (2).
###### Lemma 2.1
One has $`[[A,B]]=[[A,B]]_{discr}`$, $`Ext^{as}(A,B)=Ext_{discr}^{as}(A,B)`$.
Proof. The first equality is proved in . The second one can be proved in the same way. For an asymptotic homomorphism $`\phi =(\phi _t)_{t[1,\mathrm{})}:AQ(B)`$ one can find an infinite sequence of points $`\{t_i\}_{i𝐍}[1,\mathrm{})`$ satisfying the following properties
1. the sequence $`\{t_i\}_{i𝐍}`$ is non-decreasing and approaches infinity;
2. for every $`aA`$ one has $`\underset{i\mathrm{}}{lim}\underset{t[t_i,t_{i+1}]}{sup}\phi _t(a)\phi _{t_i}(a)=0.`$
Then $`\varphi =(\phi _{t_i})_{i𝐍}`$ is a discrete asymptotic homomorphism. It is easy to see that two homotopic asymptotic homomorphisms define homotopic asymptotic homomorphisms and that two discretizations $`\{t_i\}_{i𝐍}`$ and $`\{t_i^{}\}_{i𝐍}`$ satisfying the above properties define homotopic discrete asymptotic homomorphisms too, hence the map $`Ext^{as}(A,B)Ext_{discr}^{as}(A,B)`$ is well defined. The inverse map is given by linear interpolation of discrete asymptotic homomorphisms. $`\mathrm{}`$
Let $`(\phi _n)_{n𝐍}`$ be a discrete asymptotic homomorphism and let $`(m_n)_{n𝐍}`$ be a sequence of numbers $`m_n𝐍`$. Then we call the sequence
$$(\underset{m_1\mathrm{times}}{\underset{}{\phi _1,\mathrm{},\phi _1}},\underset{m_2\mathrm{times}}{\underset{}{\phi _2,\mathrm{},\phi _2}},\phi _3,\mathrm{})$$
a reparametrization of the sequence $`(\phi _n)_{n𝐍}`$. It is easy to see that any reparametrization does not change the homotopy class of an asymptotic homomorphism.
###### Lemma 2.2
There exists a sequence of liftings $`\phi _n^{}:AM(B)`$ for $`\phi _n`$, which is continuous uniformly in $`n`$.
Proof. It is easy to see that
$$\underset{n\mathrm{}}{lim}\underset{nk<\mathrm{}}{sup}\phi _k(a)\phi _k(b)ab$$
for any $`a,bA`$. By the Bartle–Graves selection theorem , cf. there exists a continuous selection $`s:Q(B)M(B)`$. Put $`\phi _n^{}(a)=s\phi _n(a)`$, $`aA`$. $`\mathrm{}`$
Now we are going to construct the map $`\stackrel{~}{CH}:Ext^{as}(A,B)[[SA,B]]`$. Due to Lemma 2.1 it is sufficient to define the map $`\stackrel{~}{CH}`$ as a map from $`Ext_{discr}^{as}(A,B)`$ to $`[[SA,B]]_{discr}`$.
For $`a,bA`$, $`\lambda 𝐂`$ put
$`P_n(a,b)`$ $`=`$ $`\phi _n(a)\phi _n(b)\phi _n(ab);`$
$`L_n(a,b,\lambda )`$ $`=`$ $`\phi _n(a)+\lambda \phi _n(b)\phi _n(a+\lambda b);`$
$`A_n(a)`$ $`=`$ $`\phi _n(a)^{}\phi _n(a^{})`$
and define $`P_n^{}(a,b)`$, $`L_n^{}(\lambda ,a)`$, $`A_n^{}(a)`$ in the same way but with the liftings $`\phi _n^{}`$ instead of $`\phi _n`$.
In what follows we identify $`B=B_1𝒦`$ (resp. $`M(B)`$) with the $`C^{}`$-algebra of compact (resp. adjointable) operators on the standard Hilbert $`C^{}`$-module $`B_1l_2(𝐍)=l_2(B_1)`$ and use the notion of diagonal operators in $`B`$ and $`M(B)`$ in this sense. The following Lemma shows how one has to choose a quasicentral approximate unit that makes it possible to define the map $`\stackrel{~}{CH}`$.
###### Lemma 2.3
Let $`(\phi _n)_{n𝐍}:AQ(B_1𝒦)`$ be a discrete asymptotic homomorphism. Then there exists a reparametrization of $`(\phi _n)_{n𝐍}`$ and an approximate unit $`(u_n)_{n𝐍}B_1𝒦`$ with the following properties:
1. for any $`aA`$ one has
$$\underset{n\mathrm{}}{lim}[\phi _n^{}(a),u_n]=0;$$
2. for any $`\alpha C_0(0,1)`$, for any $`a,bA`$, $`\lambda 𝐂`$ one has
$$\underset{n\mathrm{}}{lim}\alpha (u_n)P_n^{}(a,b)=\underset{n\mathrm{}}{lim}\alpha (u_n)L_n^{}(a,b,\lambda )=\underset{n\mathrm{}}{lim}\alpha (u_n)A_n^{}(a)=0;$$
3. $`lim_n\mathrm{}u_{n+1}u_n=0`$;
4. every $`u_n`$ is a diagonal operator, $`u_n=diag\{u_n^1,u_n^2,\mathrm{}\}`$, where diagonal entries $`u_n^i`$ belong to $`B_1`$ and
$$\underset{i\mathrm{}}{lim}\underset{n}{sup}u_n^{i+1}u_n^i=0.$$
Proof. Let $`\{F_n\}_{n𝐍}`$ be a generating system for $`A`$ . This means that every $`F_nA`$ is compact, $`\mathrm{}F_nF_{n+1}\mathrm{}`$ , $`_nF_n`$ is dense in $`A`$ and one has
$$F_nF_nF_{n+m(n)};F_n+\lambda F_nF_{n+m(n)},(|\lambda |1);F_n^{}F_{n+m(n)}$$
for some integer sequence $`m=(m_n)_{n𝐍}`$. Let also $`\alpha _0=e^{2\pi ix}1C_0(0,1)C_0(𝐑)`$ be a (multiplicative) generator for $`C_0(𝐑)`$.
Put
$$\epsilon _{n,k}=\underset{a,bF_k,|\lambda |1}{sup}\mathrm{max}(P_n(a,b),L_n(a,b,\lambda ),A_n(a)).$$
For every fixed $`a,b,\lambda `$ the sequences $`(P_n(a,b))`$, $`(L_n(a,b,\lambda ))`$ and $`(A_n(a))`$ vanish as $`n`$ approaches infinity, but the sequence $`(\epsilon _{n,n})_{n𝐍}`$ does not have to vanish. Nevertheless one can reparametrize the sequence $`\{F_n\}`$ by a sequence $`k=(k_n)_{n𝐍}`$, which approaches infinity slowly enough and such that $`\epsilon _{n,k(n)}`$ vanishes as $`n\mathrm{}`$. Put $`\epsilon _n=\epsilon _{n,k(n)}`$. Then
$$\underset{n\mathrm{}}{lim}\epsilon _n=0.$$
(3)
Let $`e=(e_n)_{n𝐍}B`$ be an approximate unit in $`B`$ and let $`Conv(e)`$ denote its convex hull.
By induction we can choose $`u_nConv(e)`$ in such a way that $`u_nu_{n1}`$ and that the estimates
$$[\phi _n^{}(a),u_n]<\epsilon _n;$$
(4)
and
$$\alpha _0(u_n)P_n^{}(a,b)<2\epsilon _n;\alpha _0(u_n)L_n^{}(a,b,\lambda )<2\epsilon _n;\alpha _0(u_n)A_n^{}(a)<2\epsilon _n$$
(5)
hold for any $`a,bF_n`$ and any $`|\lambda |1`$.
It is easy to see that the conditions (4-5) together with Lemma 2.2 ensure the first two items of Lemma 2.3.
The above choice of $`(u_n)_{n𝐍}`$ does not yet ensure the condition $`lim_n\mathrm{}u_{n+1}u_n=0`$. To make it hold we have to renumber the sequence $`(\phi _n)_{n𝐍}`$. At first divide every segment $`[u_n,u_{n+1}]`$ into $`n`$ equal segments $`[u_{n_i},u_{n_{i+1}}]`$, $`i=1,\mathrm{},n`$. Then as $`0u_i1`$ for all $`i`$, so we get $`u_{n_{i+1}}u_{n_i}\frac{1}{n}`$. Finally we have to change the sequences $`(\phi _1,\phi _2,\phi _3,\mathrm{})`$ and $`(u_1,u_2,u_3,\mathrm{})`$ by the sequence $`(\phi _1,\phi _2,\phi _2,\phi _3,\mathrm{})`$, where each $`\phi _n`$ is repeated $`n`$ times, and by the sequence $`(u_{1_1},u_{2_1},u_{2_2},u_{3_1},u_{3_2},u_{3_3},u_{4_1},\mathrm{})`$ respectively.
To prove the last item of Lemma 2.3 remind that an approximate unit $`(e_n)_{n𝐍}B=B_1𝒦`$ can be chosen to be diagonal, $`e_n=b_nϵ_n`$, where $`(b_n)_{n𝐍}B_1`$ and $`(ϵ_n)_{n𝐍}𝒦`$ are approximate units in $`B_1`$ and in $`𝒦`$ respectively, so the quasicentral approximate unit $`(u_n)_{n𝐍}Conv(e)`$ can be made diagonal as well, with diagonal entries from $`B_1`$.
Let $`T`$ be the right shift on the standard Hilbert $`C^{}`$-module $`l_2(B_1)=B_1l_2(𝐍)`$, $`TM(𝒦)M(B_1𝒦)`$. We can join $`S`$ to the sets $`\phi _n^{}(F_n)`$ in (4) when constructing the sequence $`(u_n)`$. Then the sequence $`[T,u_n]B_1𝒦`$ would vanish as $`n`$ approaches infinity. Hence
$$\underset{n\mathrm{}}{lim}\underset{i}{sup}u_n^{i+1}u_n^i=0$$
(6)
and the operators
$$diag\{u_n^2u_n^1,u_n^3u_n^2,u_n^4u_n^3,\mathrm{}\}$$
are compact, so $`lim_i\mathrm{}u_n^{i+1}u_n^i=0`$. Take $`\epsilon >0`$. By (6) there exists some $`N`$ such that for all $`n>N`$ one has $`sup_iu_n^{i+1}u_n^i<\epsilon `$. Now consider the finite number of compact operators $`diag\{u_n^2u_n^1,u_n^3u_n^2,u_n^4u_n^3,\mathrm{}\}`$, $`1nN`$. Due to their compactness there exists some $`I`$ such that for $`i>I`$ one has $`u_n^{i+1}u_n^i<\epsilon `$ for $`1nN`$. Therefore for $`i>I`$ we have $`u_n^{i+1}u_n^i<\epsilon `$ for every $`n`$, i.e. $`sup_nu_n^{i+1}u_n^i<\epsilon `$. $`\mathrm{}`$
Put now
$$\stackrel{~}{CH}(\phi )_n(\alpha a)=\alpha (u_n)\phi _n^{}(a),\alpha C_0(0,1),aA,$$
where $`(u_n)_{n𝐍}`$ satisfies the conditions of Lemma 2.3. Items $`i)iii)`$ of Lemma 2.3 ensure that $`(\stackrel{~}{CH}(\phi )_n)_{n𝐍}`$ is a discrete asymptotic homomorphism from $`SA`$ to $`B`$. If $`(u_n)_{n𝐍}`$ and $`(v_n)_{n𝐍}`$ are two quasicentral approximate unities satisfying Lemma 2.3 then the linear homotopy $`(tu_n+(1t)v_n)_{n𝐍}`$ provides that the maps $`\stackrel{~}{CH}`$ defined using these approximate unities are homotopic. Finally, if $`\phi `$ and $`\psi `$ represent the same homotopy class in $`Ext_{discr}^{as}(A,B)`$ then $`\stackrel{~}{CH}(\phi )`$ and $`\stackrel{~}{CH}(\psi )`$ are homotopic. If all $`\phi _n`$ are constant, $`\phi _n=f:AQ(B)`$ with $`f`$ being a genuine homomorphism, then obviously $`CH(f)=\stackrel{~}{CH}(\phi )`$, so we have
###### Lemma 2.4
The map $`\stackrel{~}{CH}:Ext^{as}(A,B)[[SA,B]]`$ is well defined and the diagram
$$\begin{array}{ccc}Ext(A,B)& \stackrel{i}{}& Ext^{as}(A,B)\\ & CH& \stackrel{~}{CH}\\ & & [[SA,B]].\end{array}$$
is commutative. $`\mathrm{}`$
## 3 An inverse for $`\stackrel{~}{CH}`$
Let $`\alpha _0=e^{2\pi ix}1`$ be a generator for $`C_0(0,1)`$ and let $`T`$ be the right shift on the Hilbert space $`l_2(𝐍)`$. By $`q:M(B)Q(B)`$ we denote the quotient map. Define a homomorphism
$$g:C_0(0,1)Q(𝒦)\mathrm{by}g(\alpha _0)=q(T)1.$$
Remind that $`B`$ is stable and denote by $`\iota :Q(B)𝒦Q(B)`$ the standard inclusion. Put
$$j=\iota (gid_B):SBQ(𝒦)BQ(B).$$
The homomorphism $`j`$ obviously induces a map
$$j_{}:[[A,SB]]Ext^{as}(A,B).$$
Let $`S:[[A,B]][[SA,SB]]`$ denote the suspension map. Then the composition $`M=j_{}S`$ gives a map
$$M:[[A,B]]Ext^{as}(SA,B).$$
Let
$$\beta =(\beta _n)_{n𝐍}:C_0(𝐑^2)𝒦$$
(7)
be a discrete asymptotic homomorphism representing a generator of $`[[C_0(𝐑^2),𝒦]]`$. For a discrete asymptotic extension $`\phi =(\phi _n)_{n𝐍}:AQ(B)`$ consider its tensor product by $`\beta `$
$$\phi \beta =(\phi _n\beta _n)_{n𝐍}:S^2AQ(B)𝒦$$
and denote its composition with the standard inclusion $`Q(B)𝒦Q(B)`$ by
$$Bott_1=\iota (\phi \beta ):Ext^{as}(A,B)Ext^{as}(S^2A,B).$$
In a similar way define a map
$$Bott_2:[[A,B]][[S^2A,B]].$$
###### Theorem 3.1
One has
$$M\stackrel{~}{CH}=Bott_1;\stackrel{~}{CH}M=Bott_2.$$
Proof. We start with $`M\stackrel{~}{CH}=Bott_1`$. Let $`H`$ be the standard Hilbert $`C^{}`$-module over $`B`$, $`H=Bl_2(𝐍)`$. Put $`=_{n𝐍}H_n`$, where every $`H_n`$ is a copy of $`H`$. We identify the $`C^{}`$-algebra of compact (resp. adjointable) operators on both $`H`$ and $``$ with $`B`$ (resp. $`M(B)`$). Instead of writing formulas in $`Q(B)`$ we will write them in $`M(B)`$ and understand them modulo compacts.
Let $`\phi =(\phi _n)_{n𝐍}:AQ(B)`$ represent an element $`[\phi ]Ext_{discr}^{as}(A,B)`$ and let $`\phi _n^{}:AM(B)`$ be liftings for $`\phi _n`$ as in Lemma 2.2.
If $`a_n:H_nH_n`$ is a sequence of operators then we write $`(a_1a_2a_3\mathrm{})`$ for their direct sum acting on $`=_{n𝐍}H_n`$. In what follows we use a shortcut
$$\alpha (u_n)\phi _n^{}(a)=a_n.$$
Let $`T`$ be the right shift on $``$, $`T:H_nH_{n+1}`$.
Remind that $`\alpha _0`$ is a generator for $`C_0(0,1)`$ and that it is sufficient to define asymptotic homomorphisms on the elements of the form $`\alpha a\alpha _0S^2A`$.
The composition map $`M\stackrel{~}{CH}(\phi )_n:S^2AQ(B)`$ acts by
$$M\stackrel{~}{CH}(\phi )_n(\alpha a\alpha _0)=\left(a_na_na_n\mathrm{}\right)(T1)$$
modulo compacts on $``$.
Let
$$v_n=\left(𝐯_n^1𝐯_n^2𝐯_n^3\mathrm{}\right)M(B_1𝒦)$$
and
$$\lambda _n=\left(\lambda _n^1\lambda _n^2\lambda _n^3\mathrm{}\right)M(B_1𝒦)$$
be a direct sum of diagonal operators $`𝐯_n^i=diag\{v_n^i,v_n^i,v_n^i,\mathrm{}\}`$, $`v_n^iB_1`$, and a direct sum of scalar operators, $`\lambda _n^i𝐑`$, ($`𝐯_n^i`$ and $`\lambda _n^i`$ act on $`H_i`$). Let the numbers $`\lambda _n^i`$ satisfy the properties
1. $`\lambda _n^1=0`$ and $`lim_i\mathrm{}\lambda _n^i=1`$ for every $`n`$;
2. $`lim_n\mathrm{}sup_i|\lambda _n^{i+1}\lambda _n^i|=0`$;
3. $`lim_n\mathrm{}sup_i|\lambda _{n+1}^i\lambda _n^i|=0`$.
We assume that the elements $`v_n^i`$ are selfadjoint and satisfy the following properties:
1. $`lim_n\mathrm{}sup_iv_n^{i+1}v_n^i=0`$;
2. $`lim_n\mathrm{}sup_iv_{n+1}^iv_n^i=0`$;
3. there exists a set of scalars $`\lambda _n^i𝐑`$, $`n,i𝐍`$, satisfying the conditions $`i)iii)`$ above and such that
$$\underset{n\mathrm{}}{lim}(v_n^i\lambda _n^i)b=0$$
(8)
for every $`i`$ and for every $`bB_1`$.
Let $`p`$ be a projection onto the first coordinate in $`H=B_1l_2(𝐍)`$ and let $`P=(ppp\mathrm{})`$. Then $`P\lambda _n^i=diag\{\lambda _n^i,0,0,\mathrm{}\}`$ and the map $`\beta _n`$ (7) can be written as
$$\beta _n(\alpha \alpha _0)=P\alpha (\lambda _n)(T1)1_{B_1}𝒦M(B_1𝒦)$$
and the map $`Bott_1(\phi ):S^2AQ(B)`$ can be written in the form
$$(Bott_1(\phi ))_n(\alpha a\alpha _0)=\left(\alpha (\lambda _n^1)\phi _n^{}(a)\alpha (\lambda _n^2)\phi _n^{}(a)\alpha (\lambda _n^3)\phi _n^{}(a)\mathrm{}\right)(T1).$$
Consider also the path of asymptotic homomorphisms $`(\mathrm{\Phi }_n(t))_{n𝐍}`$, $`t[0,1]`$, given by the formula
$$\mathrm{\Phi }_n(t)(\alpha a\alpha _0)=\left(\alpha (𝐯_n^1(t))\phi _n^{}(a)\alpha (𝐯_n^2(t))\phi _n^{}(a)\alpha (𝐯_n^3(t))\phi _n^{}(a)\mathrm{}\right)(T1),$$
where for every $`i`$ $`v_n^i(t)`$ is a piecewise linear path connecting $`v_n^i(\frac{1}{k})=v_{n1+k}^i`$, $`k𝐍`$, and $`v_n^i(0)=\lambda _n^i`$. In view of (8) it is easy to check that $`\mathrm{\Phi }_n(t)`$ is a homotopy connecting the asymptotic homomorphisms $`(Bott_1(\phi ))_n`$ and $`\mathrm{\Phi }_n=\mathrm{\Phi }_n(0)`$.
One of the obvious choices for $`v_n^i`$ is to put $`(v_n^i)^{i𝐍}=(0,\frac{1}{n},\frac{2}{n},\mathrm{},1,1,\mathrm{})`$. But for our purposes it is better to use another choice. We take $`v_n^i=u_i^n`$ for all $`n`$ and $`i`$. Lemma 2.3 ensures that the properties $`i)iii)`$ are satisfied.
Now we have to connect the asymptotic homomorphisms $`(Bott_1(\phi )_n)_{n𝐍}`$ and $`(M\stackrel{~}{CH}(\phi )_n)_{n𝐍}`$ by a homotopy in the class of asymptotic homomorphisms. In fact we are going to do more and to connect each of these asymptotic homomorphisms with a genuine homomorphism $`f:S^2AQ(B)`$ defined modulo compacts by
$$f(\alpha a\alpha _0)=\left(a_1a_2a_3\mathrm{}\right)(T1),$$
$`\alpha C_0(0,1)`$, $`aA`$. Lemma 2.3 ensures that $`f`$ is indeed a homomorphism.
At first we connect $`f`$ with $`(M\stackrel{~}{CH}(\phi ))_n`$ by a path $`F_n(t)`$, $`t[0,1]`$. Let $`F_n(1)=(M\stackrel{~}{CH}(\phi ))_n`$. Denote $`\alpha (u_n)\phi _n^{}(a)`$ by $`a_n`$ and put (modulo compacts)
$$F_n\left(\frac{1}{2}\right)(\alpha a\alpha _0)=\left(\underset{n\mathrm{times}}{\underset{}{a_n\mathrm{}a_n}}a_{n+1}a_{n+1}a_{n+1}\mathrm{}\right)(T1),$$
$$F_n\left(\frac{1}{3}\right)(\alpha a\alpha _0)=\left(\underset{n\mathrm{times}}{\underset{}{a_n\mathrm{}a_n}}a_{n+1}a_{n+2}a_{n+2}\mathrm{}\right)(T1),$$
etc. Finally put
$$F_n(0)(\alpha a\alpha _0)=\left(\underset{n\mathrm{times}}{\underset{}{a_n\mathrm{}a_n}}a_{n+1}a_{n+2}a_{n+3}\mathrm{}\right)(T1)$$
and connect $`F_n(1)`$, $`F_n(\frac{1}{2})`$, $`F_n(\frac{1}{3})`$, $`\mathrm{}`$ and $`F_n(0)`$ by a piecewise linear path $`F_n(t)`$, $`t[0,1]`$.
It is easy to see that for every $`t>0`$ the sequence $`(F_n(t))_{n𝐍}`$ is an asymptotic homomorphism. And as the maps $`F_n(0)`$ and $`f`$ differ by compacts, so they coincide as homomorphisms into $`Q(B)`$. Continuity in $`t`$ is also easy to check. So the asymptotic homomorphism $`(M\stackrel{~}{CH}(\phi )_n)_{n𝐍}`$ is homotopic to the homomorphism $`f`$.
Now we are going to construct a homotopy $`F_n^{}(t)`$, $`t[0,1]`$, which connects $`f`$ with $`(Bott_1(\phi )_n)_{n𝐍}`$.
For each $`k𝐍`$ consider the following sequence $`(𝐮_n^k)_{n𝐍}`$ of diagonal operators, each of which acts on the corresponding copy of $`H=H_n`$ in their direct sum $``$:
$$𝐮_n^k=diag\{u_n^1,u_n^2,\mathrm{},u_n^{k1},u_n^k,u_n^k,u_n^k\mathrm{}\}.$$
Put $`u_n(\frac{1}{k})=𝐮_n^k`$, $`u_n(0)=u_n`$ and connect them by a piecewise linear path $`u_n(t)`$, $`t[0,1]`$. Then we get a strictly continuous path of operators $`u_n(t)`$, which gives a homotopy
$$F_n^{}(t)(\alpha a\alpha _0)=\left(\underset{n\mathrm{times}}{\underset{}{a_{1,n}(t)\mathrm{}a_{n,n}(t)}}a_{n+1,n+1}(t)a_{n+2,n+2}(t)\mathrm{}\right)(T1),$$
where $`a_{i,n}(t)=\alpha (u_i(t))\phi _n^{}(a)`$.
As
$$\mathrm{\Phi }_n(\alpha a\alpha _0)=\left(a_{1,n}(1)a_{2,n}(1)a_{3,n}(1)\mathrm{}\right)(T1),$$
so for every $`\alpha a`$ one has
$$\underset{n\mathrm{}}{lim}F_n^{}(1)(\alpha a\alpha _0)\mathrm{\Phi }_n(\alpha a\alpha _0)=0,$$
hence the asymptotic homomorphisms $`F_n^{}(1)`$ and $`\mathrm{\Phi }_n`$ are equivalent. But we already know that $`\mathrm{\Phi }_n`$ is homotopic to $`Bott_1(\phi )_n`$. On the other hand, it is easy to see that $`F_n^{}(0)`$ coincides with $`f`$ modulo compacts, so we can finally conclude that $`M\stackrel{~}{CH}=Bott_1`$ up to homotopy.
The second identity of Theorem 3.1 is much simpler to prove. For $`\psi =(\psi _n)_{n𝐍}:AB`$ we have (modulo compacts)
$$M(\psi )_n(\alpha _0a)=\left(\psi _1(a)\psi _2(a)\psi _3(a)\mathrm{}\right)(T1),aA.$$
But as every $`\psi _n(a)B_1𝒦`$, i.e. is compact, so when choosing a quasicentral approximate unit $`\{w_n\}_{n𝐍}`$ for the map $`(M(\psi ))_{n𝐍}`$ we can define it by
$$w_n=\left(𝐰_1^{(n)}𝐰_2^{(n)}𝐰_3^{(n)}\mathrm{}\right),$$
where each $`𝐰_i^{(n)}`$ is a finite rank diagonal operator of the form
$$𝐰_i^{(n)}=diag\{\underset{m_n\mathrm{times}}{\underset{}{\lambda _ib_n,\mathrm{},\lambda _ib_n}},0,0,\mathrm{}\}$$
for some numbers $`(m_n)_{n𝐍}`$, where $`(b_n)_{n𝐍}`$ is a quasicentral approximate unit for $`B_1`$ and the scalars $`\lambda _i`$ are defined by
$$\lambda _i=\{\begin{array}{cc}\frac{ni+1}{n}\hfill & \mathrm{for}i<n,\hfill \\ \lambda _i=0\hfill & \mathrm{for}in.\hfill \end{array}$$
But after such a choice of $`w_n`$ the map $`(\stackrel{~}{CH}M)(\psi )_n`$ differs from the map $`Bott_2(\psi )_n`$ only by the presence of $`b_n`$, hence these maps are equivalent. $`\mathrm{}`$
## 4 Case of $`A`$ being a suspension
As there exists a homomorphism $`C_0(𝐑)C_0(𝐑^3)M_2`$ that induces an isomorphism in $`K`$-theory and an asymptotic homomorphism $`C_0(𝐑^3)C_0(𝐑)𝒦`$, which are inverse to each other, so the groups $`[[SA,B]]`$ and $`[[S^3A,B]]`$ are naturally isomorphic to each other and the same is true for the groups $`Ext^{as}(SA,B)`$ and $`Ext^{as}(S^3A,B)`$. Hence we obtain
###### Corollary 4.1
If $`A`$ is a suspension then the map $`\stackrel{~}{CH}:Ext^{as}(A,B)[[SA,B]]`$ is an isomorphism. $`\mathrm{}`$
It was proved in that if $`A`$ is a suspension then the map
$$CH:Ext(A,B)[[SA,B]]$$
is surjective and the group $`[[SA,B]]`$ is contained in $`Ext(A,B)`$ as a direct summand. Hence from Corollary 4.1 we immediately obtain
###### Corollary 4.2
Let $`A`$ be a suspension. Then
1. the map $`i:Ext(A,B)Ext^{as}(A,B)`$ is surjective, hence every asymptotic extension $`\phi =(\phi _t)_{t[1,\mathrm{})}:AQ(B)`$ is homotopic to a genuine extension;
2. the group $`Ext^{as}(A,B)`$ is contained in $`Ext(A,B)`$ as a direct summand. $`\mathrm{}`$
###### Problem 4.3
Is Corollary 4.2 true when $`A`$ is not a suspension ?
For $`C^{}`$-algebras $`A`$ and $`B`$ consider the set of all extensions $`f:AQ(B)`$ that are homotopy trivial as asymptotic homomorphisms and denote by $`Ext^{ph}(A,B)`$ the set of homotopy classes of such homomorphisms. As usual this set becomes a group when $`A`$ is a suspension. We call the elements of $`Ext^{ph}(A,B)`$ phantom extensions because they constitute the part in $`Ext(A,B)`$ which vanishes under the suspension map $`S:Ext(A,B)Ext(SA,SB)`$, cf. .
###### Corollary 4.4
If $`A`$ is a second suspension then there is a natural decomposition
$$Ext(A,B)=Ext^{ph}(A,B)Ext^{as}(A,B).$$
$`\mathrm{}`$
###### Remark 4.5
If $`A`$ is both a nuclear $`C^{}`$-algebra and a suspension then the groups $`Ext(A,B)`$ and $`[[A,B]]`$ coincide , therefore there is a one-to-one correspondence between homotopy classes of genuine and asymptotic homomorphisms into the Calkin algebras $`Q(B)`$ and one has $`Ext^{ph}(A,B)=0`$.
###### Problem 4.6
Does there exist a separable $`C^{}`$-algebra $`A`$ such that the $`Ext^{ph}(A,B)`$ is non-zero for some $`B`$ ?
Our definition of homotopy in $`Ext^{as}(A,B)`$ is weaker than the homotopy of asymptotic homomorphisms in $`[[A,Q(B)]]`$, so there is a surjective map
$$p:[[A,Q(B)]]Ext^{as}(A,B).$$
(9)
It would be interesting to compare the composition $`\stackrel{~}{CH}p`$ with the would-be boundary map $`[[A,Q(B)]][[SA,B]]`$ which would exist if the exact sequences of the $`E`$-theory could be generalized to the non-separable short exact sequence $`BM(B)Q(B)`$.
###### Problem 4.7
Is the map $`p`$ (9) injective ?
V. M. Manuilov
Dept. of Mech. and Math.,
Moscow State University,
Moscow, 119899, RUSSIA
e-mail: manuilov@mech.math.msu.su |
warning/0002/cond-mat0002430.html | ar5iv | text | # Bouncing–ball tunneling in quantum dots
## Abstract
We show that tunneling through quantum dots can be completely dominated by states quantized on stable bouncing–ball orbits. The fingerprints of bouncing–ball tunneling are sequences of Coulomb blockade peaks with strongly correlated peak height and asymmetric peak line shape. Our results are in agreement with the striking correlations of peak height and transmission phase found in recent interference experiments with quantum dots.
Transport experiments with semiconductor quantum dots yield important information on the energy levels and the wave functions of confined many–electron systems . In the Coulomb blockade regime current can flow only if two different charge states of the quantum dot are tuned to have the same energy. This produces large nearly equally spaced conductance peaks as a function of the dot potential (tuned via electrostatic gates). At low temperatures, the magnitude of each conductance peak directly measures the magnitude of a single resonant wave function near the contacts to the leads. Some years ago, a statistical theory for the peak height distribution was developed based on the assumption that the wave functions are completely uncorrelated and described by random matrix theory. Recent experiments found good agreement with the predicted distributions. At the same time, one of these experiments reported large correlation between the heights of adjacent peaks. Much stronger correlations of both the peak height and the phase of the transmission amplitude through a quantum dot were observed in novel interference experiments . The correlations in these experiments comprised all Coulomb blockade peaks found in a gate voltage scan, yielding sequences of more than 10 correlated peaks.
Theoretical work demonstrated that nonzero temperature can partially account for but not fully explain the peak correlations observed in Ref. . More recently , it was shown that additional correlations can arise from short periodic orbits. Such orbits lead to deviations from random matrix predictions even for dots whose classical dynamics is completely chaotic. The predicted peak modulations are of the type observed in Ref. , but can not account for the much stronger correlations in the interference experiments . Despite of considerable theoretical efforts the origin of these correlations has not been understood. In particular, the most striking feature of the experiments , that the correlations comprise all observed Coulomb peaks, has not been explained. We address this feature below.
We show that tunneling through quantum dots can be completely dominated by wave functions quantized on stable bouncing–ball orbits. Bouncing–ball tunneling (BBT) dominates the transport if (i) the leads are attached opposite to each other and (ii) if they are sufficiently wide providing near normal injection of electrons into the dot. Under the conditions (i) and (ii) the dot conductance exhibits sequences of strongly correlated Coulomb peaks. All peaks within a sequence are dominated by the same bouncing–ball mode. The unique fingerprint of BBT is a characteristic asymmetry of the peak line shape for the peaks in the tails of the sequences. This asymmetry is caused by the breaking of the particle–hole symmetry as the bouncing–ball state moves away from the Fermi energy. We find the asymmetry in the line shapes measured in Ref. , and suggest that the inspection of line shapes be used as a diagnostic for BBT in other experiments on quantum dots.
Our starting point is the transmission probability $`T_{\mathrm{ring}}`$ through an Aharonov–Bohm ring with a quantum dot embedded in one arm. We assume that the dot is in the Coulomb blockade regime. The tunnel barriers around the quantum dot suppress multiple reflections of electrons across the ring, and $`T_{\mathrm{ring}}`$ is given by
$`T_{\mathrm{ring}}={\displaystyle 𝑑E\left(\frac{f}{E}\right)\left|t_0+t(E)\mathrm{exp}[2\pi i\mathrm{\Phi }/\mathrm{\Phi }_0]\right|^2}.`$ (1)
Here, $`\mathrm{\Phi }`$ is the magnetic flux through the ring, $`\mathrm{\Phi }_0`$ is the flux quantum, and $`t_0`$, $`t(E)`$ are the transmission amplitudes through the free arm and the arm with the quantum dot, respectively ($`t_0`$ is assumed to be independent of injection energy). The ring is connected to two reservoirs labeled up (U) and down (D) occupied according to the Fermi distribution $`f`$. Below, we investigate the phase $`\varphi _{\mathrm{QD}}\mathrm{arg}𝑑E(f/E)t(E)`$ and the transmission coefficient $`T_{\mathrm{QD}}𝑑E(f/E)|t(E)|^2`$. Both can be measured in quantum dot interference experiments .
The transmission amplitude $`t(E)`$ is related to the retarded Green function $`G_{pq}`$ of the quantum dot,
$`t(E)={\displaystyle \underset{pq}{}}V_p^UV_q^DG_{pq}(E),`$ (2)
where $`V_p^l`$ is the matrix element for tunneling between level p in the dot and the reservoir $`l=U,D`$. For a weakly coupled dot the width of the level $`p`$ is given by $`\mathrm{\Gamma }_p=\mathrm{\Gamma }_p^U+\mathrm{\Gamma }_p^D`$, where $`\mathrm{\Gamma }_p^l=|V_p^l|^2`$ is the partial width for decay into the reservoir $`l`$. The matrix element is expressed as
$`V_p^l=\left({\displaystyle \frac{\mathrm{}^2}{m^{}}}\right){\displaystyle 𝑑s\psi _l(s,z)_z\psi _p^{}(s,z)}|_{z=0},`$ (3)
where the integration is performed along the edge between the potential barrier and the quantum dot ($`_z`$ denotes the derivative normal to the barrier). The wave function $`\psi _p`$ corresponds to Dirichlet boundary conditions in the dot, while the barrier tunneling is fully included in the lead wave function $`\psi _l`$. The transverse potential in the tunneling region can be taken quadratic yielding $`\psi _lc_l\mathrm{exp}[(ss_l)^2/2a_{\mathrm{eff}}^2]`$, where $`s`$ is the transverse coordinate, $`s_l`$ the center of the constriction and $`a_{\mathrm{eff}}`$ its effective width. We can restrict ourselves to the lowest transverse mode since higher modes are suppressed by the barrier penetration factor (included in $`c_l`$). The Gaussian form of the lead wave function is convenient but not crucial for the results presented below.
We obtain the dot wave functions $`\psi _p`$ for an effective potential which accounts for the dot confining potential and the interactions in the dot in a mean–field sense. We use a billiard approximation and model the dot by a hard–wall potential at the boundary, parameterized in polar coordinates by $`R(\varphi )=R_0[1+ϵ\mathrm{cos}(2\varphi )]`$. The parameter $`ϵ`$ measures the quadrupolar deformation out of circular shape. We use nonzero $`ϵ`$ to mimic the dots used in . As in the experiments, we assume that the leads are attached opposite to each other at the boundary points closest to the origin (the points with $`\varphi =\pm \pi /2`$). In contrast to most studies of quantum dots which assume the dynamics in the dot to be either regular or fully chaotic, our boundary parameterization yields mixed classical motion. This is illustrated for $`ϵ=0.2`$ in Fig. 2, showing chaotic motion in most of phase space except for two large islands associated with stable bouncing–ball motion between the contacts to the leads. We note that billiards of quadrupolar or similar shape have recently been used in studies of optical microresonators .
In Fig. 1 we show the coupling strength $`g_p\mathrm{\Gamma }_p^U\mathrm{\Gamma }_p^D/(\mathrm{\Gamma }_p^U+\mathrm{\Gamma }_p^D)`$ over a sequence of 100 billiard states for two different values of $`a_{\mathrm{e}\mathrm{f}\mathrm{f}}`$. Note that $`g_p`$ is proportional to the conductance peak height $`G_p=(e^2/h)(\pi /2kT)g_p`$ measured in low temperature Coulomb blockade experiments . We calculated $`g_p`$ from Eq. (3) using quantum dot wave functions obtained numerically with the boundary integral method . For narrow leads the peak heights fluctuate with the level index $`p`$, occasionally one observes systematic peak modulations as reported in Ref. . In striking contrast, the results for wide leads show a few isolated large peaks, separated by $`1525`$ levels with much smaller peak height (not visible on the scale of Fig. 1(b)). All large peaks are associated with states quantized on stable bouncing–ball orbits. This is illustrated in the inset for the state $`p=337`$. The height of the small peaks not resolved in Fig. 1(b) is typically two or more orders of magnitude smaller than the maximum peak height. Such tiny peaks are difficult to resolve in standard Coulomb blockade experiments. The interference experiments face additional noise from the transmission through the free arm of the ring. For wide leads, interference experiments are therefore only sensitive to the strongly coupled bouncing–ball modes.
The crucial role of the transverse barrier width $`a_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ can be understood from the following argument: Confinement in a lead of width $`a_{\mathrm{eff}}`$ yields the transverse momentum spread $`\mathrm{}/a_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ for electrons injected in the quantum dot. Wide leads therefore result in near normal injection and provide exceptionally strong coupling to the bouncing–ball states. To quantify this argument we express the partial width $`\mathrm{\Gamma }_p^l=|V_p^l|^2`$ in terms of the Husimi function $`H_p`$. Recalling the Gaussian form of $`\psi _l`$ and using the arc length $`s`$ as integration variable in Eq. (3), we find
$`\mathrm{\Gamma }_p^l=\left({\displaystyle \frac{|c_l|\mathrm{}^2}{m^{}}}\right)^2H_p(s^l,0),`$ (4)
where $`H_p`$ is calculated for the normal derivative of the wave function $`\psi _p`$. The Husimi function is evaluated at the phase space point ($`s=s^l,p_s=0)`$ reflecting the position $`s_l`$ of the lead and zero transverse momentum. We note that the Husimi function can be written as an integral over the Wigner function, smoothed with Gaussians both in arc length and transverse momentum. Their width is given by $`a_{\mathrm{eff}}/\sqrt{2}`$ and $`\mathrm{}/\sqrt{2}a_{\mathrm{eff}}`$, respectively. We express these widths in terms of the uncertainties $`\mathrm{\Delta }\varphi `$, $`\mathrm{\Delta }\mathrm{sin}\chi `$ in polar angle $`\varphi `$ and injection angle $`\chi `$, respectively ($`\chi `$ is the angle between the momentum of the injected electrons and the normal to the boundary). Using the relation $`p_s=p\mathrm{sin}\chi `$, we find $`\mathrm{\Delta }\varphi 0.7(ka_{\mathrm{eff}})/(kR_0)`$ and $`\mathrm{\Delta }\mathrm{sin}\chi 0.7/(ka_{\mathrm{eff}})`$, where $`R_0`$ is the average radius of the dot.
The phase space portrait of the dot with deformation $`ϵ=0.2`$ is shown in Fig. 2 using $`(\varphi ,\mathrm{sin}\chi )`$ coordinates. Superimposed is the Husimi function of the bouncing–ball mode $`p=337`$ calculated for $`ka_{\mathrm{e}\mathrm{f}\mathrm{f}}=5.0`$ and $`kR_040`$, corresponding to approximately $`400`$ electrons on the dot. The Husimi function attains maximum value at the phase space points $`(\varphi ,\mathrm{sin}\chi )=\pm (\pi /2,0)`$ representing normal injection from the leads.
In order to obtain the Husimi function and hence $`\mathrm{\Gamma }_p`$ analytically, we must solve for the wave function $`\psi _p`$. For states localized on a chain of stable islands the calculation can be carried out using the semiclassical theory developed in Refs. . The result can be expressed in terms of the stable periodic orbit at the center of the island,
$`H_p(s,p_s)=`$ $`{\displaystyle \underset{\mu }{}}{\displaystyle \frac{A_\mu }{\mathrm{\Delta }s_\mu \mathrm{\Delta }p_\mu }}\mathrm{exp}\left[\left({\displaystyle \frac{ss_\mu }{\mathrm{\Delta }s_\mu }}\right)^2\right]`$ (6)
$`\times \mathrm{exp}\left[\left({\displaystyle \frac{p_sp_\mu }{\mathrm{\Delta }p_\mu }}\right)^2\right],`$
where $`s_\mu `$ is the arc length and $`p_\mu `$ the transverse momentum at the bounce points $`\mu `$ of the periodic orbit. Moreover, $`\mathrm{\Delta }s_\mu =[a_{\mathrm{eff}}^2+l_\mu ^2]^{1/2}`$ and $`\mathrm{\Delta }p_\mu =[(\mathrm{}/a_{\mathrm{eff}})^2+(\mathrm{}/l_\mu )^2]^{1/2}`$ where $`l_\mu =\sqrt{2\mathrm{}|m_{12}^\mu |}|4\mathrm{Tr}^2[M_\mu ]|^{1/4}`$ is determined by the monodromy matrix $`M_\mu (m_{ij}^\mu )`$. The amplitude $`A_\mu `$ depends only weakly on $`a_{\mathrm{eff}}`$. The result (6) allows us to express the width of island states in terms of the phase space distance of the underlying periodic orbit from the lead injection points. No comparable analytical result is known for states quantized in the chaotic sea. However, we find numerically that for wide leads the Husimi function of such states has exponentially small support near the injection points (see the inset of Fig. 2).
We now turn to the calculation of the transmission coefficient $`T_{\mathrm{QD}}`$ and the phase $`\varphi _{\mathrm{QD}}`$. We assume wide leads and temperature $`kT<\mathrm{\Delta }`$. The calculation for the weak coupling limit $`\mathrm{\Gamma }_pkT,\mathrm{\Delta }`$ is straightforward and can be done using the Green function derived in Ref. ; the results will be presented elsewhere . Here, we address the regime of a more open dot as realized in the experiment . This regime is characterized by $`\mathrm{\Gamma }_{\overline{p}}\mathrm{\Delta }`$ where we identify $`\overline{p}`$ with the bouncing–ball state closest to the Fermi energy. All other states $`p\overline{p}`$ in the vicinity of $`E_F`$ have a width $`\mathrm{\Gamma }_p\mathrm{\Delta }`$. The Green function is found using the equations–of–motion method . It is diagonal up to small off–diagonal corrections $`𝒪(\sqrt{\mathrm{\Gamma }_p\mathrm{\Gamma }_{\overline{p}}}/\mathrm{\Delta })`$, and given by
$`G_{\overline{p}\overline{p}}(E)`$ $`=`$ $`{\displaystyle \underset{N=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{P_N}{E[E_{\overline{p}}eV_g+UN]+i\mathrm{\Gamma }_{\overline{p}}/2}}.`$ (7)
Here $`N`$ counts the total number of electrons in all dot levels except for the level $`\overline{p}`$, $`P_N`$ is the respective occupation probability, $`E_{\overline{p}}`$ is the single–particle energy of state $`\overline{p}`$, $`V_g`$ is the gate voltage, and $`U=e^2/C`$ is the charging energy. The transmission through the states $`p\overline{p}`$ is negligible. Equation (7) is the generalization of the weak–coupling result to the case of a single strongly conducting quantum state. The probability $`P_N=Z^1\mathrm{exp}[\mathrm{\Omega }(N)/kT]`$ with $`Z=_N\mathrm{exp}[\mathrm{\Omega }(N)/kT]`$ is related to the thermodynamic potential $`\mathrm{\Omega }(N)`$ of the dot. To evaluate $`P_N`$ we replace $`\mathrm{\Omega }_N`$ by $`\mathrm{\Omega }_N^0+[E_{\overline{p}}eV_g+UNE_F]n_{\overline{p}}_N`$, where $`\mathrm{\Omega }_N^0`$ is calculated for the dot with level $`\overline{p}`$ excluded from the spectrum, and
$`n_{\overline{p}}_N={\displaystyle \frac{1}{\pi }}{\displaystyle 𝑑E\mathrm{Im}\frac{f(E)}{E[E_{\overline{p}}eV_g+UN]+i\mathrm{\Gamma }_{\overline{p}}/2}}`$ (8)
is the canonical occupation probability of level $`\overline{p}`$.
In Fig. 3 we show the transmission coefficient $`T_{\mathrm{QD}}`$ and the phase $`\varphi _{\mathrm{QD}}`$ vs. gate voltage $`V_g`$ calculated for $`kT=0.2\mathrm{\Delta }`$, $`\mathrm{\Gamma }_{\overline{p}}=1.5\mathrm{\Delta }`$ and $`U=12\mathrm{\Delta }`$. All peaks shown result from transmission through the level $`p=337`$. The peaks have comparable height and similar phase in qualitative agreement with the experiments . The central peak has a Lorentzian shape of width $`\mathrm{\Gamma }_{\overline{p}}`$. Note that the transmission peaks develop a peculiar asymmetry as the conducting level moves away from the Fermi energy: Each peak to the left and to the right of the central peak falls off more rapidly on the side facing the central peak. This pattern extends over the whole sequence and becomes more pronounced for the peaks in the tails. This asymmetry results from the breaking of the particle–hole symmetry in the transmission through the bouncing–ball state and is a unique fingerprint of BBT. Inspection of the data of the experiment reveals the same asymmetry , providing strong evidence that the peak correlations in this experiment are due to BBT. We note that the number of correlated peaks in Fig. 3 is less than observed in experiment. The origin for this is addressed below. A detailed discussion of this issue and of the phase lapse between the peaks will be given in a future publication .
We finally address the universality of the results presented above. The stability of bouncing–ball motion does not rely on the effective dot potential considered here. Similar stability if found for dots of other shapes as well as for smooth confining potentials . We assumed that the dot shape does not change upon the addition of electrons. Some robustness of the shape has indeed been demonstrated in a recent experiment . In general, however, variation of the gate voltage is expected to change the electrostatic potential and modify the shape of the dot. It was shown that shape deformations can enhance peak correlations by “pinning” conducting levels close to the Fermi energy. To study the effect of deformations, we varied the shape by a parameter linear in gate voltage $`V_g`$. Over a range of $`V_g`$ we observed the predicted enhancement of correlations and sequences of more than 10 correlated peaks. At the same time, the bouncing–ball modes changed very little with $`V_g`$ as substantial modification of the shortest stable modes typically requires a large change in potential.
In conclusion, we have demonstrated a new mechanism for transport through quantum dots in the Coulomb blockade regime. Current flows by tunneling through bouncing–ball modes. The fingerprints of bouncing–ball tunneling are sequences of correlated conductance peaks with asymmetric peak line shape. We find the peak asymmetries in recent Coulomb blockade interference experiments, providing strong evidence that the peak correlations in these experiments are caused by bouncing–ball tunneling.
We acknowledge helpful conversations with F. Haake and H. A. Weidenmüller. The work was supported by the Deutsche Forschungsgemeinschaft, DGAPA–UNAM, and CONACyT–Mexico. |
warning/0002/astro-ph0002023.html | ar5iv | text | # The relativistic iron line profile in the Seyfert 1 Galaxy IC4329a
## 1 INTRODUCTION
Accretion of material onto a black hole is known to produce hard X–ray ($`E2`$ keV) emission: satellite observations over the last 30 years have conclusively shown this to be the case for both stellar mass black holes in our own Galaxy, and the supermassive black holes which power the Active Galactic Nuclei (AGN). However, it is still the case that the mechanism by which the gravitational potential energy is converted into high energy radiation is not understood. Standard models of an accretion disk (Shakura & Sunyaev 1973, hereafter SS) produce copious UV and even soft X–ray radiation, but are completely unable to explain the observed higher energy X–ray emission which extends from the disk spectrum out to 200 keV. Clearly something other than the standard model is required, and the two currently favored candidates are either magnetic reconnection above an accretion disk (e.g. Galeev, Rosner & Vaiana 1979; Haardt, Maraschi & Ghisellini 1994; di Matteo 1998), or that some part of the accretion flow is not given by the standard disk configuration.
The recent rediscovery of another stable solution of the accretion flow equations lent plausibility to this second possibility. The standard SS accretion disk model derives the accretion flow structure in the limit where the gravitational energy released is radiated locally in an optically thick, geometrically thin disk. This contrasts with the new accretion models, where below some critical mass accretion rate, $`\dot{m}\dot{m}_{\mathrm{crit}}`$, the material is not dense enough to thermalise and locally radiate all the gravitational potential energy which is released. Instead the energy can be carried along with the flow (advected), eventually disappearing into the black hole. These solutions give an X–ray hot, optically thin, quasi–spherical flow (Narayan & Yi 1995; Esin, McClintock & Narayan 1997).
Clearly it would be nice to know which (if any!) of these models for the hard X–ray emission is correct. In its most general form the problem comes down to understanding the geometry of the cool, optically thick accretion flow. If this extends down to the last stable orbit around the black hole then it is unlikely that the advective flow can exist (e.g. Janiuk, Życki & Czerny 1999). Conversely, if the optically thick disk truncates before this point then a composite model with an outer disk and inner, hot, advective flow (ADAF) may be favored.
The accretion flow can be tracked via X–ray spectroscopy. Wherever hard X–rays illuminate optically thick material then this gives rise to a Compton reflected continuum and associated iron K$`\alpha `$ fluorescence line (e.g. Lightman & White 1988; George & Fabian 1991; Matt, Perola & Piro 1991). The amplitude of the line and reflected continuum is dependent on the solid angle subtended by the disk to the X–ray source (as well as on elemental abundances, inclination and ionisation). The two models outlined above can then be distinguished by the amount of reflection and line, since an untruncated disk should subtend a rather larger solid angle than a truncated one.
A survey of AGN showed that Seyfert 1 spectra are consistent with a power law X–ray spectrum illuminating an optically thick, (nearly) neutral disk, which subtends a solid angle of $`2\pi `$ (Pounds et al. 1990; Nandra & Pounds 1994). This then seems to favor the magnetic reconnection picture. However, this contrasts with the situation in the Galactic Black Hole Candidates (GBHC). These are also thought to be powered by accretion via a disk onto a black hole, and, in their low/hard state, show broad band spectra which are rather similar to those from AGN, but have an apparently rather smaller amount of reflection (e.g. Gierliński et al. 1997; Done & Życki 1999). One potential drawback of this comparison is that AGN inhabit a more complex environment than the GBHC. Unification schemes for Seyfert galaxies propose that there is a molecular torus which enshrouds the nucleus. The molecular torus can also contribute to the reflected spectrum, perhaps distorting our view of the very innermost regions in AGN.
These two potential sites for the reflected component can be distinguished spectrally: any features from the accretion disk should be strongly smeared by the combination of special and general relativistic effects expected from the high orbital velocities in the vicinity of a black hole (Fabian et al. 1989), whereas the molecular torus is at much larger distances so its reflected features should be narrow. The seminal ASCA observation of the AGN MCG–6–30–15 showed that the line is so broad that it requires that the accretion disk extends down to at least the last stable orbit in a Schwarzschild metric, with no narrow component from the molecular torus (Tanaka et al. 1995; Iwasawa et al. 1996), and that the relativistically smeared material subtended a solid angle of $`2\pi `$ with respect to the X–ray source. This very clear cut result then seems to rule out the advective flows, at least in their simplest form.
However, the GBHC again show a rather different picture: they have a line which is detectably broad, but not so broad as might be expected for a disk extending all the way down to 3 Schwarzschild radii ($`R_{\mathrm{Schw}}=2GM/c^2`$; Życki, Done, & Smith 1997, 1998, 1999; Done & Życki 1999). This is consistent with the truncated disk geometry, and so perhaps with the advective flow models. Only in the soft/high state do the GBHC seem to show the extreme relativistic smearing and large amount of reflection seen in the MCG–6–30–15 spectra (Życki et al. 1998, Gierliński et al. 1999).
Is there a real difference in accretion geometry between GBHC and AGN, pointing to a difference in radiation mechanism ? This seems unlikely, since both classes are ultimately accreting black holes. Are subtle ionisation effects masking the derived disk parameters in the GBHC (Ross, Fabian & Young 1999). Or is MCG–6–30–15 unusual among AGN (and GBHC) in having such a relativistic disk? Perhaps MCG–6–30–15 is an AGN in a state which corresponds to the soft state GBHC?
One factor supporting the latter is a recent study by Zdziarski, Lubiński, & Smith (1999) which showed that there is a correlation between the intrinsic spectral slope, $`\mathrm{\Gamma }`$, and the solid angle subtended by the reflecting material, $`\mathrm{\Omega }/2\pi `$. In their plots, MCG–6–30–15 is the AGN with the steepest spectrum, and highest amount of reflection. This correlation also holds for individual objects (such as NGC 5548 for AGN and Nova Muscae for the GBHC), where the intrinsic spectrum hardens as the amount of reflection decreases. This suggests that there is a universal physical mechanism and/or geometry for both classes, with perhaps a single parameter determining the state of a given source through a feedback between the geometry and physical conditions in the X–ray emitting region. Such a feedback could be provided by e.g. soft photons from the thermalized fraction of the hard X–rays intercepted by the reprocessing medium, and the parameter could be the inner disk radius. Perhaps for MCG–6–30–15 (and other soft AGN and the soft state GBHC) the cool accretion disk extends down to the innermost stable orbit around the black hole with the X–rays being powered by magnetic reconnection above the disk, while for harder spectra AGN (and the low state GBHC) the inner disk recedes outwards, being replaced by an X–ray hot (advective ?) flow. As the disk recedes it subtends a smaller solid angle, so there is less reflection (and less relativistic smearing), but there are also fewer seed photons from the disk (both from intrinsic emission and by reprocessing) for Compton scattering into the intrinsic power law, giving a harder intrinsic power law (Życki et al. 1999; Zdziarski et al. 1999).
If this is true, then this clearly predicts that the Fe K$`\alpha `$ line in AGN is not always as broad as in the extreme case of MCG-6-30-15 (Tanaka et al. 1995). Previous ASCA studies on a sample of objects (Nandra et al. 1997) hint towards such a possibility, since a whole range of geometries were inferred for various objects.
Here we test this idea using ASCA, XTE and OSSE data from IC 4329a, the brightest ’typical’ Seyfert 1 in the X–ray band (cf. Madejski et al. 1995). It lies towards the middle of the $`\mathrm{\Gamma }\mathrm{\Omega }/2\pi `$ plot of Zdziarski et al. (1999), and (consequently) has a spectrum very close to that of the mean Seyfert 1 spectrum compiled by Zdziarski et al. (1995). IC 4329a may then be used as a template for Seyfert galaxies as a class, unlike MCG–6–30–15 which has a rather steep X–ray spectral index.
## 2 DATA REDUCTION
### 2.1 ASCA
The 1997 campaign for IC 4329a included four ASCA pointings, on August 7, 10, 12, and 16, each nominally providing 20 ks of data. The resulting data were extracted using the standard screening procedures, yielding total exposures of 61 ks for SIS0 and SIS1 (using the BRIGHT2 mode), and 78 ks for GIS2 and GIS3. The source data were extracted from circular regions with radii of 3 arc min for the SISs and 4 arc min for the GISs, while background was taken from a source-free regions of the same images. The source showed a clear variability between the four pointings, with GIS2 counting rates of $`1.77\pm 0.009`$, $`1.27\pm 0.008`$, $`1.53\pm 0.009`$, and $`1.78\pm 0.009`$, matching the variability seen in the simultaneous RXTE data (see below and Fig. 1). The PHA data were then grouped so as to have more than 20 counts per bin.
### 2.2 RXTE: PCA
IC4329a was observed a total of 66 times over a period of 58 days with RXTE. The PCA standard 2 data were extracted from all layers of detectors 0,1 and 2, using standard selection criteria (Earth elevation angle $`>10^{}`$, offset between the source and the satellite pointing direction $`<0.02^{}`$, electron rates in each detector $`<0.1`$, excluding data taken within 30 minutes of the SAA passage). This gave a total of 73 ks of good data. The background for these data were then modeled using the ’L7’ model (see Zhang et al. 1993; Madejski et al. 1999; and Jahoda et al. 2000, in preparation), and the corresponding response matrix for each observation was generated using version 3.0 of the channel to energy conversion file. We use data from 3–20 keV, since lower energies are affected by response matrix uncertainties and the background may not be well modeled above 20 keV. A 1% systematic error is applied to data in all PHA channels.
### 2.3 RXTE: HEXTE
The HEXTE instrument onboard RXTE (cf. Rothschild et al. 1998) consists of two scintillator modules, sensitive in the hard X–ray band. The background is measured via chopping the detectors to off-axis, source free locations every 16 seconds. The HEXTE data were extracted using standard data reduction procedures, which, after appropriate dead-time correction, yielded a total net exposure of $`24`$ ks. The source was clearly detected in each pointing over the range of 15 – 100 keV range, at a level consistent with the PCA, but the statistical errors in each pointing were too large to study variability. The resultant data were binned such that channels 17–28 (16.9-28.3 keV) were grouped by a factor 3, 29–48 (28.3–47.8 keV) by a factor 4 and 49–120 (47.8–123.4 keV) by a factor 6. The effective area of both HEXTE clusters was scaled by 0.7, the current best normalization to the Crab spectrum.
### 2.4 OSSE
OSSE pointed at IC 4329a during the CGRO viewing period 625, over the epoch of 1997 August 8 to 18, with the total exposure of 567 ks. The data were reduced in a standard manner (see Johnson et al. 1997 and references therein), resulting in a net counting rate of $`0.18\pm 0.08`$ counts s<sup>-1</sup> over the 50 – 500 keV range. The resulting data were binned such that channels 9–18 were grouped by a factor 5, while 19–48 were grouped by a factor 10.
### 2.5 Lightcurves
Figure 1 shows the 2–10 keV light curve obtained from the RXTE PCA data. The source is clearly variable on timescales of a few days, with an r.m.s variation of 13%. There is no significant short timescale variability within each observations. These typically have 3$`\sigma `$ upper limits of 0.06 to the fractional r.m.s variability of the 1000–2000 second lightcurve binned on 16 seconds. The horizontal lines on Figure 1 mark the times at which ASCA and OSSE data were taken.
## 3 SPECTRAL FITTING
The data were fitted using XSPEC v10.0, with errors quoted as 90% confidence intervals ($`\mathrm{\Delta }\chi ^2=2.7`$). We use elemental abundances and cross-sections of Morrison & McCammon (1983).
### 3.1 ASCA SIS and GIS
Data from SIS0 and SIS1 are showing increasing divergence with time from the GIS2 and GIS3 (and from each other) at low energies. The reasons for this are not yet well understood (Weaver & Gelbord 1999, in preparation). A current pragmatic approach is to allow excess absorption in the SIS detectors to account for this effect to zeroth order.
Another source of low energy complexity is the partially ionized absorber, first seen in the ROSAT PSPC spectrum (Madejski et al. 1995). A previous ASCA observation has shown that in IC4329a this complex absorption is better modeled by two edges (corresponding to H and He–like Oxygen at rest energies of 0.739 and 0.871 keV, respectively) rather than a full ionized absorber code (Cappi et al. 1996; Reynolds 1997 see also the discussion in George et al. 1998). We use this description here, but we caution that the determination of this absorption in our data will be somewhat dependent on the way the low energy calibration problems are treated.
We first use a phenomenological model for the ASCA data, consisting of an underlying power law and its reflected continuum from a neutral disk (pexrav: Magdziarz & Zdziarski 1995) inclined at $`30^{}`$, together with a separate Gaussian iron line, with the two edge description for the warm absorber. This model gives a good fit to the data ($`\chi _\nu ^2=2533/2392`$), for an intrinsic power law spectrum of $`\mathrm{\Gamma }=1.85\pm 0.03`$, and reflector solid angle (for an inclination of $`30^{}`$) of $`\mathrm{\Omega }/2\pi =0.48_{0.31}^{+0.34}`$. The associated iron K$`\alpha `$ fluorescence line is at a (rest frame) energy of $`6.37\pm 0.06`$ keV, with equivalent width of $`180\pm 50`$ eV and intrinsic width of $`\sigma =0.39\pm 0.10`$ keV (hereafter all intrinsic line widths are the gaussian $`\sigma `$).
The iron line physical width is similar to that seen in a previous observation of this AGN (Mushotzky et al. 1995; Cappi et al. 1996; Reynolds 1997), but is much smaller than that seen from the archetypal relativistically smeared line in MCG–6–30–15 (Tanaka et al. 1995). One way to show this is to fit the spectrum in the 2.5–10 keV band (where the effects of the ionised absorption is much less) with a simple power law and its expected reflected continuum (with solar abundances, and fixed solid angle $`\mathrm{\Omega }/2\pi =1`$ at $`30^{}`$ inclination). The fit excludes the 5–7 keV iron line range. Figure 2a shows the resulting shape of the line residuals from IC4329a, while figure 2b shows those from MCG–6–30–15 for comparison.
Replacing the Gaussian line with a diskline relativistic model (fixing the disk inclination at $`30^{}`$ and line emissivity at $`r^\beta `$ with $`\beta =3`$) gives that the smearing implies an inner disk radius of $`48_{20}^{+33}R_\mathrm{g}`$ ($`\chi _\nu ^2=2540/2392`$: a slightly worse fit than the simple broad Gaussian line). This is significantly larger than the last stable orbit at $`6R_g`$ ($`R_\mathrm{g}=GM/c^2`$, so the last stable orbit at 3 Schwarzchild radii is at $`6R_\mathrm{g}`$). The derived radius is dependent on the form of the emissivity, but fits to two separate observations of MCG–6–30–15 require $`\beta =3.4_{0.8}^{+1.3}`$ and $`\beta =4.4_{1.1}^{+3.0}`$, respectively (Nandra et al 1997), showing that the line emission is strongly weighted towards the innermost radii. We might also expect this on theoretical arguments. The energy emitted per unit area of a disk goes as $`r^3`$, so this should give the time averaged emissivity from a magnetic corona, as well as being a good approximation to the illumination expected from a central spherical source (see appendix A in Życki et al 1999). Thus we fix $`\beta =3`$ in all our fits, which gives $`\mathrm{\Delta }\chi ^2>20`$ for an inner radius equal to the last stable orbit at any inclination.
The reflection description used above has several drawbacks. Firstly it allows the line to vary in intensity (and energy) without reference to the reflected continuum. Secondly, the relativistic smearing is only applied to the iron line, and not to the reflected continuum also. We replace these components with the reflection model described by Życki et al. (1999), rel-repr, which calculates the self–consistent line associated with the reflected continuum, and then applies the relativistic smearing (including gravitational light bending) to this total reprocessed spectrum. The ionisation state is a free parameter, and the models are calculated for iron abundance between $`12\times `$ solar. The reprocessor in AGN is generally assumed to be neutral but an accretion disk temperature of $`10`$ eV can give a thermal population dominated by ions with ionization potential $`1030\times kT`$ (Rybicki & Lightman 1979) i.e. between Fe V – FeX, which in our model corresponds to $`0.3<\xi <20`$. Photo–ionisation by the X–ray source can increase these estimates considerably.
We replace the pexrav and diskline components with our combined model for the reprocessed component. We allow the ionisation, iron abundance, and inner disk radius to be fit parameters, for a fixed inclination of $`30^{}`$, and for a fixed line emissivity of $`r^3`$. This again gave a disk inner radius of $`45_{18}^{+32}R_\mathrm{g}`$ ($`\chi _\nu ^2=2540/2392`$).
### 3.2 ASCA–XTE Simultaneous Data
We extracted RXTE data which were taken exactly simultaneously with the ASCA observation (datasets 4,5,6,11,12,13,16,16–01,18,19 and 20), giving a total PCA exposure of 13 ks. The corresponding HEXTE data over this short time interval have very low signal to noise so do not add any appreciable constraints. Figure 3a shows the residuals resulting from an absorbed power law fit to the PCA data, showing the classic signature of Compton reflection and its associated iron fluorescence line ($`\chi _\nu ^2=162.3/42`$). Figure 3b shows residuals including a broad gaussian line in the fit. Clearly there are still systematic residuals, with a decrement at the expected energy of the iron edge and a rise to higher energies ($`\chi _\nu ^2=50.5/39`$). Including a reflected continuum component (the pexrav model in xspec, assuming solar abundances and inclination of $`30^{}`$) gives $`\chi _\nu ^2=11.7/38`$, showing that the reflected continuum component is significantly detected independently of the iron line.
Since the data are simultaneous, we can fix the absorption (warm and cold) to that seen in the GIS for the same model fit to the ASCA data. This gives an excellent fit to the data with $`\chi _\nu ^2=12.6/39`$ (rather too good in fact, showing that the statistics are dominated by the 1 per cent systematic error) for an intrinsic power law spectrum of $`\mathrm{\Gamma }=1.94\pm 0.04`$, and reflector solid angle (for an inclination of $`30^{}`$) of $`\mathrm{\Omega }/2\pi =0.49_{0.14}^{+0.17}`$. The associated iron K$`\alpha `$ fluorescence line is at a (rest frame) energy of $`6.33\pm 0.13`$ keV, with equivalent width of $`210\pm 45`$ eV and intrinsic width (gaussian $`\sigma `$) of $`0.49_{0.18}^{+0.19}`$ keV. The PCA 2–10 keV flux is 25 per cent higher than that from ASCA due to absolute flux calibration problems in both instruments: ASCA gives a Crab 2–10 keV flux of $`1.8\times 10^8`$ ergs s<sup>-1</sup> (Makishima et al., 1996), while RXTE gives $`2.4\times 10^8`$ ergs s<sup>-1</sup> (see PCA Crab spectrum at http://lheawww.gsfc.nasa.gov/users/keith/pcarmf.html) and the original Crab calibration is $`2.15\times 10^8`$ ergs s<sup>-1</sup> (Toor & Seward 1974). A more serious discrepancy is in the spectral index, which is $`\mathrm{\Delta }\alpha 0.1`$ steeper than that seen in ASCA. It is known that the current RXTE PCA calibration gives results for the Crab which are roughly $`\mathrm{\Delta }\alpha 0.1`$ steeper than the index assumed for the calibration of other instruments (K. Jahoda, private communication). However, reassuringly, the relative amount of reflection equivalent width, energy and intrinsic width of the line are consistent with those from the ASCA data, though of course the absolute normalisations are different due to the calibration discrepancies. Tieing the absorption across the two datasets, together with the relative reflection parameters (the reflector solid angle, its inner radius and ionisation state, and the iron line energy, width and equivalent width) gives $`\chi _\nu ^2=2548/2435`$, not significantly different from the $`\chi _\nu ^2=2546/2432`$ obtained from the separate fits. Thus in what follows we tie the relative reflection (and absorber) parameters across the two datasets, and let only the power law spectral index and normalization be free.
We replace the pexrav and broad Gaussian line with our rel-repr model. Table 1 shows the results of a joint fit to the ASCA and RXTE PCA data for inclination angles of $`30`$, $`60`$ and $`72^{}`$ and for iron abundances of $`1`$ and $`2\times `$ solar. The derived inner disk radius is highly correlated with the assumed inclination. High inclination angles give stronger Doppler effects and so a broader line. The inner disk radius then has to be larger to match the observed line width. However, none of the fits allow the disk to extend down to the innermost stable orbit, irrespective of inclination. This also means that the data are not very sensitive to the inclination ($`\mathrm{\Delta }\chi ^24`$, i.e. marginally significant preference for higher inclinations). It is only the broadest components from the very innermost disk which significantly change the skewness of the line profile (as opposed to its width) as function of inclination.
The data are sensitive to the iron abundance. They significantly prefer supersolar abundances ($`\mathrm{\Delta }\chi ^27`$), as would be expected from measurements of radial abundance gradients in spiral galaxies (e.g. Henry & Worthey 1999). Another way to see this is the phenomenological fits give a line equivalent with of $`180`$ eV, as expected for a solar abundance slab subtending a solid angle of $`2\pi `$ (e.g. George & Fabian 1991), yet the solid angle of the reflected continuum in these fits is approximately half of this. Figure 4 shows the best fit joint RXTE and ASCA data fit to a disk model with $`2\times `$ solar abundance, inclined at $`30^{}`$.
Since IC4329a is a Seyfert 1 there could be a contribution to the line/reflected continuum from a molecular torus as well as from the accretion disk (Ghisellini, Haardt & Matt 1994; Krolik, Madau & Życki 1994). Table 2 shows the results obtained including a neutral, unsmeared reflector, with assumed mean inclination of $`60^{}`$. Only the parameters of the two reflectors are included, since the continuum is similar to that derived before. This gives a significantly better fit to the data, generally with $`\mathrm{\Delta }\chi ^29`$ for the addition of 1 extra free parameter (the amount of non–relativistic reflection).
The double reprocessing model now gives a fit which is as good as or even better than those from the phenomenological (i.e. unphysical) broad Gaussian line/pexrav model. The observed line is fairly broad but also fairly symmetric, contrary to the predictions of a line from an accretion disk which must also be skewed if it is broad. Adding the second neutral, non–smeared reprocessor gives a narrow core to the line, while the line from the relativistic reprocessor then fills in a broad (and skewed) line wing. The presence of the narrow component means that the relativistic effects have to be more marked than in the single reflector fits in order to make the total line as broad as before. Thus the derived inner disk radius is always rather smaller than before but still never consistent with the 3 Schwarzschild radii. The smaller inner radius means that the lower inclination reflected spectra are significantly gravitationally redshifted, so requiring the reflector to be somewhat ionized to compensate for this.
The above discussion assumed that the torus was Compton thick, i.e. with $`N_H10^{24}`$ cm<sup>-2</sup>. However, it can still produce substantial line emission, without the accompanying strong reflected component if the torus column is $`10^{2324}`$ cm<sup>-2</sup>. Replacing the unsmeared, neutral reflected component by a narrow 6.4 keV line gives a similar series of fits as those shown in Table 2. In particular, the fits never allow the amount of smearing to be as large as expected from the innermost stable orbit of an accretion disk, and they show the same preference for twice solar iron abundance and inclination angles $`>30^{}`$.
The series of fits above show that the best physical description of the data is with two reprocessed components, one which is relativistically smeared and possibly ionized from the accretion disk, and one which is neutral and unsmeared (and possibly consisting of just line rather than line and reflected continuum) from the molecular torus. A similar composite line is seen in the Seyfert MCG–5–23-16 (Weaver et al. 1997; see also Weaver & Reynolds 1998). The data prefer models with twice solar abundances, and inclinations of $`>30^{}`$, and these solutions have the advantage that the derived ionisation of reflector is generally rather low, consistent with that expected (see section 2). However, there is a further constraint on the inclination, since the extra reprocessor cannot be along our line of sight to the nucleus. IC4329a is classified optically as a Seyfert 1 and the X–ray spectrum is not heavily absorbed. For a disk inclined at $`72^{}`$ then for our line of sight not to intercept the molecular torus severely restricts its scale height, and so the possible solid angle it can subtend. Thus while the range of double reflector fits given in Table 2 are statistically indistinguishable, these consistency arguments lead us to favor viewing angles to the accretion disk of $`45^{}`$ and supersolar abundances.
All these models assumed that the ionisation state of the disk was constant with radius. A more physical picture might be one where the ionisation varies as a function of radius (e.g. Matt et al. 1993). In such models, the inner disk might be so ionized that it produces no spectral features. The observed reflected spectrum would then arise from further out in the disk, and so not contain the highly smeared components. This might provide an alternative explanation to a truncated disk as to why the relativistic smearing observed is not compatible with a disk extending down to the last stable orbit. We test this by dividing the disk into 10 radial zones, with ionisation $`\xi (r)r^4`$ as described in Done & Życki (1999). With the inner radius fixed at $`6R_\mathrm{g}`$, and allowing for a narrow component from the molecular torus we are never able to obtain fits within $`\mathrm{\Delta }\chi ^220`$ of those in Table 2.
### 3.3 RXTE PCA Variability
There is clear intensity variability during the RXTE campaign, and also spectral variability in the sense that the spectrum becomes softer as the source brightens. This could be due to either intrinsic change in the power law spectral index, or a change in the relative contribution of the reflected spectrum (or both).
We can attempt to disentangle the power law from the reflected spectrum by fitting these components to the individual spectra from each orbit (including a 1 per cent systematic uncertainty), fixing the absorption at $`3\times 10^{21}`$ cm<sup>-2</sup>. The reflected spectrum is assumed to have twice solar iron abundance, be inclined at $`60^{}`$ and have fixed negligible ionisation ($`\xi =0.01`$, see table 1). Thus the free parameters are the power law index and normalization, the solid angle subtended by the relativistic reflector and its inner radius. Figure 5a shows these derived parameters as a function of time. Clearly the inner radius cannot be constrained, but the index and the solid angle of the reflector show some possible trend. Fitting these with a constant gives $`\chi _\nu ^2=36.0/58`$ and $`13.4/58`$ respectively i.e. they are statistically consistent with a constant value. Figure 5b shows these plotted against the 2–10 keV flux, and a linear regression (taking errors in both $`x`$ and $`y`$ into account: Press et al. 1992) shows that the power law index is significantly correlated with flux, since it gives $`\chi _\nu ^2=25.7/57`$. This corresponds to an F value of $`10.3/(25.7/57)=22.8`$, significant at $`99.9`$ per cent. Even using just the difference in $`\chi ^2`$ between the two fits gives $`F=10.3`$, significant at $`99.5`$ per cent, so the correlation is clearly present. IC4329a then becomes only the second Seyfert 1 where there is clear intrinsic spectral variability (the other is NGC 5548: Magdziarz et al. 1998), where underlying continuum changes can be unambiguously disentangled from changes in the reflected spectrum.
We use the same procedure to look for variability in the amount of reflection as a function of flux. The linear regression gives $`\chi _\nu ^2=10.6/57`$ where the fractional amount of reflection relative to the power law decreases as the flux increases. This is significant at $`>99.9`$ per cent confidence on an F test, but is only 90 per cent significant using just the change in $`\chi ^2`$.
To illustrate these points we co–add spectra when the source was at its lowest and highest intensity level (see Figure 1), and fit these spectra with a power law and single relativistically smeared reflector (with reflector parameters fixed as above, and with the disk inner radius fixed at $`60R_\mathrm{g}`$). The spectral index changes from $`1.92\pm 0.04`$ to $`2.00\pm 0.02`$, for the low and high state, respectively. Figure 6 shows the unfolded spectra, with the model extrapolated out to 100 keV. The thick and thin lines show the model components for the high and low state data, respectively. The underlying continuum is brighter at all energies $`100`$ keV in the high state despite it being steeper, with an integrated 0.01–300 keV flux of $`6.2`$ to $`10.7\times 10^{10}`$ ergs cm<sup>-2</sup> s<sup>-1</sup> for the low and high state, respectively. This behaviour is fairly easy to reproduce in comptonisation models in which the seed photons vary (see discussion), but the (marginally significant) lack of change in the absolute amount of reflection (so that the relative reflected fraction in the low state is larger than in the high state) is harder to explain. Clearly a model in which the reflected flux is produced at large distances from the source would be viable, but the reflected spectrum is broadened, so does contain at least some contribution from the relativistically smeared inner disk. For a $`10^8M_{}`$ black hole, then the inner 100 Schwarzschild radii are on scales of $`3\times 10^{15}`$ cm, i.e. less than 1 lightday. The spectra are taken at intervals of a day or more, so the relativistically smeared reprocessed component should not be appreciably lagged behind the source variability.
One possibility is that this anti–correlation of relative reflected fraction with flux is intrinsic to the source, that the geometry changes in such a way as to produce less reflection as the source brightens and steepens. However, it is very hard to see how this can be the case. A brighter, steeper source implies that there are more seed photons for the Compton scattering (see Figure 6 and the Discussion), i.e. that the geometry is such that more disk photons are intercepted by the source. Thus we can easily explain more reflected flux as the source steepens, but not less. A correlation of reflected solid angle with spectral index is indeed seen in both Galactic Black Hole Candidates (e.g. Życki et al. 1999) and AGN (Zdziarski et al. 1999). It seems much more likely that the (marginal) anti–correlation is an artifact of there being a second reflector at much larger distances from the source. The light travel time delay then means that the distant reprocessor does not have time to respond to rapid flux variations, so that as the source dims the relative contribution of the reflected spectrum from the distant reprocessor increases.
We include a reprocessed component from a torus, fixing its parameters to the best fit XTE values from the joint ASCA–RXTE fit (see table 2, again assuming twice solar abundance and inclination of $`60^{}`$, but this time also fixing the inner edge of the disk to $`60R_\mathrm{g}`$). The significance of the correlation of the spectral index with flux is unchanged ($`\chi _\nu ^2=44.7/58`$ for a constant while adding the linear term gives $`\chi _\nu ^2=29.4/57`$), while the flux/accretion disk reflection variability is now consistent with a constant solid angle ($`\chi _\nu ^2=12.4/58`$): adding a linear term gives an insignificant change ($`\chi _\nu ^2=11.4/58`$). The results are similar for a fixed Gaussian line rather than a full reprocessed spectrum.
We illustrate this again by fits to the high and low state spectra. Figure 7a shows the confidence contours for the power law spectral index, while Figure 7b shows the derived solid angle of the relativistic reflector. The diagonal lines denote solutions where the power law index and reflected fraction remain constant between the two datasets. The power law index is clearly variable between the high and low flux level datasets irrespective of how the reprocessor is modelled. With a single reprocessor then the relative amount of reflection is only consistent with remaining constant at the $`<90`$ per cent confidence level. The data prefer a larger contribution of reflected flux relative to the power law in the low state spectrum i.e. that the absolute normalization of the reflected flux is constant. The dashed and dotted lines show the same contours for a model including an unsmeared, neutral reprocessed spectrum and Gaussian line from the molecular torus, respectively. The reflected fraction is then consistent with a constant value.
### 3.4 PCA, HEXTE and OSSE Total Spectrum
All the RXTE PCA data were co–added to form a single spectrum (with 1 per cent systematic uncertainty added) and fit together with the HEXTE and OSSE data to give a broad band spectrum. Current consensus is that the continuum is formed by Compton scattering of soft seed photons by hot electrons. Such a Comptonised continuum can be approximated by a power law with exponential cutoff, but this becomes inaccurate if the spectrum extends close to the energies of either the seed photons or hot electrons. The seed photons are presumably from the accretion disk, with expected temperatures of $`10`$ eV, so the spectral curvature here is not an issue. However, the inclusion of the OSSE data means that the shape of the spectral cutoff from the electron temperature becomes important. Thus we use an analytical approximation to a Comptonised spectrum based on solutions of the Kompaneets equation (Lightman & Zdziarski 1987), where the shape of the high energy cutoff is sharper than an exponential rollover.
The relativistic reflection and line from this incident continuum are calculated, and included in the model fit. The results are detailed in Table 3 for inclinations of the relativistic reprocessor of $`30,60`$ and $`72^{}`$ and iron abundances of $`1`$ and $`2\times `$ solar. We see the same trend as in Table 1 in terms of the spectra preferring higher iron abundance, and these solutions constrain the electron temperature to be $`kT_e40100`$ keV (corresponding to an exponential e-folding energy of $`120300`$ keV). Figure 8 shows the best fit solution for twice solar abundance, inclined at $`60^{}`$. However, the poor signal–to–noise of the data at the highest energies means that this temperature is dependent on details of the reflection spectrum. For solar abundances the temperature is generally unconstrained.
We note that these data are consistent with previous observations of this source (Madejski et al. 1995), and with the mean Seyfert spectrum at high energies (Zdziarski et al. 1995). However, the temperatures of the Comptonizing medium derived from these data by Zdziarski et al. 1994) are rather higher ($`kT250`$ keV), due to their assumption of an exponential rollover as an approximation to a Comptonised cutoff. This also means that the derived plasma optical depths in these previous papers of $`\tau 0.1`$ are too low. Modelling the spectrum with more accurate comptonised spectra gives $`kT130`$ keV and $`\tau 1`$ (Zdziarski et al. 1996).
## 4 DISCUSSION
### 4.1 Fe K Line and Overall Spectral Shape
Firstly, we clearly see that not all AGN are consistent with a substantial solid angle of extreme, relativistically smeared reflection. The reprocessed component seen in MGC–6–30–15 is not necessarily typical of AGN in general. A similar result is seen in a recent analysis of the ASCA spectrum of NGC 5548 (Chiang et al. 1999), where the line is broad, but not so broad as expected from a disk which extends down to the last stable orbit around a black hole (for the June 15th data set they obtain $`R_{\mathrm{in}}=18.7_{9.5}^{+29.1}R_\mathrm{g}`$ for emissivity $`r^3`$; J. Chiang, private communication). Crucially, our data allow us to constrain a reflected component from a molecular torus. A torus with column $`10^{23}`$ cm<sup>-2</sup> can produce a strong, narrow 6.4 keV line, accompanied by a reflected continuum for columns $`10^{24}`$ cm<sup>-2</sup>. We do significantly detect such a contribution to the iron K$`\alpha `$ line, which may also be accompanied by a reflected continuum. However, even with a narrow line from the torus, the remaining line from the accretion is not as broad as that seen in MCG–6–30–15.
The amount of relativistically smeared reflection is rather less than unity for any inclination $`60^{}`$. Larger inclinations are not expected since this object is classified as a Seyfert 1 and no strong absorption from the torus is seen in the X–ray spectrum. Thus the reflection from the accretion disk in IC4329a looks very like that seen in the low state spectra of the galactic black hole systems (Życki et al. 1998; Done & Życki 1999) in having $`\mathrm{\Omega }/2\pi <1`$, relativistically smeared by velocities which are inconsistent with the reflecting material extending down to 3 Schwarzschild radii. There is then no intrinsic difference between the Galactic and extragalactic accreting black holes, but there is a spread in source properties in both classes (intrinsic spectral index, amount of reflection and amount of relativistic smearing). The AGN results to date show that steep spectra have a larger amount of reflection (Zdziarski et al. 1999) and more relativistic smearing. This is exactly the sequence seen in the Galactic Black Hole Candidate Nova Muscae 1991 (Życki et al. 1998; Życki et al. 1999) as a function of decreasing mass accretion rate.
There are currently two ways to explain the lack of extreme relativistic line. The first is to say that the inner accretion disk is simply not present, that it has been replaced by an X–ray hot flow. These composite truncated disk/hot X–ray source models were first proposed by Shapiro, Lightman and Eardley (1976) when they discovered a hot, two temperature, optically thin solution to the accretion flow equations, though this was subsequently shown to be unstable. Such models were given new impetus by the rediscovery of a related stable solution of the accretion flow equations (Narayan & Yi 1995), which include advective as well as radiative cooling (ADAFs). These ideas are clearly consistent with our results.
The alternative is that the inner disk is present, as required by the magnetic reconnection models for the X–ray flux, but that it cannot be seen in the reflected spectrum. One way to do this is if the upper layers of the disk are so ionized that they produce almost no atomic spectral features (Ross et al. 1999). Simple models for this, where the ionisation state of the disk varies as a smooth function of radius, do not match the data. However, the ionisation structure could be highly complex, with rapid transition between complete ionisation and relatively cool material (Różańska 1999; S. Nayakshin, private communication). Alternatively, if the X–ray source is moving away from the disk at transrelativistic velocities, perhaps because of plasma ejection from expanding magnetic loops, then its radiation pattern does not strongly illuminate the inner disk (Beloborodov 1999).
It is currently very difficult to distinguish observationally between these models, and all have some remaining theoretical problems. For the ADAF solutions, it is not yet known whether the fundamental assumptions underlying the solutions can hold, or how a transition from the cool disk to a hot flow can occur, while for the disk–corona geometry the uncertainties are mainly in the detailed outcome of magnetic reconnection, and in the ionisation structure of the illuminated disk.
### 4.2 Spectral Variability
Our data sample the source variability, which gives another way to investigate the underlying radiation mechanisms. This is only the second AGN where the reflected and intrinsic spectrum can be disentangled (the other is NGC 5548: Magdziarz et al. 1998; Chiang et al. 1999). The results show that the power law itself clearly gets intrinsically steeper as the source brightens, which allows to constrain the variability process. If there were merely more dissipation in the X–ray hot corona without an accompanying change in soft seed photon flux then the spectrum would harden as it got brighter (e.g. Ghisellini & Haardt 1994). Thus the observed spectral index–flux correlation implies that the soft seed photons also increase, and by somewhat more than the increase in the hard flux. Seed photons are thought to arise primarily through reprocessing, since the hard X–rays illuminating the disk which are not reflected are thermalised, emerging as soft photons. In this model the change in soft photons is commensurate with the change in flux dissipated in the hot corona. To change the seed photons by more than the change in hard dissipation requires either a change in geometry, such that the hard X–ray source intercepts a larger fraction of the disk radiation, or a decrease in reflection albedo, so that more of the incident hard X–ray radiation is thermalised rather than reflected. The former can be linked to the composite hot flow/cool accretion disk models as a result of varying the inner disk edge, while the latter could be produced in the disk–corona models if the ionisation state of the disk decreases for steeper spectra, so that the reflection albedo decreases and the thermalised soft flux increases. However, recent simultaneous observations of EUV and X–ray variability cast doubt on the simple scenario where the EUV seed photon flux is primarily reprocessed (Nandra et al. 1998; Chiang et al. 1999). The variability that we see could equally well be the result of a variable soft photon flux irradiating the X–ray region.
These three possibilities predict some differences in the behavior of the reflected continuum. If the disk geometry is changing to give more soft seed photons then we expect more solid angle of reflection as the spectrum steepens (as seen in the AGN/XRB compilation of Zdziarski et al. 1999). If the increased soft photons are from increasing thermalisation in the disk due to decreasing ionisation, then we should also see more cold reflection (as opposed to unobservable, completely ionized reflection) for steeper spectra. If it is simply the irradiating soft flux which is changing, without changing disk geometry then the reflected fraction should remain constant.
What we see is marginally (90 per cent confidence contour) consistent with a constant reflected fraction, although the data prefer that the relative amount of reflection decreases as the source increases and steepens. The resulting spectrum is consistent with the absolute normalisation of the reflected spectrum remaining constant as the source changes. Some part of the reflected spectrum could be contaminated by a line or reprocessed component from a molecular torus, which would be constant due to light travel time delays on the timescales of the monitoring campaign. Allowing for this results in the reflected fraction from the accretion disk being more convincingly consistent with a constant value, but still does not permit much of an increase with increasing flux or spectral index. The data then support the idea of a variable soft flux which is not reprocessed as the driver for the hard X–ray variability, but could also allow small changes in reflection geometry/ionisation.
We speculate that both variability mechanisms operate in Seyferts i.e. that there are spectral changes linked to changes in the geometry/ionisation (such as seen by Zdziarski et al. 1999), but that the soft seed photons can also vary independently of these changes, giving a second, subtly different source of spectral variability.
## 5 CONCLUSIONS
$``$ Not all AGN have the extreme relativistic line profiles expected from a disk extending down to the innermost stable orbit around a black hole. This is consistent with either the inner disk being truncated before the last stable orbit, or with an inner disk which produces no significant reflected features either through anisotropic illumination or extreme ionisation. Simple photo–ionisation models, where the ionisation varies smoothly as a function of radius can be ruled out by the data, but these may differ substantially from more detailed models of the ionisation structure.
$``$ There is intrinsic spectral variability, where the power law softens as the source brightens. This implies that the soft seed photons are increasing faster than the increase of the hard X–ray luminosity. The lack of a corresponding increase in the observed reflected spectrum implies that either the changes in disk inner radial extent/ionisation structure are small, or that the variability is actually driven by changes in the seed photons which are decoupled from the hard X–ray mechanism.
## 6 ACKNOWLEDGEMENTS
This research was supported in part by grant 2P03D01816 of the Polish State Committee for Scientific Research (KBN) and by NASA grant NAG 54106. We thank James Chiang for discussing with us their results on NGC 5548. We acknowledge the help of the OSSE team in reducing the OSSE data, and the RXTE Guest Observer Facility (and in particular, Tess Jaffe) for providing the PCA data reduction script rex.
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warning/0002/cond-mat0002279.html | ar5iv | text | # Quantum measurement of coherence in coupled quantum dots
## I Introduction
There have recently been a number of suggestions for a quantum computer architecture that use quantum dots of varying kinds. If these schemes are to be practical many important physical questions need to be answered, one of which is how to readout physical properties such as charge or spin at a single electron level. In this paper we present a quantum trajectory analysis of a general scheme to readout a single electronic qubit using a single electron transistor (SET). We adopt a general phenomenological description of the SET in which the tunneling rate through the SET is conditioned on the occupation or otherwise of a nearby quantum dot.
We consider two spatially separated quantum dots which are strongly coupled so that delocalized states of their relevant degrees of freedom can form. To be specific, we imagine each dot to have a single electronic bound state that can be occupied. Thus the average occupation number of each dot must be less than unity. This restriction can easily be removed to account for spin, or multiple electron states. We label each dot with an index $`1,2`$ and let $`c_i,c_i^{}`$ represent the Fermi annihilation and creation operators for each single electron state (see figure 1).
The two dots are strongly coupled via the tunnel coupling Hamilton
$$V=i\mathrm{}\frac{\mathrm{\Omega }}{2}(c_1^{}c_2c_2^{}c_1)$$
(1)
Thus the total Hamiltonian of the two-dot system is
$$H=\mathrm{}\underset{i=1}{\overset{2}{}}\omega _ic_i^{}c_i+V$$
(2)
In what follows we will work in an interaction picture and assume that the energies of each bound state are equal (again this can be relaxed). Coulomb blockade effects have been ignored at this stage, but can easily be included without significantly changing the results of this paper.
The single particle eigenstates of this Hamiltonian are even and odd superpositions of the bare states of each well. Such states are thus delocalized over the two-dot system and are sometimes called ’molecular states’ in the literature. The localized states can then be represented as an even and odd superposition of the delocalized states. The localized states are not stationary; rather, the system will periodically oscillate between them. That is to say the system will tunnel coherently between the two dots.
To this coherent system we add a measurement device which determines the presence of an electron on one of the dots, say dot $`1`$, which we shall refer to as the target (see figure 1). The model is based on a SET tunnel junction containing a single bound state on the island. The interaction between the target and the SET is via a Coulomb blockade. Thus the interaction Hamiltonian between the SET and the target must commute with the target electron number operator , $`c_1^{}c_1`$. This makes it a QND (quantum nondemolition measurement) of electron number. The Coulomb blockade changes the current flowing through the SET. In simple terms if there is no electron on the target the island state is biased so as to allow little or no current to flow though the SET. This is the quiescent state of the SET. However when there is an electron on the target, the coulomb blockade shifts the bound state on the island to allow a greater current to flow through the device (see figure 2)
We derive a master equation to describe the behaviour of the target system. This master equation describes the unconditional evolution of the measured system when the results of all measurement records (that is current records) are averaged over. This will tell us the rate at which coherence in the target system is destroyed by the measurement. However we also need to know how the system state depends on the actual current through the device in order to determine how quickly the conditional state of the electron becomes localized which is measure of the quality of the measurement.
One approach to this problem is to keep track of many different states of the system, corresponding to the different numbers of electrons which have tunneled through the SET. This is the approach used for example in Ref. . Here we adopt an alternate method which gives a more intuitive picture for the conditional dynamics. We use a conditional stochastic master equation which gives the evolution of the measured system, conditioned on a particular realization of the measured current. The instantaneous state of the target conditions the measured current while the measured current itself conditions the future evolution of the measured system in a self consistent manner. This approach to measurements has been variously called the quantum trajectory method or quantum monte carlo method .
## II SET Model
Consider a two dot system with the coupling in Eq. (1). If the electron is in dot-2 the quiescent rate of current tunneling though the SET is a constant which we will denote $`D_0`$. However if there is an electron on dot-1, the rate of tunneling through the SET changes to $`D_0+D_1`$ with $`D_1>0`$. If $`D_0`$ does not equal zero then a current spike (resulting from a tunnel event in the SET) does not necessarily imply that the electron in the measured system is in dot-1. In an ideal device the quiescent tunneling rate, $`D_0`$ is zero. In reality Johnson noise on the circuit containing the SET will give a non-zero quiescent tunneling current.
Assuming that the SET island state can be adiabatically eliminated, it is possible to derive a master equation for the state of the coupled dot system. This is done in the appendix, and the result is
$$\frac{d\rho }{dt}=i[H,\rho ]+\gamma _{\mathrm{dec}}𝒟[c_1^{}c_1]\rho =\rho $$
(3)
where the irreversible part is defined for arbitrary operators $`A`$ and $`B`$ by
$$𝒟[A]B=𝒥[A]B𝒜[A]B,$$
(4)
where
$`𝒥[A]B`$ $`=`$ $`ABA^{}`$ (5)
$`𝒜[A]B`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A^{}AB+BA^{}A).`$ (6)
The decoherence rate is given by
$$\gamma _{\mathrm{dec}}=2D_0+D_1.$$
(7)
The fact that the irreversible term is a function of the number operator in the target qubit is an indication that this describes a QND measurement of the occupation of the dot-1. It is easy to verify that the stationary solution of this master equation is an equal mixture of the two accessible electronic states.
The stochastic record of measurement ideally comprises a sequence of times, being the times at which electrons tunneled through the SET. In practice of course these events are not seen due to a finite frequency response of the circuit (including the SET) which averages each event over some time. However for the purpose of this paper we will take the zero response time limit. In this limit the current consists of a sequence of $`\delta `$ function spikes. Formally we can write $`i(t)=edN/dt`$, where $`dN(t)`$ is a classical point process which represents the number (either zero or one) of tunneling events seen in an infinitesimal time $`dt`$, and $`e`$ is the electronic charge. We can think of $`dN(t)`$ as the increment in the number of electrons $`N(t)`$ in the collector in time $`dt`$. It is this variable, the accumulated electron number transmitted by the SET, which is used in Ref. .
The point process $`dN(t)`$ is formally defined by the conditions
$`[dN(t)]^2`$ $`=`$ $`dN(t)`$ (8)
$`\mathrm{E}[dN(t)]/dt`$ $`=`$ $`D_0\mathrm{Tr}[(1n_1)\rho _\mathrm{c}(t)(1n_1)]+(D_0+D_1)\mathrm{Tr}[n_1\rho _\mathrm{c}(t)n_1]`$ (9)
$`=`$ $`D_0+D_1n_1_\mathrm{c}(t).`$ (10)
Here $`\mathrm{E}[x]`$ denotes a classical average of a classical stochastic process $`x`$, and
$$n_1=c_1^{}c_1$$
(11)
is the occupation number operator for the first dot. The first of these equations simply expresses the fact that $`dN(t)`$ equals zero or one. The second says that the rate of events is equal to a background rate $`D_0`$ plus an additional rate $`D_1`$ if and only if the electron is in the first dot.
In Eq. (10), the system state matrix $`\rho _\mathrm{c}(t)`$ is not the solution of the master equation (3). That is because if one has a record of the current $`dN/dt`$ through the SET then one knows more about the system than the master equation indicates. That is to say, $`\rho _\mathrm{c}(t)`$ is actually conditioned by $`dN(t^{})`$ for $`t^{}<t`$, hence the subscript c. The first way of writing Eq. (10) hints at how $`dN(t^{})`$ conditions $`\rho _\mathrm{c}(t)`$. From the appendix, the state at time $`t+dt`$ given $`dN(t)=1`$ is
$$\stackrel{~}{\rho }_1(t+dt)=dt\left[D_0(1n_1)\rho (t)(1n_1)+(D_0+D_1)n_1\rho (t)n_1\right]$$
(12)
This is an unnormalized state whose norm is equal to the probability of that event ($`dN(t)=1`$) occurring, as seen above in Eq. (10).
The normalized state can be written more elegantly as
$$\rho _1(t+dt)=\frac{(D_0+D_1𝒥[n_1]+2D_0𝒟[n_1])\rho (t)}{\mathrm{Tr}\{(D_0+D_1𝒥[n_1]+2D_0𝒟[n_1])\rho (t)\}}=\frac{(D_0𝒥[1n_1]+(D_0+D_1)𝒥[n_1])\rho (t)}{\mathrm{E}[dN(t)/dt]},$$
(13)
where $`𝒥`$ is as defined above in Eq. (5).
To write the full conditioned evolution we need to know the density operator $`\rho _0(t+dt)`$ given that $`dN(t)=0`$. This can be found from Eq. (13) plus the fact that when averaged over the observed classical point process $`dN`$,
$$\stackrel{~}{\rho }_0(t+dt)+\stackrel{~}{\rho }_1(t+dt)=(1+dt)\rho (t).$$
(14)
That is to say, on average the system still obeys the master equation (3). From this equation, we obtain
$`\stackrel{~}{\rho }_0(t+dt)`$ $`=`$ $`\rho (t)dt\{D_0𝒜[1n_1]\rho (t)+(D_0+D_1)𝒜[n_1]\rho (t)+i[H,\rho (t)]\},`$ (15)
$`=`$ $`\rho (t)dt\left\{D_0\rho (t)+D_1𝒜[n_1]\rho (t)i[H,\rho (t)]\right\},`$ (16)
where $`𝒜`$ is as defined above in Eq. (6). Once again this state is unnormalized and its norm gives the probability that $`dN(t)=0`$, that is
$$\mathrm{Tr}[\stackrel{~}{\rho }_0(t+dt)]=1\mathrm{E}[dN(t)]$$
(17)
Using the variable $`dN(t)`$ explicitly, the conditioned state at time $`t+dt`$ is
$$\rho _c(t+dt)=dN(t)\frac{\stackrel{~}{\rho }_1(t+dt)}{\mathrm{Tr}[\stackrel{~}{\rho }_1(t+dt)]}+[1dN(t)]\frac{\stackrel{~}{\rho }_0(t+dt)}{\mathrm{Tr}[\stackrel{~}{\rho }_0(t+dt)]}.$$
(18)
Since $`dN(t)`$ is almost always zero we can set $`dN(t)dt=0`$ and expand this expression to finally obtain the stochastic master equation, conditioned on the observed event in time $`dt`$
$`d\rho _\mathrm{c}`$ $`=`$ $`dN(t)\left[{\displaystyle \frac{D_0+D_1𝒥[n_1]+2D_0𝒟[n_1]}{D_0+D_1\mathrm{Tr}[\rho _\mathrm{c}n_1]}}1\right]\rho _\mathrm{c}`$ (20)
$`+dt\left\{D_1𝒜[n_1]\rho _\mathrm{c}+D_1\mathrm{Tr}[\rho _\mathrm{c}n_1]\rho i[H,\rho _\mathrm{c}]\right\}`$
Note that averaging this equation over the observed stochastic process (by setting $`dN`$ equal to its expected value) gives the unconditional master equation (3).
## III Average steady state properties
We now analyze in some detail the ensemble averaged properties of the system based on the unconditional master equation. In particular we calculate the stationary noise power spectrum of the current through the SET when there is the possibility of coherent tunneling between dot-2 and the measured dot-1. The details of how the quantum stochastic processes in the SET determine the average current though the SET are given in reference . The link with the stochastic formalism of the preceding section is that the current $`i(t)`$ through the SET is given by
$$i(t)=e\frac{dN(t)}{dt}.$$
(21)
First we calculate the steady state current
$`i_{\mathrm{}}`$ $`=`$ $`\mathrm{E}[i(t)]_{\mathrm{}}`$ (22)
$`=`$ $`e(D_0+D_1c_1^{}c_1_{\mathrm{}})`$ (23)
$`=`$ $`e(D_0+{\displaystyle \frac{D_1}{2}}),`$ (24)
where the $`\mathrm{}`$ subscript indicates that the system is at steady-state. The fluctuations in the observed current, $`i(t)`$ are quantified by the two-time correlation function:
$`G(\tau )`$ $`=`$ $`\mathrm{E}[i(t)i(t+\tau )i_{\mathrm{}}^2]_{\mathrm{}}`$ (25)
$`=`$ $`ei_{\mathrm{}}\delta (\tau )+i(t),i(t+\tau )_{\mathrm{}}^{\tau 0}.`$ (26)
Here $`A,BABAB`$. The fact that the multiplier of the shot noise is $`ei_{\mathrm{}}`$ rather than the usual $`(e/2)i_{\mathrm{}}`$ is because of the approximation we have made in treating the SET. Specifically, we have adiabatically eliminated the SET by taking the limit where any electron which tunnels onto the SET island from the emitter immediately tunnels off to the collector. This means that, on the time scales we are interested in, there is a perfect correlation between the emitter current and collector current. This leads to a doubling of the shot noise level. Of course, at very high frequencies, higher than we are interested in, the true shot noise level of $`(e/2)i_{\mathrm{}}`$ could still be seen in principle.
To relate these classical averages to the fundamental quantum processes occurring in the well we apply the theory of open quantum system to the present system. Specifically, we can relate the correlation function for the current to the following quantum averages
$$i(t)i(t+\tau )_{\mathrm{}}^{\tau 0}=e^2\mathrm{Tr}\left[(D_0+D_1𝒥[n_1]+2D_0𝒟[n_1])e^\tau (D_0+D_1𝒥[n_1]+2D_0𝒟[n_1])\rho _{\mathrm{}}\right].$$
(27)
Because $`\rho _{\mathrm{}}`$ is an equal mixture of the two electron states, it satisfies
$$𝒟[n_1]\rho _{\mathrm{}}=0.$$
(28)
In addition, the following identies for arbitrary operators $`A`$ and $`B`$ are easy to prove: $`\mathrm{Tr}\left[𝒟[A]B\right]0`$, $`\mathrm{Tr}\left[𝒥[n_1]B\right]\mathrm{Tr}[n_1B]`$, $`\mathrm{Tr}[e^\tau B]=\mathrm{Tr}[B]`$, and $`\mathrm{Tr}[Ae^\tau \rho _{\mathrm{}}]=\mathrm{Tr}[A\rho _{\mathrm{}}]`$. Using these simplfications we obtain
$$i(t),i(t+\tau )_{\mathrm{}}^{\tau 0}=D_1^2e^2\left\{\mathrm{Tr}\left[n_1e^\tau 𝒥[n_1]\rho _{\mathrm{}}\right]\mathrm{Tr}[n_1\rho _{\mathrm{}}]^2\right\}.$$
(29)
Evaluating this expression we find
$$G(\tau )=ei_{\mathrm{}}\delta (\tau )+\frac{e^2D_1^2}{8}\left(\frac{\mu _+e^{\mu _{}\tau }\mu _{}e^{\mu _+\tau }}{\sqrt{(\gamma _{\mathrm{dec}}/4)^2\mathrm{\Omega }^2}}\right),$$
(30)
where
$$\mu _\pm =(\gamma _{\mathrm{dec}}/4)\pm \sqrt{(\gamma _{\mathrm{dec}}/4)^2\mathrm{\Omega }^2},$$
(31)
and where the first term represents the shot noise component as discussed above. The power spectrum of the noise is
$$S(\omega )=_0^{\mathrm{}}𝑑\tau G(\tau )2\mathrm{cos}(\omega \tau ),$$
(32)
which evaluates to
$$S(\omega )=ei_{\mathrm{}}+\frac{e^2D_1^2\mathrm{\Omega }^2/2}{\sqrt{(\gamma _{\mathrm{dec}}/4)^2\mathrm{\Omega }^2}}\left\{\frac{1}{\mu _+^2+\omega ^2}\frac{1}{\mu _{}^2+\omega ^2}\right\}.$$
(33)
In the case that $`\mathrm{\Omega }>\gamma _{\mathrm{dec}}/4`$ the spectrum will have a double peak structure indicating that coherent tunneling is taking place between the two coupled dots. For smaller $`\mathrm{\Omega }`$ only a single peak appears in the spectrum. We can thus use the noise power spectrum of the current though the SET as a means to measure the tunnel coupling between dots if the tunnel coupling is high enough. We illustrate this in figure 3.
## IV Analytical Results for Conditional Dynamics
We now return to the stochastic master equation for the conditioned state,
$`d\rho _\mathrm{c}`$ $`=`$ $`dN\left[{\displaystyle \frac{D_0+D_1𝒥[c_1^{}c_1]+2D_0𝒟[c_1^{}c_1]}{D_0+D_1\text{Tr}[\rho _\mathrm{c}c_1^{}c_1]}}1\right]\rho _\mathrm{c}`$ (35)
$`+dt\left\{D_1{\displaystyle \frac{1}{2}}\{c_1^{}c_1,\rho _\mathrm{c}\}+D_1\text{Tr}[\rho _\mathrm{c}c_1^{}c_1]\rho i[H,\rho _\mathrm{c}]\right\}.`$
Comparing this to the unconditional master equation
$$\dot{\rho }=i[H,\rho ]+\gamma _{\mathrm{dec}}𝒟[c_1^{}c_1]\rho $$
(36)
we see that decoherence between the two coupled dots, 1 and 2, takes place at the rate $`\gamma _{\mathrm{dec}}=2D_0+D_1`$, but that the system decides between the two possibilities (electron on dot-1 or on dot-2) on a time scale that depends on $`D_1`$ and $`D_0`$ in some more complicated way. Of course this measurement time scale is necessarily at least as large as the decoherence time scale because successfully distinguishing between the two dots would by definition destroy any coherence between them.
The different measurement time scales can be derived most easily by introducing the Bloch representation of the state matrix:
$$\rho =\frac{1}{2}\left(I+x\sigma _x+y\sigma _y+z\sigma _z\right)$$
(37)
where the Pauli matrices are defined using the Fermi operators for the two dots
$`\sigma _x`$ $`=`$ $`c_1^{}c_2+c_2^{}c_1`$ (38)
$`\sigma _y`$ $`=`$ $`ic_1^{}c_2+ic_2^{}c_1`$ (39)
$`\sigma _z`$ $`=`$ $`c_2^{}c_2c_1^{}c_1`$ (40)
In this representation the means of the Pauli matrices $`\sigma _\alpha `$ are given by the respective coefficient $`\alpha `$, with $`\alpha =x,y,z`$.
The stochastic master equation can now be written as a set of coupled stochastic differential equations for the Bloch sphere variables as
$`dz_\mathrm{c}`$ $`=`$ $`\mathrm{\Omega }x_\mathrm{c}dt+{\displaystyle \frac{D_1}{2}}(1z_\mathrm{c}^2)dtdN(t){\displaystyle \frac{D_1(1z_c^2)/2}{D_0+D_1(1z_c)/2}}`$ (41)
$`dx_\mathrm{c}`$ $`=`$ $`\mathrm{\Omega }z_\mathrm{c}dt{\displaystyle \frac{D_1}{2}}z_\mathrm{c}x_\mathrm{c}dtdN(t)x_\mathrm{c}`$ (42)
$`dy_\mathrm{c}`$ $`=`$ $`{\displaystyle \frac{D_1}{2}}z_\mathrm{c}y_\mathrm{c}dtdN(t)y_\mathrm{c}.`$ (43)
Again the c-subscript is to emphasize that these variables refer to the conditional state. If we average over the noise, the ensemble dynamics is then seen to be given by
$`{\displaystyle \frac{dz}{dt}}`$ $`=`$ $`\mathrm{\Omega }x`$ (44)
$`{\displaystyle \frac{dx}{dt}}`$ $`=`$ $`\mathrm{\Omega }z{\displaystyle \frac{\gamma _{\mathrm{dec}}}{2}}x`$ (45)
$`{\displaystyle \frac{dy}{dt}}`$ $`=`$ $`{\displaystyle \frac{\gamma _{\mathrm{dec}}}{2}}y`$ (46)
where $`\alpha =\mathrm{E}[\alpha _\mathrm{c}]`$ denotes the averaging over the ensemble of conditional states. These equations are exactly what would be obtained directly from the ensemble averaged master equation Eq(3). In particular we note that the average population difference $`z`$ between the dots is a constant of the motion in the absence of any free Hamiltonian. However the stochastic differential equations enable us to calculate important averages that are not obtainable from the master equation. For example, if the model does indeed describe a measurement of $`c_1^{}c_1=(1\sigma _z)/2`$, then, in the absence of tunneling, we would expect to see the conditional state become localized at either $`z=1`$ or $`z=1`$. Indeed for $`\mathrm{\Omega }=0`$ we can see from the conditional equation for $`z_\mathrm{c}`$ that $`z_\mathrm{c}=\pm 1`$ is a fixed point.
We can take into account both fixed points by considering $`z_\mathrm{c}^2`$. In the absence of tunneling this must must approach $`1`$ for all trajectories, since the system will eventually become localized due to the measurement in one dot or the other. Therefore it is sensible to take the ensemble average $`\mathrm{E}[z_\mathrm{c}^2]`$ and find the rate at which this deterministic quantity approaches one. Noting that for a stochastic variable we have $`d(z^2)=2zdz+dzdz`$, and that $`\mathrm{E}[dN^2]=\mathrm{E}[dN]=[D_0+D_1(1z_\mathrm{c})/2]dt`$, we find that
$$\frac{d\mathrm{E}[z_\mathrm{c}^2]}{dt}=\mathrm{E}\left[\frac{D_1^2(1z_\mathrm{c}^2)^2}{4D_0+2D_1(1z_\mathrm{c})}\right].$$
(47)
If the system starts state which has an equal probability for single electron to be on each dot then $`z_\mathrm{c}(0)=0`$ and in the ensemble average this would remain the case. However if we ensemble average $`z_\mathrm{c}^2`$ over many quantum trajectories then for short times we find
$$\mathrm{E}[z_\mathrm{c}^2(\delta t)]=\frac{D_1^2}{4D_0+2D_1}\delta t$$
(48)
That is to say, the system tends towards a definite state (with $`z_\mathrm{c}=\pm 1`$ so $`z_\mathrm{c}^2=1`$) at an initial rate of $`D_1^2/(4D_0+2D_1)`$. For vanishing $`D_0`$, this is the same as the decoherence rate, $`D_1/2`$, as expected. But for $`D_0D_1`$, the rate goes to $`(D_1/2D_0)\times D_1/2D_1/2`$. That is, the rate at which the system becomes localized at one or the other dot is much less than the decoherence rate. This result cannot be obtained from the ensemble averaged master equation alone. It is a direct reflection of the fact that for $`D_00`$ a tunneling event cannot be unambiguously attributed to the location of an electron on the double dot system. As the rate of localization is a direct indication of the quality of the measurement, we can use the localization rate defined as
$$\gamma _{\mathrm{loc}}=\frac{D_1^2}{4D_0+2D_1}$$
(49)
as an important parameter defining the quality of the measurement. This parameter is related to the signal-to-noise ratio for this measurement as we now show.
For a Poisson process at rate $`R`$, the probability for $`m`$ events to occur in time $`T`$ is
$$p(m;T)=\frac{(RT)^m}{m!}e^{RT}$$
(50)
The mean and variance of this distribution are equal and given by $`\mathrm{E}(m)=\mathrm{Var}(m)=RT`$. Now consider an electron which is, with equal likelihood, in either dot, so that $`z=0`$. If the electron is in dot 1 then the rate of electrons passing through the SET is $`D_0+D_1`$; if it is in dot two then it is just $`D_0`$. These two possibilities will begin to be distinguishable when the difference in the means of the two distributions $`p(m,T)`$ is of order the square root of the sum of the variances. That is, when
$$D_1T\sqrt{D_0T+(D_0+D_1)T}$$
(51)
Solving this for $`T`$ gives a characteristic rate
$$T^1\frac{D_1^2}{2D_0+D_1}.$$
(52)
The right hand side of this expression is simply twice the $`\gamma _{\mathrm{loc}}`$ defined above. A similar conclusion is reached in reference .
In the ideal limit of no quiescent current in the SET $`D_0=0`$, the stochastic master equation can be replaced by a stochastic Schrödinger equation, and will collapse to a single possibility at a rate $`D_1`$ which is the same as the decoherence rate. The effect of $`D_0`$ is most clearly seen in the other limit, $`D_1D_0`$, as noted above. In this limit the single electron makes only a small relative change in the tunneling rate through the SET. As the rate of jumps also becomes large, the trajectories in this limit take on the appearance of diffusion rather than jumps. The rate at which the electron localizes into one well or the other scales as $`D_1^2D_0`$, which is much longer than the decoherence time scale $`D_0^1`$.
## V Numerical Simulations of Conditional dynamics
We now turn to numerical simulations of the conditional evolution and to estimate the conditions for a good measurement. Unlike traditional condensed matter measurements we wish to describe repeated measurements made on a single quantum system rather than a single measurement made upon an ensemble of systems. To do this we use the conditional dynamics of the measured system given a particular measurement record as described by the above stochastic evolution equation (35)
We return to the Bloch description defined in Eq. (37). In what follows we will assume that $`y(0)=0`$. For the form of tunneling used here the value of $`y`$ does not in fact change under either conditional or ensemble averaged dynamics. If the conditional state of the system remains in a pure state then $`x_\mathrm{c}^2+y_\mathrm{c}^2+z_\mathrm{c}^2=x_\mathrm{c}^2+z_\mathrm{c}^2=1`$. As noted previously this can only occur if the bare tunneling rate ($`D_0`$) is zero, when a tunneling event can unambiguously be attributed to the occupation of the target dot and no information is lost about the state of the system. In the more realistic case in which $`D_00`$, we can use the quantity $`x_\mathrm{c}^2+z_\mathrm{c}^2`$ as a measure of the purity of the sate, or equivalently as a measure of how much information the conditional record of measurements gives about the actual state of the two coupled dots. If the conditional state is a maximally mixed state of a two state system then $`x_\mathrm{c}^2+z_\mathrm{c}^2=0`$. We now describe in detail the numerical simulation of the conditional dynamics.
### A No Background Current
First we consider the case $`D_0=0`$, so the system is always in a pure state. Typical trajectories are shown in Fig. 4 for various values of $`\mathrm{\Omega }`$. For small $`\mathrm{\Omega }D_1/2`$ we see little evidence for coherent tunneling. Most of the time the electron is localized almost entirely in one well or the other. However, there is an asymmetry between the wells. A transition from dot-2 into dot-1 (the target) is sudden, occurring whenever an electron tunnels through the SET. A transition the other way takes a time of order $`2/D_1`$. This time is still much smaller than the average time between state-changing transitions, which can be shown analytically to be $`D_1/\mathrm{\Omega }^2`$. Thus over a long time, as shown in Fig. 4(a), the system still has the appearance of a random telegraph. This behaviour gives rise to a single-peaked noise spectrum, as shown in Fig. 3(a).
For moderate $`\mathrm{\Omega }D_1/2`$ the system is no longer well-localized in one dot or the other. Rather, the dynamics is complicated with clearly non-sinusoidal oscillations from one dot to the other interspersed with jumps into the target dot. The fact that oscillations are present gives peaks in the current noise spectrum, as shown in Fig. 3(b). The position of these peaks is a frequency less than $`\mathrm{\Omega }`$, as shown analytically in Sec. III.
For $`\mathrm{\Omega }D_1/2`$ the dynamics once again becomes simple, with nearly sinusoidal oscillations interspersed with jumps which occur with an average rate of $`D_1/2`$. This corresponds to a noise spectrum having a very sharp feature at $`\omega \pm \mathrm{\Omega }`$, as shown in Fig. 3(c).
The change in behaviour as $`\mathrm{\Omega }`$ increases is summarized in Fig. 5. There we plot E$`[z_\mathrm{c}^2]`$ versus $`\mathrm{\Omega }/D_1`$. The quantity E$`[z_\mathrm{c}^2]`$ measures how well localized the electron is at one well or the other, and would be $`1`$ if the electron were always localized and $`0`$ if it were never localized. It is actually possible to calculate this quantity numerically without using a stochastic ensemble, as follows.
With no background current, every time an electron tunnels through the SET the electron on the dots is known to be on dot $`1`$. If there are no further SET tunneling events for a time $`t`$ later then from Eq. (15), the system evolves up to that time by the equation
$$d\stackrel{~}{\rho }_0(t)=dt\{D_1\{c_1^{}c_1,\stackrel{~}{\rho }_0(t)\}/2+i[H,\stackrel{~}{\rho }_0(t)]\}.$$
(53)
Because there is no background current, and because the initial state is pure, it is possible to rewrite this in terms of a non-Hermitian Schrödinger equation
$$d|\stackrel{~}{\psi }_0(t)=dt(iH+D_1c_1^{}c_1/2)|\stackrel{~}{\psi }_0(t).$$
(54)
Here it must be remembered that the norm of this state represents the probability for the event that no electron has passed through the SET since the last one a time $`t`$ ago:
$$p_0(t)=\mathrm{Tr}[\stackrel{~}{\rho }_0(t)]=\stackrel{~}{\psi }_0(t)|\stackrel{~}{\psi }_0(t).$$
(55)
It is not difficult to show that the solution to Eq. (54) satisfying the initial condition $`|\psi _0(0)=|1,0`$ is
$$|\stackrel{~}{\psi }_0(t)=\alpha (t)|1,0+\beta (t)|0,1,$$
(56)
where the occupation numbers refer to the dots one and two in order. Here $`\alpha `$ and $`\beta `$ are real numbers defined by
$`\alpha (t)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _+\lambda _{}}}\left(\lambda _+e^{\lambda _+t}\lambda _{}e^{\lambda _{}t}\right),`$ (57)
$`\beta (t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }/2}{\lambda _+\lambda _{}}}\left(e^{\lambda _+t}e^{\lambda _{}t}\right),`$ (58)
where
$$\lambda _\pm =\frac{1}{2}\left[\frac{D_1}{2}\pm \sqrt{\left(\frac{D_1}{2}\right)^2\mathrm{\Omega }^2}\right]$$
(59)
The conditioned quantum expectation value for $`\sigma _z`$ is
$$z_0(t)=\frac{\stackrel{~}{\psi }_0(t)|\sigma _z|\stackrel{~}{\psi }_0(t)}{p_0(t)}=\frac{\beta ^2\alpha ^2}{\beta ^2+\alpha ^2}$$
(60)
Now in steady state the probability $`p_0(t)`$ that there is no increment in $`N(t)`$ for a time $`t`$ ago is related to the probability $`q_0(t)`$ that the last increment was a time $`t`$ ago by
$$q_0(t)=\frac{p_0(t)}{_0^{\mathrm{}}p_0(s)𝑑s}.$$
(61)
Since at steady state all conditioned states are uniquely identified by how long it has been since the last SET event, the ensemble average for $`z_c^2`$ is simply given by
$`\mathrm{E}[z_c^2]`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}q_0(t)[z_0(t)]^2𝑑t`$ (62)
$`=`$ $`\left[{\displaystyle _0^{\mathrm{}}}\left(\beta ^2+\alpha ^2\right)𝑑t\right]^1{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{\left(\beta ^2\alpha ^2\right)^2}{\beta ^2+\alpha ^2}}𝑑t.`$ (63)
Unfortunately it does not appear possible to evaluate the second integral here analytically. However a numerical integration is easy. The results, shown in Fig. 5, is in agreement with the ensemble averages obtained numerically using the stochastic master equation.
### B A Finite Background Current
We next consider the case where $`D_00`$. We show two plots, both with $`\mathrm{\Omega }=D_1`$, which is a regime in which coherent tunneling is clearly evident in the current noise spectrum. The first plot, in Fig. 6, is for $`D_0=D_1`$. Here coherent oscillations are still evident in $`z`$, but $`z`$ rarely attains its extreme values of $`\pm 1`$. The conditioned state is no longer pure, even immediately after a count. Also, the conditioned state following a count now depends on the state before the count. For this reason an exact solution by the method of the preceding section is impossible.
The second plot, in Fig. 7, is for $`D_0=10D_1`$. Here coherent oscillations are no longer obvious in the condition mean of $`\sigma _z`$, even thought they are present in the spectrum (as small features above the shot noise), as calculated in Sec. III. In this regime $`D_0D_1`$ so the diffusive limit discussed at the end of Sec. IV applies.
## VI Discussion and Conclusion
The three parameters we need to compare our theoretical results with experiment are $`\mathrm{\Omega }`$, $`D_0`$ and $`D_1`$. The two incoherent tunneling rates can be obtained by considering how they determine the steady state current through the device. This is given in Eq (24). Our model implicitly assumes that the quiescent noise in the SET is shot-noise limited, based as it is on elementary tunneling events. However from the point of view of the macroscopic circuit in which the SET is placed, the tunnel junctions appear as a capacitor in series with a resistor ( see note 14 in reference ). If the resistance of the junction is $`R`$ and the capacitance is $`C`$ the fundamental time constant for the junction is $`\frac{1}{RC}`$. This sets an upper bound for the tunneling rate $`\gamma \frac{1}{RC}`$ For a typical Al/AlO<sub>x</sub>/Al junctions we have $`C0.2f`$F and $`R50k\mathrm{\Omega }`$ and thus the time constant is $`\gamma 10^{11}\text{s}^1`$. For a double tunnel junction device, the maximum conductance is achieved for a symmetric pair of tunnel junctions. The value of $`D_1`$ in this case is given by $`D_1=\gamma /2`$ where $`\gamma `$ is the tunneling rate of the SET under conditions of maximum conductance (see appendix). We thus estimate that $`D_15\times 10^{10}\text{s}^1`$. The background tunneling rate through the SET depends on temperature as well as the bias conditions (see appendix). Typical maximum and minimum conductance for the SET at different temperatures have been measured by Joyez et al. . At a temperature of $`100m\text{K}`$ the minimum conductance is approximately zero, and thus at this temperature we can safely take $`D_0=0`$, the zero temperature result. However at a temperature of $`400\text{mK}`$ the ratio of the maximum to minimum conductance is 2.2. This indicates that at 400 mK, $`D_0=4\times 10^{10}\text{s}^1`$.
The value we choose for the tunneling rate depends strongly on the particular quantum dot system. We will consider the value appropriate for the single electron measurement scheme of Kane et al.. In this model the two localized states correspond to an electron on a tellurium ion donor or at a nearby interface below the donor. The tunneling rate is then between the donor state and the interface state. Kane et al. estimate that for this system $`\mathrm{\Omega }10^9\text{s}^1`$. This value is too low to observe a peak in the noise power spectrum away from zero frequency. However the tunneling rate could be increased by changing the donor interface bias voltage.
In conclusion we have presented a simple model to describe single electron measurements on a coupled double quantum dot system using an SET. We have given both the stationary (ensemble averaged) properties of the current through the SET as well as the conditional dynamics of repeated measurements on a single system. This illustrates how quantum trajectory methods may be naturally adapted to single electronics, and aid in the interpretation of ensemble averaged properties. We believe these models will prove useful in current attempts to fabricate quantum logic gates in solid state devices.
## A Derivation of the master equation.
It will suffice to consider a single quantum dot near the SET. This allows us to remove any reference to the electron field labeled by $`c_2,c_2^{}`$. The SET is modeled as a single biased double barrier (single well) device with a single bound state on the well described by the Fermi operators $`b,b^{}`$. The total Hamiltonian for the system including the reservoirs is
$$H=H_0+H_{CB}+H_{RT}+H_{LT}$$
(A1)
The term $`H_{CB}`$ is the Coulomb blockade term and is given by
$$H_{CB}=\mathrm{}\chi c_1^{}c_1b^{}b$$
(A2)
where $`\mathrm{}\chi `$ is the Coulomb blockade energy gap (see figure 2). Note this term can only be nonzero if there is an electron on the island and on the dot, in which case the energy of the island electron is increased. The terms $`H_{RT},H_{LT}`$ described the tunneling between the many modes in the left and right ohmic contacts and the bound state on the SET
$`H_{LT}`$ $`=`$ $`{\displaystyle \underset{k}{}}T_{Lk}a_{Lk}^{}b+T_{Lk}^{}a_{Lk}b^{}`$
$`H_{RT}`$ $`=`$ $`{\displaystyle \underset{k}{}}T_{Rk}a_{Rk}^{}b+T_{Rk}^{}a_{Rk}b^{}`$
where $`a_{Lk},a_{Rk}`$ are respectively the Fermi field annihilation operators for the left and right reservoir states at momentum $`k`$. The tunneling matrix elements between respectively the left and right Ohmic contacts and the island are $`T_{Lk},T_{Rk}`$. The free Hamiltonian for the the system is
$$H_0=\mathrm{}\underset{k}{}\omega _k^La_{Lk}^{}a_{Lk}+\omega _k^Ra_{Rk}^{}a_{Rk}+\mathrm{}\omega _1c_1^{}c_1+\mathrm{}\omega _0b^{}b$$
(A3)
We now transform to an interaction picture to remove the terms $`H_0+H_{CB}`$. The dynamics in the Schrödinger picture is now described by the time dependent Hamiltonian
$`H_I(t)`$ $`=`$ $`{\displaystyle \underset{k}{}}T_{Lk}a_{Lk}^{}be^{i\chi tc_1^{}c_1}e^{i(\omega _k^L\omega _0)t}`$
$`+T_{Lk}^{}a_{Lk}b^{}e^{i\chi tc_1^{}c_1}e^{i(\omega _k^L\omega _0)t}`$
$`T_{Rk}a_{Rk}^{}be^{i\chi tc_1^{}c_1}e^{i(\omega _k^R\omega _0)t}`$
$`+T_{Rk}^{}a_{Rk}b^{}e^{i\chi tc_1^{}c_1}e^{i(\omega _k^R\omega _0)t}`$
Using the fact that $`(c_1^{}c_1)^n=c_1^{}c_1`$ we find
$`H_I(t)=H_1(t)+H_2(t)`$
where
$`H_1(t)`$ $`=`$ $`(1c_1^{}c_1){\displaystyle \underset{k}{}}(T_{Lk}a_{Lk}^{}be^{i(\omega _k^L\omega _0)t}+\mathrm{H}.\mathrm{c}.)+(T_{Rk}a_{Rk}^{}be^{i(\omega _k^R\omega _0)t}+\mathrm{H}.\mathrm{c}.)`$
$`H_2(t)`$ $`=`$ $`c_1^{}c_1{\displaystyle \underset{k}{}}(T_{Lk}a_{Lk}^{}be^{i(\omega _k^L\omega _0\chi )t}+\mathrm{H}.\mathrm{c}.)+(T_{Rk}a_{Rk}^{}be^{i(\omega _k^R\omega _0\chi )t}+\mathrm{H}.\mathrm{c}.)`$
Notice that if there is no electron on the dot and $`c_1^{}c_10`$ then the second term is zero and the first term is a standard tunneling interaction onto a bound state with energy $`\mathrm{}\omega _0`$. On the other hand if there is an electron on the dot $`c_1^{}c_10`$ and the first term is zero and the second term is a standard tunneling interaction onto a bound state with energy $`\mathrm{}(\omega _0+\chi )`$ as expected.
The derivation of the master equation for the state matrix $`R`$ for the system (SET and quantum dot) can now proceed using standard techniques which we will sketch. The objective is to obtain a semigroup evolution in Lindblad form (that is to say positivity-preserving irreversible dynamics) for the state of the SET island and the dot alone with no reference to the ohmic contacts. The ohmic contacts are treated as perfect Fermi thermal reservoirs with a very fast relaxation constants. Each ohmic contact (left and right) remains in thermal equilibrium with chemical potentials $`\mu _L,\mu _R`$, but the total system is not in thermal equilibrium due to the external bias potential, $`V`$ with $`eV=\mu _L\mu _R`$ (see references for further discussion). We first define a time interval $`\delta t`$ which is slow compared to the dynamics of the island and the dot but very long compared to the time scale in which the ohmic contacts relax back to their steady state. The change in the state matrix $`W`$ of the system (SET and dot) and environment (Ohmic contacts) from time $`t`$ to $`t+\delta t`$, to second order in the tunnel coupling energy, is given by
$$W(t+\delta t)=W(t)i\delta t[H_I(t),W(t)]\delta t_t^{t+\delta t}𝑑t_1[H_I(t),[H_I(t_1),W(t_1)]]$$
(A4)
We now make the first Markov approximation and assume that at any time the state of the total system may be approximated by $`W(t)=R(t)\rho _L\rho _R`$ , that is to say the left and right ohmic contacts instantaneously relax back to Fermi distributions. We now obtain an evolution equation for $`R(t)`$, the state of the island and the dot by tracing over the reservoirs. The result is
$`{\displaystyle \frac{dR(t)}{dt}}`$ $`=`$ $`[\gamma _L(1f_L(\omega _0))+\gamma _R(1f_R(\omega _0)]𝒟[b(1c_1^{}c_1)]R`$
$`+\left[\gamma _Lf_L(\omega _0)+\gamma _Rf_R(\omega _0)\right]𝒟[b^{}(1c_1^{}c_1)]R`$
$`+\left[\gamma _L^{}(1f_L(\omega _0+\chi ))+\gamma _R^{}(1f_R(\omega _0+\chi ))\right]𝒟[bc_1^{}c_1]R`$
$`+\left[\gamma _L^{}f_L(\omega _0+\chi )+\gamma _R^{}f_R(\omega _0+\chi )\right]𝒟[b^{}c_1^{}c_1]R`$
where for arbitrary operators $`A`$ and $`B`$, $`𝒟[A]B=ABA^{}\frac{1}{2}(A^{}ABBA^{}A)`$ and where $`f_{L,R}(\omega )`$ is the Fermi filling probability for the left/right ohmic contact at the energy $`\mathrm{}\omega `$. The rates $`\gamma _{L,R}`$ and $`\gamma _{L,R}^{}`$ determine the rate of injection from the left ohmic contact into the island or emission from the island into the right ohmic contact under the conditions of no electron on the dot (unprimed) and with an electron on the dot (primed). These are evaluated using the second markov approximation as
$`\gamma _L`$ $`=`$ $`|T_{Lk_0}|^2,`$ (A5)
$`\gamma _R`$ $`=`$ $`|T_{Rk_0}|^2,`$ (A6)
$`\gamma _L^{}`$ $`=`$ $`|T_{Lk_0^{}}|^2`$ (A7)
$`\gamma _R^{}`$ $`=`$ $`|T_{Rk_0^{}}|^2`$ (A8)
where $`k_0=\sqrt{2m\omega _0/\mathrm{}}`$ and $`k_0^{}=\sqrt{2m(\omega _0+\chi )/\mathrm{}}`$.
The ideal quiescent state of the SET is defined as $`f_L(\omega _0)=1,f_R(\omega _0)=1`$ while $`f_L(\omega _0+\chi )=1,f_R(\omega _0+\chi )=0`$. Under these conditions the master equation reduces to
$$\frac{dR}{dt}=\gamma _R𝒟[b(1c_1^{}c_1)]R+\gamma _L𝒟[b^{}(1c_1^{}c_1)]R+\gamma _R^{}𝒟[bc_1^{}c_1]R+\gamma _L^{}𝒟[b^{}c_1^{}c_1]R$$
(A9)
We now wish to derive a master equation for the state matrix $`\rho `$ for the dot alone. This is easiest if we assume that $`\gamma _R,\gamma _R^{}`$ are much larger than all other rates in the system. In this case it is possible to adiabatically eliminate the SET island using techniques similar to that in Ref. . We expand the state matrix $`R`$ in powers of $`1/\gamma _R`$ or $`1/\gamma _R^{}`$ as
$$R=\rho _0|00|+\rho _1|00|.$$
(A10)
The equations of motion for $`\rho _1`$ and $`\rho _0`$ are
$`\dot{\rho }_1`$ $`=`$ $`\gamma _R𝒜[1n_1]\rho _1+\gamma _L𝒥[1n_1]\rho _0\gamma _R^{}𝒜[n_1]\rho _1+\gamma _L^{}𝒥[n_1]\rho _0`$ (A11)
$`\dot{\rho }_0`$ $`=`$ $`\gamma _R𝒥[1n_1]\rho _1\gamma _L𝒜[1n_1]\rho _0+\gamma _R^{}𝒥[n_1]\rho _1\gamma _L^{}𝒜[n_1]\rho _1`$ (A12)
Here $`n_1=c_1^{}c_1`$ and $`𝒥`$ and $`𝒜`$ are as defined in Eqs. (5), (6). Under the above conditions, we can slave $`\rho _1`$ to $`\rho _0`$ so that
$$(\gamma _R𝒜[1n_1]+\gamma _R^{}𝒜[n_1])\rho _1=(\gamma _L𝒥[1n_1]+\gamma _L^{}𝒥[n_1])\rho _0$$
(A13)
Operating on both sides alternately by $`𝒥[n_1]`$ and $`𝒥[1n_1]`$ it is easy to show that
$`\gamma _R^{}𝒥[n_1]\rho _1`$ $`=`$ $`\gamma _L^{}𝒥[n_1]\rho _0`$ (A14)
$`\gamma _R𝒥[1n_1]\rho _1`$ $`=`$ $`\gamma _L𝒥[1n_1]\rho _0.`$ (A15)
Substituting these into Eq. (A12) yields
$$\dot{\rho }_0=(\gamma _L^{}+\gamma _L)𝒟[n_1]\rho _0$$
(A16)
Since the probability or their being an electron on the SET is very small we can say $`\rho \rho _0`$. Hence we have derived The master equation (3) (without the Hamiltonian term) for the dot alone
$$\dot{\rho }=(2D_0+D_1)𝒟[c_1^{}c_1]\rho .$$
(A17)
Here we have defined
$`D_1`$ $`=`$ $`\gamma _L^{}\gamma _L,`$ (A18)
$`D_0`$ $`=`$ $`\gamma _L.`$ (A19)
Because the SET collector reservoir in the two cases (an electron on the dot and an electron not on the dot) are independent (due to the SET energy shift), the state conditioned on an electron entering the collector is an incoherent mixture of the two possible paths. From quantum trajectory theory , the unnormalized state conditioned on this event is
$$dt\left(\gamma _R𝒥[b(1n_1)]+\gamma _R^{}𝒥[bn_1]\right)R.$$
(A20)
The norm of this state matrix gives the probability for this event, and is equal to the norm of
$$dt\left(\gamma _R𝒥[1n_1]+\gamma _R^{}𝒥[n_1]\right)\rho _1.$$
(A21)
¿From the adiabatic elimination procedure above, this is equal to
$$dt\left(\gamma _L𝒥[1n_1]+\gamma _L^{}𝒥[n_1]\right)\rho .$$
(A22)
This is the unnormalized state $`\stackrel{~}{\rho }_1(t+dt)`$ of the dot alone conditioned on an electron tunneling through the SET. From this it is easy to derive the rate of such tunnelings as
$$\gamma _L1n_1+\gamma _L^{}n_1=D_0+D_1c_1^{}c_1.$$
(A23) |
warning/0002/astro-ph0002531.html | ar5iv | text | # Models of Wave Supported Clumps in Giant Molecular Clouds
## 1 Introduction
Giant molecular cloud complexes (GMCs), the birth places of stars, are typically many tens of parsecs in linear extent and have masses from $`10^4`$ to $`10^6M_{\mathrm{}}`$ and temperatures of 10–30 K (see Hartquist et al. (1998) for a recent review). Observations of CO emission from GMCs (Blitz & Thaddeus (1980); Williams et al. (1995)) show them to be composed of many smaller clumps that are a few parsecs in extent and contain $`\stackrel{<}{}10^3M_{\mathrm{}}`$.
The widths of CO emission lines originating in individual clumps are supersonic and have been attributed to the presence of Alfvén waves having subAlfvénic velocity amplitudes (Arons & Max (1975)). The Alfvén waves contribute to the support of a clump along the direction of the large-scale magnetic field; the damping of the waves affects the degree of support that they provide. An important and well understood mechanism for the damping of linear Alfvén waves in a partially ionized medium is that due to ion-neutral friction which depends on the ionization structure (Kulsrud & Pearce (1969)). Ruffle et al. (1998, hereafter R98) and Hartquist et al. (1993) have emphasised that the dependence of the ionization structure on total visual extinction, $`A_\mathrm{V}`$, should greatly influence the density profiles of clumps if Alfvén waves contribute to their support. To quantify the assertion of Hartquist et al. (1993) and R98, we present in this paper models of plane-parallel, wave-supported GMC clumps like those identified in the work of Williams et al. (1995), who made a detailed analysis of the CO maps of the Rosette Molecular Cloud (RMC), identifying more than 70 clumps. The models that we have constructed are for RMC-type clumps in equilibrium, a restriction justified by the fact that clear spectral signatures of collapse have been found only when much smaller scale features have been resolved (see, e.g., Hartquist et al. (1998)).
We have adopted a WKB description of the wave propagation as did Martin et al. (1997) in their work on wave-supported clumps. Their work differs substantially from ours in that they used an ionization structure appropriate for dark regions. Also, we have considered inwardly rather than outwardly propagating waves, as many of the clumps mapped by Williams et al. (1995) do not contain detected stars and may have no internal means of generating waves. Indeed, the waves may be produced at the surface of a clump by its interaction with an interclump medium.
Other authors have addressed the importance of photoabsorption for the effects that the ionization structure will have upon a clump’s dynamics. These authors have been concerned primarily with dense cores and/or envelopes around them; cores are much smaller-scale objects than the clumps identified in Williams et al. (1995). McKee (1989) addressed the possibility that collapse in a system of dense cores is a self-regulating process due to the ionization of metals such as Magnesium and Sodium by photons emitted by stars formed in the collapse; he was concerned with infall due to ambipolar diffusion of a large-scale magnetic field. Ciolek & Mouschovias (1995) have shown that the large-scale magnetic field can support a photoionized envelope around a dense core for a time that is very long compared to the ambipolar diffusion timescale in the center of the dense core. In contrast to McKee (1989) and Ciolek & Mouschovias (1995), Myers & Lazarian (1998) addressed the effect of photoabsorption on support by waves rather than by the large-scale magnetic field. They stressed that observed infall of dense core envelopes is slower than that expected due to gravitational free-fall and more rapid than collapse due to the reduction by ambipolar diffusion of support by an ordered large-scale magnetic field. They considered collapse of material supported primarily by waves and subjected to an external radiation field. While they made clear comments about the importance of the $`A_\mathrm{V}`$ dependence of the ionization structure for their model, they did not perform any detailed calculations in which a realistic dependence of the ionization fraction on $`A_\mathrm{V}`$ was used.
Several sets of authors have considered nonlinear effects in the dissipation of waves supporting a clump. Gammie & Ostriker (1996) investigated models of plane-parallel clumps and from their “1 2/3-dimensional” models found dissipation times due to nonlinear effects to be longer than the Alfvén crossing times for a fairly large range of parameters. The three dimensional investigations of Mac Low et al. (1998) and Stone et al. (1998) suggest the more restrictive condition that the angular frequency of the longest waves be no more than a few times $`2\pi /t_\mathrm{A}`$ (where $`t_\mathrm{A}`$ is the Alfvén crossing time) in order for the dissipation timescale due to nonlinear damping to be roughly the Alfvén crossing time or more. In most cases addressed in this paper we have restricted our attention to such angular frequencies so that we are justified to lowest order in focusing on only the damping due to ion-neutral friction. It should be noted that the above three dimensional studies of nonlinear effects concerned homogeneous turbulence and did not include ion-neutral damping for a realistic ionization structure. If we are correct in supposing that the waves in clumps are driven externally, then the turbulence is not homogeneous and its nature depends on both the viscous scale set by ion-neutral damping and the exact boundary conditions. The effects of nonlinear damping and multiple dimensions will be considered in subsequent work.
In Sect. 2 we present the equations for the wave energy, the static equilibrium clump structure, and the gravitational field. In Sect. 3 we give a description of the calculations of the ionization structure for various values of the clump density and $`A_\mathrm{V}`$ while Sect. 4 contains details of the models considered here. Finally, in Sect. 5, we present conclusions.
## 2 Equations of Wave Propagation and Static Equilibrium
We consider plane-parallel clumps with $`z=0`$ corresponding to the clump midplane and $`z=z_\mathrm{b}`$ (with $`z_\mathrm{b}`$ defined as positive) corresponding to boundary between the clump and the interclump medium. The large-scale magnetic field is taken to be $`B_0\widehat{z}`$ with $`\widehat{z}`$ normal to the surface of a plane-parallel clump. We study waves of angular frequency $`\omega `$ propagating from $`z=z_\mathrm{b}`$ in the $`\widehat{z}`$ direction.
The ion velocity can be expressed as
$$v_\mathrm{i}=Ve^{i{\scriptscriptstyle k_\mathrm{r}dz}},$$
(1)
where $`V`$ is defined below and $`k_\mathrm{r}`$ is the real component of the complex wave vector. Hartquist & Morfill (1984) used a two-fluid treatment to examine a related problem and showed that for an inwardly propagating linear Alfvén wave
$$\frac{\mathrm{d}}{\mathrm{d}z}\left(k_\mathrm{r}v_\mathrm{i}v_\mathrm{i}^{}\right)=\frac{V_2}{v_{A_\mathrm{i}}^2}v_\mathrm{i}v_\mathrm{i}^{},$$
(2)
where $`v_\mathrm{i}^{}`$ is the complex conjugate of the ion velocity perturbation and
$$V_2\frac{\nu _0\omega ^3\rho _\mathrm{n}^2}{\nu _0^2\rho _\mathrm{i}^2+\omega ^2\rho _\mathrm{n}^2},$$
(3)
where $`\rho _\mathrm{i}`$ is ion mass density and $`\rho _\mathrm{n}`$ is the neutral mass density. The ion Alfvén velocity is given by
$$v_{A_\mathrm{i}}^2=\frac{B_0}{\sqrt{4\pi \rho _\mathrm{i}}}.$$
(4)
The ion-neutral coupling frequency, $`\nu _0\rho _\mathrm{i}`$, is such that the momentum transfer per unit volume per unit time from ions to neutrals is given by $`\nu _0\rho _\mathrm{n}\rho _\mathrm{i}(v_\mathrm{i}v_\mathrm{n})`$ where $`v_\mathrm{n}`$ is the velocity of the neutrals. To a reasonably good approximation, C<sup>+</sup> is the dominant ion and we may take
$$\nu _0\rho _\mathrm{i}2.1\times 10^9\mathrm{sec}^1\left(\frac{n_\mathrm{i}}{1\mathrm{cm}^3}\right)$$
(5)
where $`n_\mathrm{i}`$ is the ion number density (Osterbrock (1961)). If the dominant ion species are very massive, as occurs at large $`A_\mathrm{V}`$ and densities, the constant would approach $`2.3\times 10^9`$s<sup>-1</sup>. However, in this work we ignore this dependence and use $`12m_\mathrm{H}`$ as the mass per ion.
If $`V`$ in Eq. 1 is written as $`e^{{\scriptscriptstyle k_\mathrm{i}dz}}`$, where $`k_\mathrm{i}`$ is the imaginary component of the wave vector, and the linearized version of the induction equation (cf. Eq. 3 of Hartquist & Morfill 1984) is used, Eq. 2 yields
$$\frac{\mathrm{d}}{\mathrm{d}z}\left(\frac{k_\mathrm{r}}{k_\mathrm{r}^2+k_\mathrm{i}^2}bb^{}\right)=\frac{V_2}{v_{A_\mathrm{i}}^2}\frac{bb^{}}{k_\mathrm{r}^2+k_\mathrm{i}^2},$$
(6)
where $`b`$ is the perturbation magnetic field and $`b^{}`$ is its complex conjugate. In the WKB approximation, d$`{}_{}{}^{2}V/`$d$`z^2`$ is taken to be equal to zero, which is equivalent to assuming
$$k_\mathrm{i}^2+\mathrm{d}k_\mathrm{i}/\mathrm{d}zk_\mathrm{r}^2$$
(7)
Then it follows (cf. Eq. 7b of Hartquist & Morfill 1984) that in the WKB approximation
$$k_\mathrm{r}^2\frac{\omega ^2+V_1}{v_{A_\mathrm{i}}^2},$$
(8)
with
$$V_1\frac{\nu _0^2\omega ^2\rho _\mathrm{n}\rho _\mathrm{i}}{\nu _0^2\rho _\mathrm{i}^2+\omega ^2\rho _\mathrm{n}^2}.$$
(9)
As is consistent with the WKB approximation, we assume $`k_\mathrm{i}^2k_\mathrm{r}^2`$ and substitute Eq. 8 into Eq. 6 to find
$$\frac{\mathrm{d}}{\mathrm{d}z}\left(\frac{v_{A_\mathrm{i}}U}{\sqrt{\omega ^2+V_1}}\right)=\frac{V_2}{\omega ^2+V_1}U,$$
(10)
where $`U=bb^{}/16\pi `$ is the time-averaged energy density of the perturbation magnetic field.
We solve Eq. 10 along with the static equilibrium equation
$$c_s^2\frac{\mathrm{d}(\rho _\mathrm{n}+\rho _\mathrm{i})}{\mathrm{d}z}+\frac{\mathrm{d}U}{\mathrm{d}z}=(\rho _\mathrm{n}+\rho _\mathrm{i})g$$
(11)
(Martin et al. (1997)) and the gravitational equation
$$\frac{\mathrm{d}g}{\mathrm{d}z}=4\pi G(\rho _\mathrm{n}+\rho _\mathrm{i}),$$
(12)
where $`c_s`$, $`g`$, and $`G`$ are the isothermal sound speed, the strength of the gravitational field, and the gravitational constant, respectively.
We verify the assumption that $`k_\mathrm{i}^2k_\mathrm{r}^2`$ a posteri by checking that
$$k_\mathrm{i}=\frac{1}{V}\frac{\mathrm{d}V}{\mathrm{d}z}\frac{V_2}{v_{A_\mathrm{i}}\sqrt{V_1}}\frac{V_1}{v_{A_\mathrm{i}}^2}.$$
(13)
Note that if Eq. 13 is satisfied, Eq. 7 is as well.
## 3 Calculations of the Ionization Structure
The ionization structure determined by R98 and presented in their Fig. 1 was calculated on the assumption that, due to shielding of the CO by itself and by H<sub>2</sub>, the rate of CO dissociation by photons of external origin is negligible. For a plane-parallel semi-infinite cloud with constant Hydrogen nucleus number density, $`n_\mathrm{H}=10^3`$cm<sup>-3</sup>, we assume an $`A_\mathrm{V}`$-dependent dissociation rate that results in a CO abundance relative to $`n_\mathrm{H}`$, x(CO), that is in harmony with the measurements shown in Fig. 6 of van Dishoeck (1998). For the $`n_\mathrm{H}=10^3`$cm<sup>-3</sup> model, Table 1 gives x(CO) and the photodissociation rate as a function of $`A_\mathrm{V}`$. Note that the total abundance of carbon nuclei relative to $`n_\mathrm{H}`$ is fixed at $`10^4`$. In the work reported here, we used an ionization fraction that depends on both $`A_\mathrm{V}`$ and $`n_\mathrm{H}`$. Using the $`A_\mathrm{V}`$-dependent CO photodissociation rate from Table 1, we calculated the fractional ionization as a function of $`A_\mathrm{V}`$ for $`n_\mathrm{H}=3\times 10^2,3\times 10^3,`$ and $`10^4`$cm<sup>-3</sup> at various $`A_\mathrm{V}`$ values. A bilinear interpolation in $`A_\mathrm{V}`$ and $`n_\mathrm{H}`$ is used to find the actual ionization fraction used for a given point in the clump. For densities above $`10^4`$cm<sup>-3</sup> we assume that the total ionization fraction, $`\xi \rho _\mathrm{i}/\rho _\mathrm{n}`$, goes as $`n_\mathrm{H}^{1/2}`$, as expected in a dark region. Note that in all models, $`A_\mathrm{V}`$ is always greater than 4 whenever $`n_\mathrm{H}>10^4`$cm<sup>-3</sup>.
Since the depletions in RMC-type clumps are very uncertain, we present results for both depletion cases given in R98. As discussed therein (also see Shalabiea & Greenberg 1995), case A abundances resemble those seen in dark cores with $`A_\mathrm{V}\stackrel{>}{}5`$ while case B, with higher fractional abundances of lower ionization potential elements, is more appropriate for more diffuse clouds.
## 4 Details of the Models
Many RMC-type clumps are not bound by their own gravity and must be confined by interclump media (Bertoldi & McKee (1992)). We shall assume that an interclump medium is sufficiently tenuous that it does not shield the clump from the standard interstellar background radiation field used in the calculation of $`\xi `$ so that $`A_\mathrm{V}=0`$ at $`z=z_\mathrm{b}`$. The material on either side of the interface between a clump and the interclump medium is in two distinct phases, and we may assume that at its outer boundary a clump has a substantial density; we take $`n({}_{}{}^{}\mathrm{H}_{2}^{})=n_\mathrm{H}/2`$ everywhere in the clump and at $`z=z_\mathrm{b}`$ set $`n({}_{}{}^{}\mathrm{H}_{2}^{})=n_\mathrm{b}`$. Waves may exist in the interclump medium and be partially transmitted into the clump, or, as mentioned earlier, may be generated near the interface by the interaction between the clump and the interclump medium. Consequently, we assume that the magnitude of the amplitude of the perturbation magnetic field, $`\delta B=\sqrt{bb^{}}`$, is, at $`z=z_\mathrm{b}`$, a substantial fraction, $`f_\mathrm{b}`$, of the large-scale field, $`B_0`$.
In all models we have taken the mean mass per neutral particle, $`\mu _\mathrm{n}`$, to be 2.3 amu, corresponding to $`14\%`$ of the neutral particles being He and $`86\%`$ being $`{}_{}{}^{}\mathrm{H}_{2}^{}`$. As is consistent with data given by Savage & Mathis (1979), we have assumed
$$A_\mathrm{V}=\frac{N_\mathrm{H}}{1.9\times 10^{21}\mathrm{cm}^2},$$
(14)
where $`N_\mathrm{H}`$ is the column density of Hydrogen nuclei. Also, the temperature throughout a clump was taken to be 20 K.
For a given setup, an initial boundary value of $`g(z=z_\mathrm{b})g_\mathrm{b}`$, was selected and Eqs. 1011, and 12 were numerically integrated using an adaptive Gear algorithm (Gear (1971)). The value of $`g_\mathrm{b}`$ was changed by iteration until the inner boundary condition $`g(z=0)=0`$ was satisfied. Note that as the total visual extinction (or, equivalently for a plane-parallel cloud, the column density) for a cloud is increased, the clump reaches a maximum size, $`z_{\mathrm{max}}`$, and starts shrinking as a larger and larger thermal pressure (and therefore density) is required to balance gravity. Thus, as long as $`z_\mathrm{b}<z_{\mathrm{max}}`$, there are two values of $`g_\mathrm{b}`$ for each setup that satisfy the boundary conditions. One solution has a relatively flat density profile and a small total visual extinction while the other solution is more centrally condensed. We deal here exclusively with the latter solutions.
We have considered other models but present full results only for models which have a total edge-to-center visual extinction of 5 magnitudes since, as discussed in R98, it is in the region of $`A_\mathrm{V}`$ of a few that clumps appear to begin to contain detected stars, while many dense cores may have $`A_\mathrm{V}\stackrel{<}{}5`$ (McKee (1989)). For our canonical model (Model 1) we require a velocity amplitude, $`V`$, of 2 km sec<sup>-1</sup> and an Alfvén speed, $`v_\mathrm{A}`$, of 3 km sec<sup>-1</sup> at $`A_\mathrm{V}=2`$. Since the concentration of CO at $`A_\mathrm{V}\stackrel{<}{}2`$ is very low and measurements of GMC clump CO profiles have a width of $`2`$ km sec<sup>-1</sup> (Williams et al. (1995)), observations require such velocities to exist well within the cloud. Also, we use a wave frequency, $`\omega `$, of $`2\times 10^{12}`$ sec<sup>-1</sup>; this results in relatively strong neutral-ion coupling while keeping the wavelength of the perturbing wave less than the size of the cloud. In other words,
$$\nu _0\xi \omega \stackrel{<}{}\frac{2\pi }{dz/v_\mathrm{A}}.$$
(15)
Within the above constraints, we find for our canonical model the solution requiring the smallest value of $`B_0`$, and, since the velocity amplitude of a linear Alfvén wave is thought to be comparable to but less than the Alfvén speed, the largest value of $`f_\mathrm{b}`$. However, throughout the clump, $`Vv_\mathrm{A}`$, consistent with our assumption of linear Alfvén waves.
We present the results of 5 models. Model 1 is for the above canonical parameters with R98’s depletion case A while Model 2 is for case B. Model 3 is the same as Model 1 but with $`\omega =1\times 10^{12}`$ to illustrate the effect of a scenario with roughly maximum ion-neutral coupling. Model 4 is as Model 1 but with an ionization profile given by $`\xi =3\times 10^{16}\rho _{\mathrm{n}}^{}{}_{}{}^{1/2}`$ (see, e.g., McKee 1989 and Myers & Lazarian 1998). Thus, Model 4 is for a dark region surrounded by interclump material. Finally, for completeness, we present a model with no turbulence; Model 5 is identical to Model 1 except with $`f_\mathrm{b}=0`$. A summary of the parameters of the 5 models is given in Table 2. For all models, $`n_\mathrm{b}=375{}_{}{}^{}\mathrm{H}_{2}^{}`$ cm<sup>-3</sup> and $`B_0=135\mu `$G.
## 5 Results and Conclusions
In Fig. 1 we present density as a function of visual extinction for each of the 5 Models. It is clear that the newer ionization profiles used in Models 1–3 result in less condensed, more extended clumps than the profile used in Model 4. In fact, $`n_\mathrm{H}`$ is roughly proportional to $`1/z`$ in Models 1–3 while $`n_\mathrm{H}`$ goes roughly as $`1/z^2`$ in Model 4. In Fig. 2 we show plots of $`\delta B`$ versus $`A_\mathrm{V}`$ for Models 1–4. Except for Model 4, which has $`k_\mathrm{i}/k_\mathrm{r}`$ approaching $`0.5`$ at the clump boundary so that Eq. 13 is not satisfactorily satisfied, the perturbing field obeys flux conservation near the surface of the clump. In Model 1, in the central region of the clump dissipation is rapid enough that $`\delta B`$ begins to decrease. In order to compensate for the loss of support, the equilibrium solution requires a complementary increase in the density, as can be seen in Fig. 1. On the other hand, in Model 4, the turbulence is dissipated much nearer to the cloud boundary, thus requiring a steeper overall density profile. The higher ionization fractions of the case B depletions result in very little dissipation even in the center of the clump for Model 2. Note that even though observations (Williams et al. (1995)) suggest that the temperature of RMC-type clumps is closer to 10 K rather than the 20 K used here, thermal support is insignificant except in the centre of Model 1, so the effect of a lower temperature on the models would be merely to enhance slightly any central condensations.
The ionization profiles used in the Models are shown in Fig. 3. The ionization profiles described in Sec. 3 result in $`\xi `$ for Models 1–3 being more than 50 times greater near the surface of the clump than in Model 4. However, in the center of the clump the ionization fraction drops, resulting in more dissipation. Again, this leads to clumps that are overall more diffuse but with small condensed cores. Clearly, clumps with $`A_\mathrm{V}\stackrel{>}{}5`$ will have distinct central condensations with $`n/n_\mathrm{b}\stackrel{>}{}100`$ and central fractional ionizations of $`\stackrel{<}{}5\times 10^7`$. Though dense cores may be formed during the fairly rapid collapse (as envisaged by Fielder & Mouschovias 1993) of more extended objects (i.e. RMC-like clumps) that become unstable, even in our equilibrium models we find central cores having densities and fractional ionizations similar to those measured for dense cores and their envelopes (Williams et al. (1995); Williams et al. (1998); Bergin et al. (1999)).
Fig. 4 shows the flux of magnetic energy through the clumps for Models 1–4. Near the clump center, Model 1 is nearly thermally supported due to the dissipation of the turbulence. The higher ionization fraction for the case B depletions used in Model 2 results in less dissipation and thus more turbulent support for the clump. Consequently, as can be seen in Fig. 1, Model 2 has no central condensation and is more extended. Unfortunately, we can only speculate about how the depletions of Sulphur, metals, and some other species behave in RMC-like clumps (Ruffle et al. (1999)). Thus, cases A and B are merely representative; as can be clearly seen in the figures, the clump profiles are very sensitive to the choice of abundances and the subsequent fractional ionizations. In addition, compared to Model 1, the stronger ion-neutral coupling in Model 3 results in less dissipation and subsequently the clump has little central condensation, as expected. Note however that for Model 3 the lower limit in Eq. 15 is not adequately satisfied.
Fig. 4 also shows the effect of external wave generation. If the fractional ionization is too low, as in Model 4, dissipation occurs close to the surface of the clump. Conversely, if the fractional ionization is too large, as in Model 2, significant dissipation occurs only at the clump’s very centre. Both extremes produce density profiles which lack a central condensation. Note that if our externally generated wave model is correct, one should not see turbulence within a condensed core if there is no turbulence in its surrounding envelope.
In Fig. 5 we present curves which map the visual extinction to the spatial extent of the clumps. Clouds with larger extents which match the observed 2–3 pc size of RMC-type clumps (Williams et al. (1995)) cannot be reproduced within the constraints given in Sect. 4. However, observations generally measure the largest linear extent of a clump. Thus, since the waves only support the model clumps parallel to the large-scale field, it is not surprising that the model sizes given here are less than the observed sizes.
Similarly, the models require high boundary densities and magnetic field strengths in order for the Alfvén speed and wave velocity amplitude at visual extinctions where CO is abundant to be large enough to be compatible with observed linewidths. For Model 1, $`n_\mathrm{b}=375{}_{}{}^{}\mathrm{H}_{2}^{}`$ cm<sup>-3</sup>. This is rather higher than the typical value of $`n({}_{}{}^{}\mathrm{H}_{2}^{})=220`$ cm<sup>-3</sup> given by Williams et al. (1995) for RMC-type clumps but, given the uncertainties, it is within a reasonable range of the Williams et al. (1995) value. For Model 1, $`B_0=135\mu `$G, significantly higher than the value of $`30\mu `$G suggested by observations (Heiles (1987)) and expected from robust theoretical arguments (Mouschovias (1987)). In order to determine whether the values of $`n_\mathrm{b}`$ and $`B_0`$ could be lower and still allow model properties to be consistent with observed linewidths, we constructed models for clumps with total edge-to-center extinctions of 3 magnitudes. The model giving $`V=2`$ km sec<sup>-1</sup> and $`v_\mathrm{A}=3`$ km sec<sup>-1</sup> at $`A_\mathrm{V}=2`$ had $`n_\mathrm{b}=325{}_{}{}^{}\mathrm{H}_{2}^{}`$ cm<sup>-3</sup>, $`B_0=105\mu `$G, $`f_\mathrm{b}=0.49`$, and $`z_{\mathrm{max}}=0.565`$ pc; although $`n_\mathrm{b}`$ and $`B_0`$ were smaller and $`z_{max}`$ larger, the agreement with observations is nonetheless poor.
Thus, the next step in the modelling of clumps in which wave support is important is the inclusion of wave support in models analogous to the axisymmetric models of magnetically and thermally suported clumps described in classic papers by Mouschovias (1976a,1976b). It is possible that the inclusion of magnetic tension, as well as pressure, will allow the reduction of $`B_0`$ to a value more like that expected and the construction of models of clumps having larger linear extents. |
warning/0002/astro-ph0002302.html | ar5iv | text | # Spirals and the size of the disk in EX Dra
## 1 Introduction
Spiral shocks in accretion disks have been predicted from numerical simulations (Sawada et al. 1986, 1987, Różyczka & Spruit 1993, Yukawa et al. 1997), and analytic considerations (Spruit 1987, Spruit et al. 1987, Larson 1990). They are excited by the tidal field of the secondary if the disk extends far enough into the Roche lobe and can result in two prominent spiral arms. If shock dissipation is the main mechanism damping the wave, it extends over the entire disk and causes accretion at an effective $`\alpha `$-value of $`0.01(H/r)^{3/2}`$ (Spruit 1987, Larson 1990, Godon 1997).
The first observational evidence for shock waves in accretion disks of cataclysmic variables (CVs) was the detection of a clear two-armed structure in the disk of IP Peg during rise to outburst (Steeghs et al. 1997). The spiral pattern, interpreted as evidence for shock waves, has also been seen during outburst maximum (Harlaftis et al. 1999) and early decline of outburst (Morales-Rueda et al. 2000).
At the temperatures expected from dwarf nova disks models, whether in outburst or in quiescence, the predicted spirals are tightly wound and would be hard to detect observationally (Bunk et al. 1990), so that their presence in the observations is somewhat unexpected. Spirals this strong are most naturally explained if the disk temporarily extends rather far into the primary Roche lobe, so that the tidal force of the secondary causes a strong disturbance. A strong non-axisymmetric disturbance, however, would also cause the gas to loose angular momentum quickly (transfered to the secondary), so that the disk would shrink to a smaller radius where the tidal force is weaker. Spirals in disks of CVs would then be understandable if they are a temporary phenomenon, perhaps restricted to outbursts. To test this, more observations of different systems at sufficient spectral resolution and signal-to-noise are needed (Steeghs & Stehle 1999).
With high quality spectroscopic studies of different CVs it should also be possible to answer the question if spiral shocks in CV accretion disks are a common phenomenon. Systems suitable for this purpose would be bright and have frequent outbursts, such as SS Cyg and EX Dra.
EX Dra is a double–eclipsing dwarf nova with a 5-hour orbit (Barwig et al. 1993, Billington et al. 1996, system parameters by Fiedler et al. 1997). There is suggestive evidence for asymmetric structures in the He i Doppler map reconstructed from outburst data taken in 1993 (Joergens et al. 2000). From eclipse maps obtained at various stages in the outburst cycle Baptista & Catalan (1999) claim that spiral waves form at the early stages of an outburst. We report in this paper on observations at high spectral and temporal resolution during an outburst in 1996.
## 2 Observations and reduction
EX Dra was observed on the nights of July 27 and 28 1996, with the ISIS spectrograph on the 4.2 m William Herschel Telescope. The red ISIS arm was equipped with the TEK 5 CCD, the blue arm with TEK 1 CCD, covering the wavelength ranges 6375-6778 Å and 4585-4993 Å respectively, at a dispersion of 0.4 Å/pixel. The spectra were observed during seeing of 1.4 – 1.7 arc sec, using an exposure time of 60s. The spectra were optimally-extracted, including the elimination of cosmic ray hits but not the correction of Pixel-to-pixel variations of the detector sensitivity. Slit losses due to variable atmospheric conditions were corrected using a faint comparison star on the slit. The red spectra, but not the blue spectra, were flux calibrated.
Our spectra were taken roughly in the middle of an outburst. Photometric observation at the Wendelstein observatory showed that the system was already in outburst on 26 July, and in quiescence on 4 August (Barwig, private comm.). VSNET (1998) records show that the system was in quiescence on 24 July and in outburst one day later and in quiescence again on 1 August. The outburst therefore started on 25 August, three days before our observations.
## 3 Interpretation of the Doppler maps
Doppler maps were computed from the phase-folded spectra of 28 July, using the IDL-based fast-maximum entropy package described by Spruit (1998)<sup>1</sup><sup>1</sup>1available at http:www.mpa-garching.mpg.de/$``$henk. For further details on Doppler tomography see Marsh & Horne (1988). The eclipsed part of the data (phases -0.12 to 0.12) has been excluded from the reconstruction process, and does not affect the maps produced. For information, however, these parts of the data are included in the spectra shown in Fig. 1. They were used separately to derive disk sizes in the lines and the continuum.
Fig. 1 shows the results for the four strongest lines in the spectra, H<sub>α</sub>, H<sub>β</sub>, He i$`\lambda 6678`$ and He ii$`\lambda 4686`$. The phase-folded spectra are shown in the left column, the corresponding Doppler maps in the right and in the middle the spectra reconstructed from the maps. The theoretical trajectory of the mass transfering stream has been plotted in the He ii image, together with the Keplerian velocity along the stream path (cp. Marsh & Horne 1988). Bars connecting the two arcs indicate correspondence in physical space, and are annotated with radius r/a and the azimuth $`\mathrm{\Phi }`$ relative to the primary. The system parameters and the ephemeris are from Fiedler et al. (1997).
The He i image shows the spirals clearest: asymmetric disk emission is concentrated in two arms in the first and third quadrant. We interpret this deviation from a Keplerian disk as indication of the presence of spiral shock waves in the accretion disk. The image looks quite similar to the He i map of IP Peg published by Steeghs et al. (1997). As in the case of IP Peg, the arms do not follow a circle centered on the white dwarf, but are somewhat elongated along the V<sub>y</sub>-axis. This is the pattern expected from spiral shocks (Steeghs & Stehle 1999). The asymmetry of the spirals in the EX Dra disk is smaller than in the IP Peg observations (Steeghs et al. 1997, Harlaftis et al. 1999, Morales-Rueda et al. 2000), perhaps indicating that the shocks are weaker. The spiral arm in the upper right is stronger in intensity as well as in asymmetry than that one in the lower left. This pattern is also seen in the Doppler maps of IP Peg. As in the case of IP Peg, the Balmer lines show similar, but less clearly defined structures.
The Doppler images, with the exception of He ii, show strong emission from the secondary star, also visible during the outburst in 1993 (Joergens et al. 2000). This indicates heating of the secondary by radiation from the inner disk. The temperatures are obviously not high enough to excite the He ii line. The He ii line also differs from the other lines by the presence of a prominent emission patch at the theoretical gas stream trajectory. The center of this hot spot emission is somewhat below the trajectory, consistent with quiescence observations of EX Dra (Billington et al. 1996, Joergens et al. 2000).
## 4 The disk radius
The large width of the He ii line is as expected for a high excitation line produced near the center of a disk in outburst. The eclipse of the He ii line is visible in the spectrum up to a velocity of about 1500 km/s. The orbital velocity at the disk edge is about 4.3 times lower than this. Assuming Keplerian rotation, the disk region eclipsed at 1500 km/s is about 20 times smaller in radius than the outer radius of the disk. At our phase resolution, the 1500 km/s emission is thus point-like, and can be used to measure the width of the secondary’s shadow on the orbital plane, $`\mathrm{\Delta }\varphi =0.08\pm 0.005`$. From the position of the hot spot in the He ii map, we can estimate a disk size of $`r/a=0.32\pm 0.04`$.
The eclipse of the Balmer and He i lines shows that the blue disk emission reappears (at V$`350`$km/s) at the same phase as the last red emission disappears, within the measurement uncertainty. The size of the disk as seen in these lines is thus, by a coincidence, nearly as wide as the occulting shadow of the secondary. Hence the disk radius as seen in the Balmer lines is $`r_{\mathrm{dB}}/a2\pi \mathrm{\Delta }\varphi /2=0.25\pm 0.02`$.
We can compare this with the disk size as measured from the eclipse in the continuum. The light curve in the continuum between 6400 and 6500Å, extracted from our spectra, is shown in Fig. 2. The width of the eclipse is $`w=0.17\pm 0.005`$ orbits. With $`\mathrm{\Delta }\varphi =0.08`$, this implies a disk radius $`r_\mathrm{d}/a=0.31\pm 0.01`$.
Fiedler et al. (1997) find a value of 0.11 for the eclipse width of the white dwarf, from analysis of a mean blue light curve in quiescence. The disk as seen in the continuum thus appears to be larger in outburst than in quiescence, as expected. The eclipse width values of 0.17 in outburst and 0.11 in quiescence are also compatible with the eclipse light curves of Baptista & Catalan (1999). These authors measured the eclipses of EX Dra at various stages in the outburst cycle.
## 5 Discussion and conclusions
We find evidence for spiral structures in the outburst accretion disk of EX Dra similar to those found in IP Peg by Steeghs et al. (1997). The pattern of intensity and velocity perturbations agrees with that predicted from numerical simulations of spiral shock waves (Steeghs & Stehle 1999). It is best seen in the He i line, somewhat less clearly in the Balmer and He ii lines. In EX Dra the pattern appears somewhat less clearly and asymmetric than in IP Peg.
In particular, it is less clear in the He ii line, suggesting lower temperatures and shock strengths of the spirals of EX Dra. Possibly the observability of spirals in the He ii map is hampered by strong hot spot emission visible in this line.
We derive a disk radius of $`r/a=0.31\pm 0.01`$ from the red continuum eclipse light curve. From numerical simulations, Steeghs & Stehle find that disk sizes $`r_\mathrm{d}/a=\mathrm{\hspace{0.17em}0.3}0.4`$ are needed to excite spirals that are strong enough to generate an observable pattern in the spectra. (Transformation from $`r_\mathrm{d}/a`$ to $`r_\mathrm{d}/R_{L_1}`$ is given by R$`{}_{L_1}{}^{}=0.53a`$ for q=0.75, cp. Plavec & Kratochvil 1964.) Since the tidal force is a very steep function of $`r/a`$, the spirals rapidly become weak at smaller disk sizes. The disk size we find here in EX Dra is at the lower limit of the required size.
Further evidence for the size of the disk in EX Dra in outburst was obtained by Baptista & Catalan (1999). Radial intensity distributions presented there show disk radii of 0.30 $`a`$ in quiescence and 0.49 $`a`$ in outburst. The authors see hints of spirals in their eclipse maps, during the early outburst stages. This may be compared with the observations presented here, which show spirals and a disk size of 0.31 $`a`$ three days after the beginning of an outburst.
The eclipses of the Balmer lines in our spectra yield significantly smaller disks sizes than the continuum eclipse, $`r_{\mathrm{dB}}/a=0.25\pm 0.02`$. Since it is known that the line emission is produced primarily in the outer parts of the disk (e.g. Rutten et al. 1994), one might have expected the disk as seen in the lines to be larger, if anything, than in the continuum.
A possible resolution of this conflict may lie in the optical depth effects affecting the emission lines in systems seen at high inclination. The low central intensity of the Balmer lines in high-inclination CVs (often below the continuum) shows that such effects are strong. As shown by Horne & Marsh (1986), the effects are strongest for lines of sight parallel and perpendicular to the orbital motion, leading to reduced line emission from these directions compared to intermediate lines of sight (near $`45^{}`$ to the orbit). The bias towards intermediate angles will give the appearance of a somewhat smaller disk size. It is still to be determined if this suggestion also works quantitatively.
###### Acknowledgements.
We thank Dr. Heinz Barwig for providing the photometric information on EX Dra. We thank the anonymous referee for the comments, which helped to improve the presentation of the data and to correct an error. This work was done in the context of the research network ‘Accretion onto black holes, compact objects and protostars’ (TMR Grant ERB-FMRX-CT98-0195). The Isaac Newton Telescope is operated on the island of La Palma by the Isaac Newton Group of Telescopes in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias. |
warning/0002/cond-mat0002162.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In the first paper of this series we have shown how to diagonalize the Hamiltonian of the XX–model with boundaries
$$H=\frac{1}{2}\underset{j=1}{\overset{L1}{}}\left[\sigma _j^+\sigma _{j+1}^{}+\sigma _j^{}\sigma _{j+1}^+\right]+\frac{1}{\sqrt{8}}\left[\alpha _{}\sigma _1^{}+\alpha _+\sigma _1^++\alpha _z\sigma _1^z+\beta _+\sigma _L^++\beta _{}\sigma _L^{}+\beta _z\sigma _L^z\right]$$
(1)
by introducing an auxiliary Hamiltonian
$`H_{\mathrm{long}}={\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{L1}{}}}\left[\sigma _j^+\sigma _{j+1}^{}+\sigma _j^{}\sigma _{j+1}^+\right]`$
$`+{\displaystyle \frac{1}{\sqrt{8}}}\left[\alpha _{}\sigma _0^x\sigma _1^{}+\alpha _+\sigma _0^x\sigma _1^++\alpha _z\sigma _1^z+\beta _+\sigma _L^+\sigma _{L+1}^x+\beta _{}\sigma _L^{}\sigma _{L+1}^x+\beta _z\sigma _L^z\right]`$ (2)
which in turn may be diagonalized in terms of free fermions. The parameters $`\alpha _\pm ,\beta _\pm ,\alpha _z`$ and $`\beta _z`$ are arbitrary complex numbers. Note that $`H_{\mathrm{long}}`$ commutes with $`\sigma _0^x`$ and $`\sigma _{L+1}^x`$. Hence the spectrum of $`H_{\mathrm{long}}`$ decomposes into four sectors $`(+,+),(+,),(,),(,+)`$ corresponding to the eigenvalues $`\pm 1`$ of $`\sigma _0^x`$ and $`\sigma _{L+1}^x`$. The spectrum of $`H`$ is obtained by projecting onto the $`(+,+)`$–sector.
While in our focus was on the diagonalization of the finite chain, here we are going to examine the asymptotic behaviour of the energy levels for large values of $`L`$. In the case of conformal invariance this information is usually encoded in the partition function, which we define by
$$𝒵=\underset{L\mathrm{}}{lim}\mathrm{tr}z^{\frac{L}{\xi }(He_{\mathrm{}}Lf_{\mathrm{}})}$$
(3)
where $`z<1`$ and $`e_{\mathrm{}}`$ and $`f_{\mathrm{}}`$ denote the free bulk and the free surface energy respectively. $`\xi `$ may be considered as a normalization constant here. We will also consider the partition function for $`H_{\mathrm{long}}`$.
In the case of periodic boundary conditions the continuum limit of the XXZ–chain can be described by a free boson field $`\mathrm{\Phi }`$ (see e.g. ), with action
$$S=\frac{1}{2}𝑑x_1𝑑x_2\left[\left(_1\mathrm{\Phi }\right)^2+\left(_2\mathrm{\Phi }\right)^2\right]$$
(4)
being compactified on a circle of radius $`r`$, i.e. we identify $`\mathrm{\Phi }`$ with $`\mathrm{\Phi }+2\pi r`$ (we follow the notation of ). The radius of compactification $`r`$ is related to the value of the anisotropy in the XXZ–model. At the free fermion-point the anisotropy is zero and we have $`r=1/\sqrt{4\pi }`$.
Consider the boson field on a cylinder of length $`L`$ and circumference $`N`$ imposing periodic boundary conditions in the direction of the cylinders length. Quantizing the field with time in $`N`$-direction yields the partition function
$$𝒵_{\mathrm{periodic}}=\frac{1}{\eta (q)^2}\underset{m,n}{}q^{\frac{1}{2}\left(\frac{m^2}{2\pi r^2}+2\pi r^2n^2\right)}$$
(5)
where $`q=\mathrm{}^{2\pi N/L}`$ and
$$\eta (q)=q^{\frac{1}{24}}\underset{n=1}{\overset{\mathrm{}}{}}(1q^n).$$
(6)
This expression coincides with the partition function for the periodic XXZ-chain, which has been studied extensively in . The partition function for $`H_{\mathrm{long}}`$ in the periodic case is just $`4𝒵_{\mathrm{periodic}}`$.
Suppose now the case of open boundaries at both ends of the cylinder. There are two different boundary conditions preserving conformal invariance, i.e. the Dirichlet and the von Neumann (see e.g. ). This yields three different partition functions corresponding to the various possible combinations of the two boundary conditions. The respective partition functions can be found in . Imposing Dirichlet boundary conditions at both ends of the cylinder results in
$$𝒵_{\mathrm{DD}}(\mathrm{\Delta })=\frac{1}{\eta (q)}\underset{n}{}q^{\frac{1}{2}(2\sqrt{\pi }rn+\mathrm{\Delta })^2}$$
(7)
where $`q=\mathrm{}^{\pi N/L}`$ and $`\mathrm{\Delta }=(\mathrm{\Phi }_0\mathrm{\Phi }_L)/\sqrt{\pi }`$. By $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_L`$ we denote the values of the boson field at the boundaries. This type of partition function has also been obtained for the open XXZ-chain with diagonal, hermitian boundary fields .
The Neumann-Neumann partition function is given by
$$𝒵_{\mathrm{NN}}(\mathrm{\Delta })=\frac{1}{\eta (q)}\underset{m}{}q^{2(m/(2\sqrt{\pi }r)+\mathrm{\Delta }/2)^2}$$
(8)
where $`\mathrm{\Delta }=(\stackrel{~}{\mathrm{\Phi }}_0\stackrel{~}{\mathrm{\Phi }}_L)/\sqrt{\pi }`$ and $`\stackrel{~}{\mathrm{\Phi }}_0,\stackrel{~}{\mathrm{\Phi }}_L`$ denote the values of the dual field $`\stackrel{~}{\mathrm{\Phi }}`$ at the boundaries . The Dirichlet-Neumann boundary condition yields
$$𝒵_{\mathrm{DN}}=\frac{1}{2\eta (q)}\underset{k}{}q^{\frac{1}{4}(k1/2)^2}.$$
(9)
In this case no free parameter appears.
Which boundary conditions of the Hamiltonian $`H`$ lead to the partition functions (8) and (9) is yet not known. We will close this gap. We are going to see that, as long as the Hamiltonian $`H`$ is hermitian, the partition function of the chain corresponds to one of the three boson partition functions for open boundaries. However, this may also be the case for certain non-hermitian boundary terms. See section 2 for a survey over the boundary conditions which we will consider.
In it has been argued that non-diagonal boundary terms correspond to von Neumann boundary conditions for the boson field. Our results will show that this assumption is correct in the hermitian case, but not necessarily for non-hermitian boundaries.
In this paper we also study the asymmetric XX–chain with boundaries (see equation (3)). This includes the quantum chain with Dzyaloshinsky-Moriya interactions, which was already studied in for periodic boundary conditions. The partition function for real values of $`p`$ and $`q`$ in (3) and periodic boundary conditions was obtained in . The special case $`p=1`$ and $`q=0`$ with boundaries has already been studied in .
This paper is organized as follows: We will start with a summary of our results for the symmetric chain with hermitian and non-hermitian boundaries in section 2. The asymmetric XX-chain will be discussed in section 3. We will give our conclusions in section 4. The appendix is dedicated to the computation of the fermion energies which yield the energy gaps of the chain.
## 2 Summary of results
The structure of our results is most simply encoded in terms of new parameters $`F,C,G,K,J`$, where
$`F=4\alpha _{}\alpha _+\beta _+\beta _{}`$
$`C=\alpha _{}^2\beta _+^2+\alpha _+^2\beta _{}^2`$
$`G=2\alpha _{}\alpha _+(1+2\beta _z^2)+2\beta _{}\beta _+(1+2\alpha _z^2)`$
$`K=2(\alpha _{}\alpha _++\beta _{}\beta _+)+(1+2\alpha _z^2)(1+2\beta _z^2)`$
$`J=8\alpha _z\beta _z4\alpha _z^2\beta _z^2+2(\alpha _{}\alpha _++\beta _{}\beta _+)+2\alpha _z^2+2\beta _z^21.`$ (10)
The six different situations we considered are given in table 1. This table also indicates which kind of partition function we obtained for $`H`$. For hermitian boundaries the cases in this table cover the whole parameter space. This is not the case for non-hermitian boundary terms.
For the cases in table 1 we computed analytically the energy gaps in the leading order as shown in appendix A. The exact ground-state energies for certain cases being listed in table 2 have already been obtained in .
However, the information of the energy gaps is already sufficient to obtain the partition functions for $`H_{\mathrm{long}}`$ upto a factor $`z^{\frac{1}{24}+h_{\mathrm{min}}}`$, where $`h_{\mathrm{min}}`$ denotes the lowest highest weight appearing in the representation of the Virasoro algebra. This information has to be extracted from the ground-state energies
$$E_0=e_{\mathrm{}}L+f_{\mathrm{}}\frac{1}{L}\left(\frac{1}{24}h_{\mathrm{min}}\right)+\mathrm{}$$
(11)
The value of the free bulk-energy is already known from the periodic chain, i.e. $`e_{\mathrm{}}=\frac{1}{\pi }`$ . The values of the free surface energies $`f_{\mathrm{}}`$ and of $`h_{\mathrm{min}}`$ have been obtained for the cases given in table 2 expanding the exact expressions for the finite chain given in .
The results confirm the partition functions we will give for $`H_{\mathrm{long}}`$ in the following. For the cases not being listed in table 2 we checked our expressions numerically. The partition functions for $`H`$ have been obtained afterwards by projecting onto the $`(+,+)`$-sector (see for details on the projection mechanism). We are now going to discuss the cases given in table 1 in detail.
### 2.1 Neumann-Neumann boundary conditions
The condition $`F0`$ simply implies that all non-diagonal boundary terms are present. The diagonal terms may be arbitrary. It is quite instructive to introduce a new parametrisation of the boundary parameters $`\alpha _\pm `$ and $`\beta _\pm `$ as follows:
$$\alpha _+=R_\alpha \mathrm{}^{\mathrm{}\pi \phi }\alpha _{}=R_\alpha \mathrm{}^{\mathrm{}\pi \phi }\beta _+=R_\beta \mathrm{}^{\mathrm{}\pi (\chi +\phi )}\beta _{}=R_\beta \mathrm{}^{\mathrm{}\pi (\chi +\phi )}$$
(12)
where $`R_\alpha ,R_\beta ^+,\chi ,\phi ,1<\chi 1`$. The value of $`\mathrm{\Delta }`$ which enters the partition function (8) is then given by
$$\mathrm{\Delta }=\{\begin{array}{cc}\chi \hfill & \text{for }L\text{ odd}\hfill \\ \chi +1\hfill & \text{for }L\text{ even}\hfill \end{array}.$$
(13)
The partition function for $`H_{\mathrm{long}}`$ is given by
$$𝒵_{\mathrm{long}}=2𝒵_{NN}(\chi )+2𝒵_{NN}(\chi +1)$$
(14)
for even and odd values of $`L`$. Note that the parametrization (12) does not work for non-hermitian boundaries. In this case we may define the parameter $`\chi `$ via
$$\chi =\frac{\mathrm{}}{2\pi }\mathrm{ln}\left(\frac{\alpha _{}\beta _+}{\alpha _+\beta _{}}\right)$$
(15)
where the logarithm is taken in a way such that $`0|\text{Re}(\chi )|\frac{1}{2}`$ for Re$`(\alpha _+\beta _{}+\alpha _{}\beta _+)>0`$ and $`1|\text{Re}(\chi )|\frac{1}{2}`$ for Re$`(\alpha _+\beta _{}+\alpha _{}\beta _+)<0`$. This rule is consistent with the parametrisation (12), which implies Re$`(\alpha _+\beta _{}+\alpha _{}\beta _+)=2R_\alpha R_\beta \mathrm{cos}(\chi \pi )`$. However, we cannot perform the projection onto the $`(+,+)`$-sector in general if the boundaries are non-hermitian. In this case our result for the partition function of $`H`$ is restricted to the cases listed in table 2 which satisfy the conditions $`\alpha _z=\beta _z=0`$ or $`\alpha _+=\alpha _{},\beta _+=\beta _{}`$ . The result for $`H_{\mathrm{long}}`$ is valid in general.
### 2.2 Dirichlet-Neumann boundary conditions
The conditions for the Dirichlet-Neumann case yield
$$\beta _\pm =0\alpha _\pm 0\beta _z\pm \frac{\mathrm{}}{\sqrt{2}}\text{or}\beta _\pm 0\alpha _\pm =0\alpha _z\pm \frac{\mathrm{}}{\sqrt{2}}.$$
(16)
The value of $`\alpha _z`$ respectively of $`\beta _z`$ may be arbitrary. The partition function for $`H_{\mathrm{long}}`$ is just $`𝒵_{\mathrm{long}}=4𝒵_{DN}`$ for $`L`$ even and odd. The case of one non-diagonal boundary at the end of a semi-infinite XX–chain has already been considered in . The author has shown, that the model can be decoupled into two Ising models with different boundary conditions. One of them being subject to a free boundary condition at the end, the other one being subject to the fixed boundary condition. The partition functions for the Ising model with free and mixed boundary conditions are well known . Taking the product of them also results in (9).
### 2.3 Dirichlet-Dirichlet boundary conditions
The Dirichlet-Dirichlet partition function is obtained for two different types of boundary terms. The first type is given by
$$\beta _z\pm \frac{\mathrm{}}{\sqrt{2}}\alpha _z\pm \frac{\mathrm{}}{\sqrt{2}}\beta _{}=\alpha _{}=0.$$
(17)
where instead of the last equation we might have also chosen $`\beta _+=\alpha _+=0`$. However, for this type of boundaries the non-diagonal boundary terms have no influence on the energies even on the finite chain. The second type is given by
$$\beta _\pm =0\alpha _\pm 0\beta _z=\pm \frac{\mathrm{}}{\sqrt{2}}\text{or}\beta _\pm 0\alpha _\pm =0\alpha _z=\pm \frac{\mathrm{}}{\sqrt{2}}$$
(18)
where the values of $`\alpha _z`$ respectively $`\beta _z`$ may be chosen arbitrarily. Note that for this second type the boundary terms are non-hermitian. For both types of boundaries the value of $`\mathrm{\Delta }`$ which enters (7) is given by
$$\mathrm{\Delta }=\frac{1}{2\pi }\text{arccos}\left((1)^{L+1}\frac{J}{K}\right).$$
(19)
The partition function for $`H_{\mathrm{long}}`$ is $`𝒵_{\mathrm{long}}=4𝒵_{DD}(\mathrm{\Delta })`$ for even and odd $`L`$.
### 2.4 $`F=C=G=K=J=0`$
For the case $`F=C=G=K=J=0`$ the boundary parameters are nearly fixed, i.e.
$$\alpha _z=\pm \frac{\mathrm{}}{\sqrt{2}}\beta _z=\frac{\mathrm{}}{\sqrt{2}}$$
(20)
and $`\alpha _+=\beta _+=0`$ or $`\alpha _{}=\beta _{}=0`$. This case is special. The partition function we obtained does not fit into the picture suggested by the analogy to the free boson. Here an additional zero mode appears. The value of $`\mathrm{\Delta }`$ is given by $`\mathrm{\Delta }=0`$ for odd $`L`$ and by $`\mathrm{\Delta }=1/2`$ for even $`L`$. The partition function for $`H_{\mathrm{long}}`$ is $`𝒵_{\mathrm{long}}=8𝒵_{DD}(\mathrm{\Delta })`$ for $`L`$ even and $`L`$ odd.
### 2.5 Anomalous behaviour
For the last two cases in table 1 we obtained logarithmic corrections to the free surface energy. Note that there is no case in table 2 which corresponds to the boundaries in question. However, we computed the ground-state energy $`E_0`$ of $`H_{\mathrm{long}}`$ numerically for one example of boundaries which lead to the anomalous behaviour. Our computations suggest
$$E_0\frac{L}{\pi }+f_{\mathrm{}}+\frac{1}{L}\left(\rho _1\mathrm{ln}L+\rho _2\right)$$
(21)
where $`\rho _1,\rho _2`$ are complex numbers. Figure 1 and figure 2 show the real respectively the imaginary part of the Casimir amplitude. The computation has been done for $`\alpha _z=1/\sqrt{2},\beta _z=\mathrm{}/\sqrt{2},\alpha _\pm =\beta _\pm =0`$ which corresponds to the last case in table 1. In both cases the fermionic energies scale as
$$2\mathrm{\Lambda }\frac{1}{L}\left[k\pi \pm \frac{1}{2}\left(\mathrm{arg}\mathrm{\Delta }\mathrm{}\mathrm{ln}L\right)\right]k=0,1,2,\mathrm{}.$$
(22)
where the value of $`\mathrm{\Delta }`$ is given by
$$\mathrm{\Delta }=\frac{4C}{G}$$
(23)
for $`F=0,C,G0`$ respectively by
$$\mathrm{\Delta }=(1)^L\frac{2J}{14\alpha _z^2\beta _z^2}$$
(24)
for $`F=C=G=K=0,J0`$. Our result for the last case makes only sense for $`(\alpha _z\beta _z)^21/4`$. The factor 2 in (22) is just a remnant from our notation in . Note the integer spacing of the energy gaps (22).
## 3 The asymmetric XX-chain
We also considered the asymmetric XX-chain which is defined by the Hamiltonian
$`H_\mathrm{a}={\displaystyle \underset{j=1}{\overset{L1}{}}}\left[p\sigma _j^+\sigma _{j+1}^{}+q\sigma _j^{}\sigma _{j+1}^+\right]`$
$`+{\displaystyle \frac{1}{\sqrt{8}}}\left[\alpha _{}^{}\sigma _1^{}+\alpha _+^{}\sigma _1^++\alpha _z\sigma _1^z+\beta _{}^{}\sigma _L^{}+\beta _+^{}\sigma _L^++\beta _z\sigma _L^z\right].`$ (25)
Without loss of generality we restrict ourselves to $`\sqrt{pq}=\frac{1}{2}`$. Under this condition $`H_\mathrm{a}`$ can be mapped on the symmetric chain (1) with boundary parameters
$$\alpha _{}=Q^{\frac{1L}{2}}\alpha _{}^{}\alpha _+=Q^{\frac{L1}{2}}\alpha _+^{}\beta _{}=Q^{\frac{L1}{2}}\beta _{}^{}\beta _+=Q^{\frac{1L}{2}}\beta _+^{}$$
(26)
by a similarity transformation, where $`Q=\sqrt{\frac{q}{p}}`$. The diagonal terms remain unchanged. We will separate two cases in the following. First we will consider hermitian bulk terms, i.e.
$$p=\mathrm{}^{\mathrm{}\pi \kappa }/2,q=\mathrm{}^{\mathrm{}\pi \kappa }/2.$$
(27)
This corresponds to Dzyaloshinsky-Moriya type interactions in the bulk. Thereafter we will turn to non-hermitian bulk terms with $`p,q`$ in $``$.
### 3.1 Dzyaloshinsky-Moriya interactions
We restricted ourselves to hermitian boundaries. In this case the value of $`\kappa `$ has only an effect onto the spectra if non-diagonal boundary terms are present at both ends of the chain. We introduce the parametrization
$$\alpha _+^{}=R_\alpha ^{}\mathrm{}^{\mathrm{}\pi \phi ^{}}\alpha _{}^{}=R_\alpha ^{}\mathrm{}^{\mathrm{}\pi \phi ^{}}\beta _+^{}=R_\beta ^{}\mathrm{}^{\mathrm{}\pi (\chi ^{}+\phi ^{})}\beta _{}^{}=R_\beta ^{}\mathrm{}^{\mathrm{}\pi (\chi ^{}+\phi ^{})}.$$
(28)
According to (26) and (12) the mapping onto symmetric bulk terms yields
$$R_\alpha ^{}=R_\alpha R_\beta ^{}=R_\beta \phi =\phi ^{}+\kappa \frac{1L}{2}\chi =\chi ^{}+\kappa (L1)\text{mod}\mathrm{\hspace{0.17em}2}.$$
(29)
Remember that in the case of length independent boundary terms, the value of $`\chi `$ contains the full information about the partition function. Observe now, that it is possible to choose values of $`\kappa `$ such that the value of $`\chi `$ in (29) becomes independent of $`L`$ as long as one considers only certain sequences of lattice lengths. This is exactly the case if $`\kappa /2`$ is rational, i.e. $`\kappa /2=m/n,m,n^+`$. In this case $`\chi `$ is independent of $`L`$ for $`L=ln+r,0r<n,l`$, i.e.
$$\chi =\chi ^{}+\frac{2m(r1)}{n}.$$
(30)
The results obtained for the symmetric case can be adopted immediately. This is not possible for irrational values of $`\kappa /2`$. Note that this kind of commensurability and incommensurability has also been observed for the periodic chain with this type of interaction .
### 3.2 $`p,q`$
Without loss of generality, we are going to restrict ourselves to $`p>q`$ which implies $`Q<1`$. We have to distinguish two different situations. If $`\alpha _{}^{}\beta _+^{}`$ equals zero we can adopt the results for the symmetric case with $`F=C=0`$ (see appendix A for details). One has just to exchange $`\alpha _\pm `$ and $`\beta _\pm `$ by $`\alpha _\pm ^{}`$ and $`\beta _\pm ^{}`$. If otherwise $`\alpha _{}^{}\beta _+^{}0`$ the situation changes. The energy gaps suggest the partition function for the long chain (2) with boundary terms given by (26), i.e.
$`𝒵_{\mathrm{long}}=\mathrm{tr}z^{\frac{2L}{(Q+1/Q)\pi }(He_{\mathrm{}}Lf_{\mathrm{}})}={\displaystyle \frac{2}{\eta (z)}}{\displaystyle \underset{m/2}{}}z^{2(m+\frac{\mathrm{\Delta }_x}{2})^2+2mL\mathrm{\Delta }_y}`$ (31)
where
$$\mathrm{\Delta }_x=\frac{\mathrm{}}{2\pi }\mathrm{ln}\mathrm{\Gamma }\mathrm{\Delta }_y=\frac{\mathrm{}}{\pi }\frac{QQ^1}{Q+Q^1}$$
(32)
and
$`\mathrm{\Gamma }`$ $`=`$ $`(1Q^2(12\alpha _z^22\alpha _{}^{}\alpha _+^{})2\alpha _z^2Q^4)`$ (33)
$`\times (1Q^2(12\beta _z^22\beta _{}^{}\beta _+^{})2\beta _z^2Q^4)/(\alpha _{}^2\beta _+^2Q^4(1+Q^2)^2).`$
Note that the value of $`\mathrm{\Delta }_y`$ is purely imaginary. Hence the length dependent term given in the partition function is a phase. Such a term already appeared in the toroidal partition function for the asymmetric model with periodic boundary conditions . Note also that the result (31) simplifies to the partition function we obtained in the Neumann-Neumann case for $`H_{\mathrm{long}}`$ if one sets $`Q=1`$ (see (14)).
In we computed the exact ground-state energy on the finite chain for boundaries defined by
$$\alpha _z=\beta _z=0\alpha _+^{}\alpha _{}^{}=\beta _+^{}\beta _{}^{}=1.$$
(34)
If we introduce the parameter $`\chi ^{}`$ via $`\beta _+^{}=\mathrm{}^{\mathrm{}\pi \chi ^{}}\alpha _+^{}`$, then (33) simplifies to
$$\mathrm{\Gamma }=\left(\mathrm{}^{\mathrm{}\pi \chi ^{}}Q^2\right)^2.$$
(35)
Expanding the exact expression obtained in leads to
$`E_0{\displaystyle \frac{Q+Q^1}{2\pi }}L{\displaystyle \frac{Q+Q^1+(QQ^1)\frac{1}{2}\mathrm{ln}\mathrm{\Gamma }}{2\pi }}{\displaystyle \frac{(Q+Q^1)\pi }{2L}}\left({\displaystyle \frac{1}{24}}+{\displaystyle \frac{(\mathrm{ln}\mathrm{\Gamma })^2}{8\pi }}\right).`$ (36)
It is only for this case that we can perform the projection onto the $`(+,+)`$-sector. We obtain
$$𝒵=\frac{1}{\eta (z)}\underset{m}{}z^{2(m+\chi ^{}/2\mathrm{}\mathrm{ln}Q/\pi )^2+2mL\mathrm{\Delta }_y}\text{for odd }L$$
(37)
$$𝒵=\frac{1}{\eta (z)}\underset{m\frac{2+1}{2}}{}z^{2(m+\chi ^{}/2\mathrm{}\mathrm{ln}Q/\pi )^2+2mL\mathrm{\Delta }_y}\text{for even }L.$$
(38)
Note that our results for $`H_{\mathrm{long}}`$ only apply if (33) is different from zero. If the nominator of $`\mathrm{\Gamma }`$ vanishes we obtain logarithmic terms in the asymptotic behaviour of the energy gaps similar the ones obtained for the symmetric chain. We obtained the fermion energies
$`2\mathrm{\Lambda }{\displaystyle \frac{Q+Q^1}{2L}}\left\{k\pi \pm {\displaystyle \frac{1}{2}}\left[\mathrm{arg}\mathrm{\Delta }\mathrm{}\mathrm{ln}L\right]\right\}\pm \mathrm{}{\displaystyle \frac{QQ^1}{2}}`$ (39)
where
$$\mathrm{\Delta }=\frac{2\alpha _{}^2\beta _+^2Q^4(1+Q^2)^2}{2AQ^2+4Q^4(B+E^2)6Q^6(D+2E^2)+8E^2Q^8}.$$
(40)
The values of $`A,B,D,E`$ are given in (A). We are not going to consider the case, where the denominator in (40) vanishes.
## 4 Conclusions
In this paper we considered the spectra of the XX–model with boundary fields given by the Hamiltonian $`H`$ in (1). In order to obtain our results we also studied the Hamiltonian $`H_{\mathrm{long}}`$ defined by (2). Furthermore we considered the asymmetric XX–chain (3). Here we separated two cases. First we considered the case where the values of $`p`$ and $`q`$ in (3) are given by $`p=\mathrm{}^{\mathrm{}\pi \kappa }/2`$ and $`q=\mathrm{}^{\mathrm{}\pi \kappa }/2`$ (this corresponds to Dzyaloshinsky-Moriya interactions). Second we considered the case where $`p`$ and $`q`$ are real numbers.
For periodic boundary conditions the partition function for $`H`$ is given by the partition function of the free boson with periodic boundary conditions (5) . The partition function for $`H_{\mathrm{long}}`$ with periodic boundary conditions is just the expression (5) multiplied by 4. The spectra of the asymmetric chain with Dzyaloshinsky-Moriya interactions and periodic boundary conditions have been studied in , whereas the partition function for real values of $`p`$ and $`q`$ and periodic boundary conditions was given in .
The results we obtained for $`H`$ and $`H_{\mathrm{long}}`$ can be encoded in terms of the parameters given in (2). For reasons discussed in the text we obtained the energy gaps only for the cases given in table 1. Note that as long as we restrict ourselves to hermitian boundary terms we studied the most general case. This is not true for non-hermitian boundaries.
For $`H`$ with hermitian boundaries we obtained the partition functions corresponding to one of the three boson partition functions (7),(8) and (9). However, we found these partition functions also for certain non-hermitian boundary terms.
We obtained the Neumann-Neumann partition function (8) if all non-diagonal terms are present, where the value of $`\mathrm{\Delta }`$ in (8) is given by (13). The Dirichlet-Neumann partition function (9) is obtained for boundary terms which satisfy (16). We found the Dirichlet-Dirichlet partition function (7) for two types of boundary terms given by (18) respectively (17). The value of $`\mathrm{\Delta }`$ in (7) is given by (19) in both cases. Furthermore, for the case of non-hermitian boundaries we found a case which is special. For boundary terms given by (20) the partition function is the Dirichlet-Dirichlet partition function (7) multiplied by 2, where the value of $`\mathrm{\Delta }`$ is $`0`$ or $`1/2`$ for odd respectively even values of the lattice length.
The partition functions for $`H_{\mathrm{long}}`$ have also been considered. If all non-diagonal terms are present it is given by (14). For the other cases it is just the partition function of $`H`$ multiplied by 4. We also computed the values of the free surface energies of $`H_{\mathrm{long}}`$ and $`H`$ for certain boundary terms. They are given in table 2. In this table we have also given the values of the lowest highest weights which appear in the spectra of $`H_{\mathrm{long}}`$ for these boundary terms.
For the last two cases in table 1 we obtained logarithmic corrections to the free surface energy (equation (21)). Here we found also logarithmic terms in the asymptotic behaviour of the energy gaps (equation (22)).
In the case of Dzyaloshinsky-Moriya interactions, we restricted ourselves to hermitian chains. The value of the phase $`\kappa `$ (see (27)) has only an effect on the spectra if non-diagonal boundary terms are present at both ends of the chain. In this case we obtain the Neumann-Neumann partition function (8) if $`\kappa /2`$ is a rational number. The value of $`\chi `$ in the definition of $`\mathrm{\Delta }`$ in (13) has just to be exchanged by the expression in (30).
The partition function for the asymmetric chain (3) with real values of $`p`$ and $`q`$ has only been obtained for one special type of boundaries (see equation (34)). It is given by (37) for $`L`$ odd respectively by (38) for $`L`$ even. The asymptotic behaviour of the ground-state energy for this case is given in (36).
I would like to thank Birgit Wehefritz for her contribution to the early stage of this work. I would also like to thank Vladimir Rittenberg for many discussions and constant encouragement. I would like to thank Paul Pearce for helpful comments. I am grateful to Klaus Krebs for carefully reading the manuscript and many fruitful discussions. This work was supported by the TMR Network Contract FMRX-CT96-0012 of the European Commission.
## Appendix A Determination of the energy gaps
In we have seen that the spectrum of $`H_{\mathrm{long}}`$ is given in terms of free fermions. The fermionic energies $`2\mathrm{\Lambda }_n`$ are given by
$$2\mathrm{\Lambda }_n=\frac{1}{2}(x_n+x_n^1).$$
(41)
where the $`x_n`$ are the roots of the polynomial
$`p(x^2)=[x^{4L+8}+1A(x^{4L+6}+x^2)+(B+E^2)(x^{4L+4}+x^4)`$
$`+(D+2E^2)(x^{4L+2}+x^6)+E^2(x^{4L}+x^8)2E(x^{2L+8}+x^{2L})`$
$`+\left((ABD1)/2(1)^LC2E^2\right)(x^{2L+6}+x^{2L+2})`$
$`+(ABD1+2(1)^LC+4E4E^2)x^{2L+4}]/(x^21)^2`$ (42)
where the coefficients appearing in (A) are functions of the boundary parameters:
$`A=2(\alpha _{}\alpha _++\beta _{}\beta _++\alpha _z^2+\beta _z^21)C=\alpha _{}^2\beta _+^2+\alpha _+^2\beta _{}^2`$
$`B=(2\alpha _{}\alpha _+1)(2\beta _{}\beta _+1)+4\beta _z^2(\alpha _{}\alpha _+1)+4\alpha _z^2(\beta _+\beta _{}1)`$
$`D=\beta _z^2(4\alpha _{}\alpha _+2)+\alpha _z^2(4\beta _+\beta _{}2)E=2\alpha _z\beta _z.`$ (43)
Since there appear 4 zeros for each fermion, namely $`x_n,x_n^1,x_n`$ and $`x_n^1`$, the polynomial yields $`L+1`$ fermionic energies (see for details). In addition to these fermions there always exists a fermion with energy $`2\mathrm{\Lambda }_0=0`$, which we named ’spurious’ zero mode. All possible combinations of these $`L+2`$ fermions build up the spectrum of $`H_{\mathrm{long}}`$. The spectrum of $`H`$ is then obtained by excluding the ’spurious’ zero mode from the set of fermion energies and then taking the sector with an even or with an odd number of fermions being excited with respect to the vacuum. How to decide whether one has to pick up the even or the odd sector has been discussed in detail in .
Since in this paper we are interested into the large $`L`$ behaviour of the low lying energy levels, we look for the asymptotics of the zeros of $`p(x^2)`$ in the vicinity of the point $`x=\pm \mathrm{}`$ using the ansatz $`x=\mathrm{}^{\mathrm{}\frac{\pi }{2}\mathrm{}\frac{\varphi }{L}}`$, where we assume $`\varphi `$ to be a constant in leading order. For the first four cases given in table 1 we obtained the following equations:
* $`F0`$
$$\mathrm{cos}(2\varphi )+\frac{2C}{F}=0$$
(44)
* $`F=C=0,G0`$
$$\varphi \mathrm{sin}(2\varphi )=0$$
(45)
* $`F=C=G=0,K0`$
$$\varphi ^2\left[\mathrm{cos}(2\varphi )+(1)^{L+1}\frac{J}{K}\right]=0$$
(46)
* $`F=C=G=K=J=0`$
$$\varphi ^4\mathrm{exp}(2\mathrm{}\varphi )=(1)^L$$
(47)
Solving these equations for $`\varphi `$ yields the energy gaps (cf. (41)) in leading order. We obtained the following expressions:
* $`F0`$
$$2\mathrm{\Lambda }\frac{1}{L}\left[\frac{2n1}{2}\pi \pm \frac{1}{2}\text{arccos}\left(\frac{2C}{F}\right)\right]1n$$
(48)
* $`F=C=0,G0`$
$$2\mathrm{\Lambda }\frac{1}{L}\frac{n}{2}\pi 0n$$
(49)
* $`F=C=G=0,K0`$
$$2\mathrm{\Lambda }\frac{1}{L}\left[\frac{2n1}{2}\pi \pm \frac{1}{2}\text{arccos}\left((1)^{L+1}\frac{J}{K}\right)\right]1n$$
(50)
In addition to these modes there appears an additional zero mode.
* $`F=C=G=K=J=0`$
$$2\mathrm{\Lambda }\frac{1}{L}n\pi 1n\text{for }L\text{ even}$$
(51)
$$2\mathrm{\Lambda }\frac{1}{L}\frac{2n1}{2}\pi 1n\text{for }L\text{ odd}$$
(52)
Furthermore we found 3 additional zero modes for $`L`$ even and 2 additional zero modes for $`L`$ odd.
For the last two cases in table 1 the asymptotic zeros of the polynomial can only be found using the ansatz $`x=\mathrm{}^{\mathrm{}\frac{\pi }{2}\mathrm{}\frac{\varphi }{L}}`$ if we assume the imaginary part of $`\varphi `$ to diverge as $`L`$ goes to infinity. Hence we generalize our ansatz to $`x=\mathrm{}^{\mathrm{}\frac{\pi }{2}\mathrm{}\frac{\varphi (L)}{L}}`$, where $`\varphi (L)`$ is a complex function. Furthermore we assume $`lim_L\mathrm{}\frac{\varphi (L)}{L}=0`$ and $`lim_L\mathrm{}\mathrm{}^{\mathrm{}\varphi (L)}=0`$. The second assumption will be explained shortly. In both cases our ansatz leads to the solution of an equation of the form
$$a\mathrm{}^{2\mathrm{}\varphi (L)}\left[1+\mathrm{O}\left(\frac{\varphi (L)}{L}\right)\right]+\frac{2\mathrm{}\varphi (L)}{L}\left[b+\mathrm{O}\left(\frac{\varphi (L)}{L}\right)\right]=0.$$
(53)
From this expression one can see that our second assumption is indeed necessary to solve this equation, since otherwise the first term would be finite for all values of $`L`$, whereas the second term vanishes as $`L`$ goes to infinity. Neglecting the terms of order $`\varphi (L)/L`$ this equation is solved in terms of the so called Lambert W function $``$ which is defined by the property $`(x)\mathrm{}^{(x)}=x`$. We obtain
$$\varphi (L)=\mathrm{}(\mathrm{\Delta }L)/2\text{where}\mathrm{\Delta }=\frac{a}{b}.$$
(54)
The asymptotic behaviour of $``$ is well known , i.e.
$$(L)2\mathrm{}k\pi +\mathrm{ln}L\mathrm{ln}(\mathrm{ln}L+2\mathrm{}\pi k)+\mathrm{}$$
(55)
Note that this expression is in accordance with our assumptions we made concerning the asymptotic behaviour of $`\varphi (L)`$.
In this paper we also considered the asymmetric chain, which is similar to the symmetric chain with length dependent boundary parameters (cf. (26)). This length dependence enters the polynomial (A) only via the coefficient $`C`$, which becomes
$$C=\alpha _{}^2\beta _+^2Q^{22L}+\alpha _+^2\beta _{}^2Q^{2L2}.$$
(56)
For $`Q<1`$ (the case considered in section 3.2) the first term on the RHS of (56) diverges exponentially whereas the second term vanishes exponentially as a function of the lattice length $`L`$. Hence we have to distinguish two different situations. If $`\alpha _{}^{}\beta _+^{}`$ equals zero we may find the asymptotic zeros of the polynomial as for the symmetric case with $`F=C=0`$. If otherwise $`\alpha _{}^{}\beta _+^{}0`$ then $`C`$ diverges as $`L`$ is increased. This has to be compensated by modifying our ansatz to $`x=\mathrm{}^{\mathrm{}\frac{\pi }{2}\mathrm{}\frac{\varphi }{L}+\mathrm{ln}Q}`$ . This works as long as the nominator of $`\mathrm{\Gamma }`$ in (33) is different from zero. Otherwise we have to modify our ansatz another time to $`x=\mathrm{}^{\mathrm{}\frac{\pi }{2}\mathrm{}\frac{\varphi (L)}{L}+\mathrm{ln}Q}`$, where $`lim_L\mathrm{}\frac{\varphi (L)}{L}=0`$ and $`lim_L\mathrm{}\mathrm{}^{\mathrm{}\varphi (L)}=0`$. This leads again to an equation of the form (53) which can be solved as described above.
## References |
warning/0002/hep-lat0002012.html | ar5iv | text | # Quasi-asymptotic freedom in the two dimensional 𝑂(3) model
\[
## Abstract
The behavior of the renormalized spin 2-point function in the $`O(3)`$ and dodecahedron spin model are investigated numerically. The Monte Carlo data show excellent agreement between the two models. The short distance behavior comes very close to standard theoretical expectations, yet it differs significantly from it. A possible explanation of this situation is offered.
11.25.Bt, 11.15.Ha, 75.10.Hk
\]
Are two dimensional (2D) nonabelian $`\sigma `$-models asymptotically free? After more than two decades, the answer remains unclear, with lots of supportive evidence on both sides of the issue. The main point of this letter is to provide some fresh evidence that even though the orthodox picture provides an excellent description of physics at intermediate momemta and distances not too short, it is in fact incorrect, a situation which I belive occurs also in $`QCD_4`$.
To that end I will present the result of Monte Carlo (MC) investigations of the continuum limit of the spin 2-point function in four spin models: i) Ising model, ii) $`O(2)`$ model, iii) $`O(3)`$ model and iv) dodecahedron spin model. As I will show, there is strong numerical evidence that the continuum limit of the dodecahedron spin model is identical to that of the $`O(3)`$ model. Undergoing a freezing transition at nonzero temperature, the latter is not likely to be asymptotically free. Moreover the continuum limit of the spin 2-point function of the $`O(3)`$ model disagrees with the PT formula with the value of $`\mathrm{\Lambda }/m`$ predicted by Hasenfratz, Maggiore and Niedermayer based on the thermodynamic Bethe ansatz and the normalization constant fixed by the scaling hypothesis of the form factor approach proposed by Balog and Niedermaier (BN) .
The numerical data suggest that in all the four models investigated, the short distance behaviour of the spin 2-point function is of the form $`1/x^{1/4}`$.
As a biproduct of my investigation, I report also on the behavior of the critical exponent $`\eta `$, which probably equals 1/4 in all four models. The subject of logarithmic corrections to $`\eta =1/4`$ in the $`O(2)`$ model has been discussed in several recent papers ,, ,, since a theoretical prediction stemming from the perturbative renormalization group approach exists and it is not clear that the data are consistent with it. As I will discuss shortly, the existence of a continuum limit obtainable from the lattice model via multiplicative renormalization, an assumption which seems to be corroborated by the numerics, suggests that no such logarithmic corrections are present.
I would like to begin by reviewing the intimate connection between the short distance behavior of the continuum spin 2-point function and the critical exponent $`\eta `$. First some definitions: let $`P=(p,0)`$ and
\- Spin 2-point function $`G(p)`$:
$$G(p)=\frac{1}{L^2}|\widehat{s}(P)|^2;\widehat{s}(P)=\underset{x}{}e^{iPx}s(x)$$
(1)
\- Magnetic susceptibility $`\chi `$:
$$\chi =G(0)$$
(2)
\- Correlation length $`\xi `$:
$$\xi =\frac{1}{2\mathrm{sin}(\pi /L)}\sqrt{(G(0)/G(1)1)}$$
(3)
As the continuum spin 2-point function I will use the limit $`\xi \mathrm{}`$ of
$$G(x/\xi )s(0).s(x)\xi ^2/\chi $$
(4)
Let $`z`$ be the physical continuum distance ($`zx/\xi `$, limit $`\xi \mathrm{}`$); the exponents $`\overline{\eta }`$ and $`\overline{r}`$ are defined by requiring that in the limit $`z0`$ the expression
$$G(z)z^{\overline{\eta }}(log(z))^{2\overline{r}}$$
(5)
is finite and different from 0. Let $`\eta `$ and $`r`$ be defined so that in the limit $`\xi \mathrm{}`$ the quantity
$$c_{\mathrm{𝑘𝑡}}\frac{\chi }{\xi ^{2\eta }log(\xi )^{2r}}$$
(6)
is finite and different from 0.
I would like to claim that the existence of a continuum limit with the short distance behaviour described by (5) requires $`\eta =\overline{\eta }`$ and $`r=\overline{r}`$. Indeed in the continuum limit, the limit $`z0`$ is reached by letting both $`x\mathrm{}`$ and $`\xi \mathrm{}`$ but $`x/\xi 0`$. In fact the same consideration shows that in the same limit the lattice spin 2-point function must have the property
$$s(0).s(x)x^{\overline{\eta }}$$
(7)
goes to a finite nonzero constant (i.e. there are no logarithmic modifications).
The Kosterlitz prediction is that $`\overline{r}=1/16`$. Finite size scaling at the kt-point do indicate a negative $`r`$ though smaller in magnitude . The paradox discussed by many , , is that the thermodynamic data suggest a positive value for $`r`$ ,. As explained above this is consistent with the fact that $`r=\overline{r}`$, rather than $`r=\overline{r}`$ as incorrectly claimed in the literature.
However the fact that eq.(7) requires the lattice spin 2-point function to decay as a pure power with no logarithmic corrections suggests that most likely $`r=\overline{r}=0`$. Indeed while for $`\xi <\mathrm{}`$ eq.(7) is valid only for $`x\xi `$, baring some nonuniformity, eq.(7) ought to be describing the large distance behaviour of the lattice spin 2-point function at the critical point.
Next I would like to discuss the main point of this paper, which concerns the asymptotic behaviour of the continuum spin 2-point function $`G(z)`$ at short distances. It is known that in the Ising model $`\eta =1/4`$ $`r=0`$. I stated above what perturbative renormalization group (RG) arguments predict for $`O(2)`$. For $`O(3)`$ similar arguments predict $`\eta =0`$ $`r=1`$. However when combined with the thermodynamic Bethe ansatz and the scaling hypothesis in the form factor approach, this latter prediction becomes even more precise :
$`G(z)`$ $`=`$ $`{\displaystyle \frac{1.000034657}{3\pi ^3\mathrm{\hspace{0.17em}1.001687}^2}}(t+ln(t)`$ (8)
$`+`$ $`1.1159+ln(t)/t+.1159/t)^2,`$ (9)
where $`tln(z\mathrm{\Lambda }),\mathrm{\Lambda }=e/8.`$ I am not aware of any predictions regarding $`\eta `$ and $`r`$ for the dodecahedron spin model.
In Fig.1 I present my Monte Carlo results for $`G(z)z^{1/4}`$ for (in rising order) Ising, $`O(3)`$ and $`O(2)`$; I also show the curve for free field (lowest curve) and the PT curve with the normalization constant predicted by Balog and Niedermaier (upper curve) . The MC data were taken at $`\xi 167`$ and L=1230. They represent the on axis correlation and I will discuss shortly the lattice artefacts. For now I would like to make the following observations:
1. In all three models $`G(z)`$ approaches the free field value for $`z\mathrm{}`$, as required by the Orstein-Zernke behaviour .
2. For $`z.05`$ $`G(z)`$ in the three models behaves quite similarly.
3. The data suggest that in all three models $`\eta =1/4`$, but that compared to the Ising model the short distance behavior may be more singular in $`O(2)`$ and less singular in $`O(3)`$. This is would be consistent with a negative $`\overline{r}`$ in $`O(2)`$ but a positive one in $`O(3)`$ (although the data are also consistent with $`\overline{r}=0`$).
4. Even though for $`.02<z<.1`$ the difference with the BN refined PT curve is small (about 2$`\%`$), it is clearly there and, as I will indicate next, it is not a lattice artefact.
In Fig.2 I present the MC results for the $`O(3)`$ model for $`\xi `$ ranging from 65 to 310. The error bars are about .3$`\%`$, too small to show. The data were produced using Wolff’s method . At each $`\beta `$ and $`L`$ value, at least five independent runs were made and the error estimated using the jack-knife method. The primary source of error are the values of $`\chi `$ and $`\xi `$, whose determination involves large distances. To give a feeling for the size of the lattice artefacts, for $`\xi 310`$ the (solid) curve represents $`G(z)`$ on axis while the points $`G(z)`$ along the diagonal. Both this test and the good agreement of these data with the data at $`\xi 230`$ suggest that, barring some very slow drift, the continuum limit has been reached for $`z.03`$. The continuum values are clearly off the BN refined PT curve (upper curve).
As explained above, the critical behaviour of $`\chi `$ and $`\xi `$ carry relevant information regarding the short distance behaviour of $`G(z)`$. In Fig.3 I represent $`ln(\chi )1.75ln(\xi )`$. The upper curve represents the MC data. The lower (dashed) curve represents the ratio of the values of this quantity computed via the BN formula
$$\frac{\chi }{\xi ^2}=\frac{3\pi \mathrm{\hspace{0.17em}1.001687}^2}{\mathrm{4\hspace{0.17em}1.000034657}\beta ^2}\left\{1+\frac{.1816}{\beta }+\frac{.133}{\beta ^2}+\frac{.1362}{\beta ^3}\right\}$$
(10)
over its MC values. A correct prediction of the BN formula would require the dashed curve to approach 1 for $`\xi \mathrm{}`$. The data produce a (dashed) curve which is not flat and which seems to approache the line at 1 with a nonzero slope.
The upper (solid) curve indicates that $`r`$ could be different from 0 and in fact negative. To determine whether there is a logarithmic modification to $`\eta =1/4`$ and obtain a precise value for $`r`$ one would need data at larger $`\xi `$. However the lack of smoothness of the curve suggests that the errors are underestimated already at $`\xi 310`$ (probably due to critical slowing down).
The last piece of numerical evidence which suggests that $`O(3)`$ is not asymptotically free comes from a comparison of its spin 2-point function with that of of the dodecahedron spin model. This is shown in Fig. 4, which shows $`G(z)`$ for the latter model at $`\xi 11,19,34,66`$ and 120 (broken line) and $`G(z)`$ for the $`O(3)`$ model at $`\xi 122`$ (solid line on axis, dots on the diagonal). As the correlation length increases, $`G(z)`$ for the dodecahedron model approaches that of $`O(3)`$ and at $`\xi 120`$ the two curves are practically indistinguishable for $`z.04`$. It appears rather certain that the two models share the same $`G(z)`$ in the continuum limit. However, as already mentioned, it is highly unlikely that the dodecahedron model, undergoing freezing at finite $`\beta `$, exhibits asymptotic freedom.
$`\mathrm{𝐷𝑖𝑠𝑐𝑢𝑠𝑠𝑖𝑜𝑛}`$
1. I believe the numerical evidence presented above suggests that in the $`O(3)`$ nonlinear $`\sigma `$ model the short distance behaviour of $`G(z)`$ corresponds to $`\overline{\eta }=1/4`$, $`\overline{r}0`$.
2. Perturbative renormalization group arguments predict that in all $`O(N)`$ models $`N3`$ $`\eta =0`$, $`\overline{r}`$ is negative and varies monotonically with $`N`$ from -1 for $`O(3)`$ to -1/2 for $`O(\mathrm{})`$. Based on the observed behavior of $`O(3)`$ I would conjecture that for all $`3N<\mathrm{}`$ $`\eta =1/4`$. Since the spherical model $`O(\mathrm{})`$ has $`\eta =0`$, $`\overline{r}=1/2`$, such a scenario would correspond to a nonuniformity of the $`1/N`$ expansion. The uniformity of the $`1/N`$ expansion for $`\beta \mathrm{}`$ has been questioned before .
3. Since for $`.02z.1`$ the deviation from the BN refined PT prediction is 2$`\%`$ or less, one may wonder why that is so? I think the answer is in Fig.1, where I compare $`G(z)`$ for different models. Please note that the free field curve describes also the spherical model and $`G(z)`$ for it is quite close to that of $`O(3)`$ for $`z`$ not too small. It is well known that the $`1/N`$ expansion is legitimate at fixed $`\stackrel{~}{\beta }\frac{\beta }{N}`$ and the only issue is whether this expansion is uniform for $`\stackrel{~}{\beta }\mathrm{}`$. The reason the orthodoxy comes so close to the truth is two fold: i) the fast convergence of the $`1/N`$ expansion for moderate $`\stackrel{~}{\beta }`$ and ii) the large value of $`\beta _{crt}`$. Thus, even though I believe both the $`1/N`$ expansion and perturbation theory are nonuniform, they come very close to the truth at moderate correlation length. Since for distances not too small or momenta not too large the continuum limit is well described at moderate $`\xi `$ values, the orthodoxy will describe correctly physics at such distances and/or momenta. It is only at short distances or large momenta that the true test of the orthodoxy can be performed.
4. Until now I discussed only some new numerical information. Some readers, such as those who recently wrote ‘We believe however that the continuing accumulation of unambiguous, consistent and increasingly accurate numerical support for the RG predictions from a variety of independent approaches leaves liitle if any space for alternative pictures.’ , may not find my numerical evidence conclusive. I will invite them then to reconsider the two theoretical arguments advanced by Patrascioiu and Seiler regarding the existence of a massless phase in all $`O(N)`$ models:
i) We showed rigorously that if a certain cluster, baptized the equatorial cluster, does not percolate, the model must be massless. We also advanced some nonrigorous arguments why that cluster did not seem likely to percolate . After seven years since those arguments were advanced, no mathematical physicist has proved the contrary or given us an example of how this equatorial cluster might percolate (see Open Problems in Mathematical Physics, http://www.iamp.org). Can any skeptical reader give such an example?
ii) We also showed rigorously that among smooth configurations, the dominant ones must be the superinstantons not the much publicized instantons. But in a superinstanton configuration there is a ring of arbitrary size in which the spin points roughly in a certain direction. Hence in such a configuration, the inverse image of a small patch of the sphere will form a ring and hence the equatorial cluster will not percolate. Therefore the existence of superinstantons as the dominant configuration at weak coupling reinforces the argument for the existence of a massless phase and hence the absence of asymptotic freedom in the massive phase.
I think irrespective of any numerics, to paraphrase Butera and Comi , ‘the two arguments mentioned above leave little if any space for the standard picture to be correct’.
5. Even though strictly incorrect, the standard picture provides an excellent phenomenological description at distances not too small ($`z>.02`$) or momenta not too large ($`p<2\pi /.02300`$. This is not something never seen before. Indeed it is taught in every course in Quantum Mechanics that the probability to find a particle trapped in a potential well (say an $`\alpha `$-particle) decays exponentially in time. The time constant of that exponential is called the life time of the state, and innumerable experiments over the years have measured such life times. I am not aware of anybody ever having claimed to have detected deviations from this alleged exponential law. Yet it is a mathematical fact that if the particle decays at all, the large time asymptotic behavior cannot be exponential, but some inverse power of the time . This fact is little known in the physics community (witness the many papers by particle physicists computing the lifetime of the ‘false cosmological vacuum’) because in fact the exponential law describes very well the data. My claim is that a similar situation occurs in the 2D nonlinear $`\sigma `$ models (and probably $`QCD_4`$): even though eventually, at $`\beta `$ sufficiently large superinstantons win and render the model massless, at moderate values of $`\beta `$ such as $`\beta =2`$ for $`O(3)`$, instantons (localized defects) dominate the typical configuration, which behaves very much as predicted by perturbation theory and/or the $`1/N`$ expansion. Only by probing sufficiently small distances or large momenta should one be able to detect major deviations from the expected behavior.
These ideas stem from my long term collaboration with Erhard Seiler and were crystalized by my recent collaboration with J.Balog, M.Niedermaier, F.Niedermayer and P.Weisz. |
warning/0002/math0002017.html | ar5iv | text | # Verlinde bundles and generalized theta linear series
## 1. Preliminaries and notations
*1.1. Generalized theta divisors.* Let $`U_X(r,0)`$ be the moduli space of equivalence classes of semistable vector bundles of rank $`r`$ and degree $`0`$ on $`X`$, $`SU_X(r,L)`$ the moduli space of vector bundles with fixed determinant $`L`$ Pic$`{}_{}{}^{0}(X)`$ and $`U_X^s(r,0)`$ and $`SU_X^s(r,L)`$ the open subsets corresponding to stable bundles. We recall the construction and some basic facts about generalized theta divisors on these moduli spaces, drawing especially on . Analogous constructions work for any degree $`d`$, as we will recall in the last section.
Let $`N`$ be a line bundle of degree $`g1`$ on $`X`$. Then $`\chi (EN)=0`$ for all $`EU_X(r,0)`$. Consider the subset of $`U_X^s(r,0)`$
$$\mathrm{\Theta }_N^s:=\{EU_X^s(r,0)|h^0(EN)0\}$$
and the analogous set in $`SU_X^s(r,L)`$. One can prove (see (7.4.2)) that $`\mathrm{\Theta }_N^s`$ is a hypersurface in $`U_X^s(r,0)`$ (resp. $`SU_X^s(r,L)`$). Denote by $`\mathrm{\Theta }_N`$ the closure of $`\mathrm{\Theta }_N^s`$ in $`U_X(r,0)`$ and $`SU_X(r,L)`$. As we vary $`N`$, these hypersurfaces are called $`\mathrm{𝑔𝑒𝑛𝑒𝑟𝑎𝑙𝑖𝑧𝑒𝑑}\mathrm{𝑡ℎ𝑒𝑡𝑎}\mathrm{𝑑𝑖𝑣𝑖𝑠𝑜𝑟𝑠}`$. It is proved in , Theorem A, that $`U_X(r,0)`$ and $`SU_X(r,L)`$ are locally factorial and so the generalized theta divisors determine line bundles $`𝒪(\mathrm{\Theta }_N)`$ on these moduli spaces. We have the following important facts:
###### Theorem 1.1.
(, Theorem B) The line bundle $`𝒪(\mathrm{\Theta }_N)`$ on $`SU_X(r,L)`$ does not depend on the choice of $`N`$. The Picard group of $`SU_X(r,L)`$ is isomorphic to $`𝐙`$, generated by $`𝒪(\mathrm{\Theta }_N)`$.
The line bundle in the theorem above, independent of the choice of $`N`$, is denoted by $``$ and is called the $`\mathrm{𝑑𝑒𝑡𝑒𝑟𝑚𝑖𝑛𝑎𝑛𝑡}`$ $`\mathrm{𝑏𝑢𝑛𝑑𝑙𝑒}`$.
###### Theorem 1.2.
(, Theorem C) The inclusions $`\mathrm{Pic}(J(X))\mathrm{Pic}(U_X(r,0))`$ (given by the determinant morphism) and $`𝐙𝒪(\mathrm{\Theta }_N)\mathrm{Pic}(U_X(r,0))`$ induce an isomorphism
$$\mathrm{Pic}(U_X(r,0))\mathrm{Pic}(J(X))𝐙.$$
More generally, for any vector bundle $`F`$ of rank $`k`$ and degree $`k(g1)`$, we can define
$$\mathrm{\Theta }_F^s:=\{EU_X^s(r,0)|h^0(EF)0\}$$
and denote by $`\mathrm{\Theta }_F`$ the closure of $`\mathrm{\Theta }_F^s`$ in $`U_X(r,0)`$. It is clear that, for generic $`F`$ at least, $`\mathrm{\Theta }_F`$ is strictly contained in $`U_X(r,0)`$ (in which case it is again a divisor). It is useful to know what is the dependence of $`𝒪(\mathrm{\Theta }_F)`$ on $`F`$ and in this direction we have:
###### Proposition 1.3.
( (7.4.3) and Prop.3) Let $`F`$ and $`G`$ be two vector bundles of slope $`g1`$ on $`X`$. If $`\mathrm{rk}(F)=m\mathrm{rk}(G)`$, then
$$𝒪(\mathrm{\Theta }_F)𝒪(\mathrm{\Theta }_G)^m\mathrm{det}^{}(\mathrm{det}F(\mathrm{detG})^m)$$
where we use the natural identification of $`\mathrm{Pic}^0(X)`$ with $`\mathrm{Pic}^0(J(X))`$. In particular, if $`F`$ has rank $`k`$ and $`N`$ is a line bundle of degree $`g1`$, we get
$$𝒪(\mathrm{\Theta }_F)𝒪(\mathrm{\Theta }_N)^k\mathrm{det}^{}(\mathrm{det}FN^k).$$
*1.2. A convention on theta divisors.* When looking at generalized theta divisors $`\mathrm{\Theta }_N`$ and the corresponding linear series, it will be convenient to consider the line bundle $`N`$ to be a $`\mathrm{𝑡ℎ𝑒𝑡𝑎}\mathrm{𝑐ℎ𝑎𝑟𝑎𝑐𝑡𝑒𝑟𝑖𝑠𝑡𝑖𝑐}`$, i.e. satisfying $`N^2\omega _X`$. The assumption brings some simplifications to most of the arguments, but on the other hand this case implies all the results for arbitrary $`N`$. This is true since for any $`N`$ and $`M`$ in Pic$`{}_{}{}^{g1}(X)`$, if $`\xi :=NM^1`$, twisting by $`\xi `$ gives an automorphism:
$$U_X(r,0)\stackrel{\xi }{}U_X(r,0)$$
by which $`𝒪(\mathrm{\Theta }_M)`$ corresponds to $`𝒪(\mathrm{\Theta }_N)`$. As an example, we will use freely isomorphisms of the type $`r_J^{}𝒪_J(\mathrm{\Theta }_N)𝒪_J(r^2\mathrm{\Theta }_N)`$, where $`r_J`$ is multiplication by $`r`$ on $`J(X)`$, which in general would work only up to numerical equivalence. One can easily see that in each particular proof the arguments could be worked out in the general situation with little extra effort (the main point is that the cohomological arguments work even if we use numerical equivalence instead of linear equivalence).
It is worth mentioning another convention about the notation that we will be using. The divisors $`\mathrm{\Theta }_N`$ make sense of course on both $`U_X(r,0)`$ for $`r2`$ and $`J(X)=U_X(1,0)`$ and in some proofs both versions will be used. We will denote the associated line bundle simply by $`𝒪(\mathrm{\Theta }_N)`$ if $`\mathrm{\Theta }_N`$ lives on $`U_X(r,0)`$ and by $`𝒪_J(\mathrm{\Theta }_N)`$ if it lives on the Jacobian.
*1.3. The Fourier-Mukai transform on an abelian variety.* Here we give a brief overview of some basic facts on the Fourier-Mukai transform on an abelian variety, following the original paper of Mukai . Let $`X`$ be an abelian variety of dimension $`g`$, $`\widehat{X}`$ its dual and $`𝒫`$ the Poincaré line bundle on $`X\times \widehat{X}`$, normalized such that $`𝒫_{|X\times \{0\}}`$ and $`𝒫_{|\{0\}\times \widehat{X}}`$ are trivial. To any coherent sheaf $``$ on $`X`$ we can associate the sheaf $`p_{2}^{}{}_{}{}^{}(p_{1}^{}{}_{}{}^{}𝒫)`$ on $`\widehat{X}`$ via the natural diagram:
This correspondence gives a functor
$$𝒮:\mathrm{Coh}(\mathrm{X})\mathrm{Coh}(\widehat{X}).$$
If we denote by $`D(X)`$ and $`D(\widehat{X})`$ the derived categories of Coh$`(X)`$ and Coh$`(\widehat{X})`$, then the derived functor $`𝐑𝒮:D(X)D(\widehat{X})`$ is defined (and called the Fourier functor) and one can consider $`𝐑\widehat{𝒮}:D(\widehat{X})D(X)`$ in a similar way. Mukai’s main theorem is the following:
###### Theorem 1.4.
( (2.2)) The Fourier functor establishes an equivalence of categories between $`D(X)`$ and $`D(\widehat{X})`$. More precisely there are isomorphisms of functors:
$$𝐑𝒮𝐑\widehat{𝒮}(1_{\widehat{X}})^{}[g]$$
$$𝐑\widehat{𝒮}𝐑𝒮(1_X)^{}[g].$$
In this paper we will essentially have to deal with the simple situation when by applying the Fourier functor we get back another vector bundle, i.e. a complex with only one nonzero (and locally free) term. This is packaged in the following definition (see (2.3):
###### Definition 1.5.
A coherent sheaf $``$ on $`X`$ satisfies I.T. (index theorem) with $`\mathrm{𝑖𝑛𝑑𝑒𝑥}j`$ if $`H^i(\alpha )=0`$ for all $`\alpha \mathrm{Pic}^0(X)`$ and all $`ij`$. In this situation we have $`R^i𝒮()=0`$ for all $`ij`$ and by the base change theorem $`R^j𝒮()`$ is locally free. We denote $`R^j𝒮()`$ by $`\widehat{}`$ and call it the $`\mathrm{𝐹𝑜𝑢𝑟𝑖𝑒𝑟}\mathrm{𝑡𝑟𝑎𝑛𝑠𝑓𝑜𝑟𝑚}`$ of $``$. Note that then $`𝐑𝒮()\widehat{F}[j]`$.
For later reference we list some basic properties of the Fourier transform that will be used repeatedly throughout the paper:
###### Proposition 1.6.
Let $`X`$ be an abelian variety and $`\widehat{X}`$ its dual and let $`𝐑𝒮`$ and $`𝐑\widehat{𝒮}`$ the corresponding Fourier functors. Then the following are true:
(1)( (3.4)) Let $`Y`$ be an abelian veriety, $`f:YX`$ an isogeny and $`\widehat{f}:\widehat{X}\widehat{Y}`$ the dual isogeny. Then there are isomorphisms of functors:
$$f^{}𝐑\widehat{𝒮}_X𝐑\widehat{𝒮}_Y\widehat{f}_{}$$
$$f_{}𝐑\widehat{𝒮}_Y𝐑\widehat{𝒮}_X\widehat{f}^{}.$$
(2)( (3.7)) Let $``$ and $`𝒢`$ coherent sheaves on $`X`$ and define their Pontrjagin product by $`𝒢:=\mu _{}(p_{1}^{}{}_{}{}^{}p_{2}^{}{}_{}{}^{}𝒢)`$, where $`\mu :X\times XX`$ is the multiplication on $`X`$. Then we have the following isomorphisms:
$$𝐑𝒮(𝒢)𝐑𝒮()𝐑𝒮(𝒢)$$
$$𝐑𝒮(𝒢)𝐑𝒮()𝐑𝒮(𝒢)[g],$$
where the operations on the right hand side should be thought of in the derived category.
(3)( (3.11)) Let $`L`$ be a nondegenerate line bundle on $`X`$ of index $`i`$, i.e. $`h^i(L)0`$ and $`h^j(L)=0`$ for all $`ji`$. Then by §16, I.T. holds for $`L`$ and there is an isomorphism
$$\varphi _L^{}\widehat{L}H^i(L)L^1(\underset{|\chi (L)|}{}L^1),$$
with $`\varphi _L`$ the isogeny canonically defined by L.
(4)( (3.1)) Let $``$ be a coherent sheaf on $`X`$, $`x\widehat{X}`$ and $`P_x\mathrm{Pic}^0(X)`$ the corresponding line bundle. Then we have an isomorphism:
$$𝐑𝒮(P_x)t_x^{}𝐑𝒮(),$$
where $`t_x`$ is translation by $`x`$.
(5)( (2.4),(2.5),(2.8)) Assume that $`E`$ satisfies I.T. with index $`i`$. Then $`\widehat{E}`$ satisfies I.T. with index $`gi`$ and in this case
$$\chi (E)=(1)^i\mathrm{rk}(\widehat{E}).$$
Moreover, there are isomorphisms $`\mathrm{Ext}^k(E,E)\mathrm{Ext}^k(\widehat{E},\widehat{E})`$ for all $`k`$, and in particular $`E`$ is simple if and only if $`\widehat{E}`$ is simple.
*1.4. Global generation and normal generation of vector bundles on abelian varieties.* For the reader’s convenience, in the present paragraph we give a brief account of some very recent results and techniques in the study of vector bundles on abelian varieties – following work of Pareschi – in a form convenient for our purposes. The underlying theme is to give useful criteria for the global generation and surjectivity of multiplication maps of such vector bundles.
Let $`X`$ be an abelian variety and $`E`$ a vector bundle on $`X`$. Building on earlier work of Kempf , Pareschi proves the following cohomological criterion for global generation:
###### Theorem 1.7.
( (2.1)) Assume that $`E`$ satisfies the following vanishing property:
$$h^i(E\alpha )=0,\alpha \mathrm{Pic}^0(X)\mathrm{and}i>0.$$
Then for any ample line bundle $`L`$ on $`X`$, $`EL`$ is globally generated.
In another direction, in order to attack questions about multiplication maps of the form
(1)
$$H^0(E)H^0(F)H^0(EF)$$
for $`E`$ and $`F`$ vector bundles on $`X`$, the right notion turns out to be that of skew Pontrjagin product:
###### Definition 1.8.
Let $`E`$ and $`F`$ be coherent sheaves on the abelian variety $`X`$. Then the $`\mathrm{𝑠𝑘𝑒𝑤}\mathrm{𝑃𝑜𝑛𝑡𝑟𝑗𝑎𝑔𝑖𝑛}\mathrm{𝑝𝑟𝑜𝑑𝑢𝑐𝑡}`$ of $`E`$ and $`F`$ is defined by
$$E\widehat{}F:=p_{1}^{}{}_{}{}^{}((p_1+p_2)^{}Ep_2^{}F)$$
where $`p_1,p_2:X\times XX`$ are the two projections.
The following is a simple but essential result relating the skew Pontrjagin product to the surjectivity of the multiplication map (1). It is a restatement of (1.1) in a form convenient to us and we reproduce Pareschi’s argument for the sake of completeness.
###### Proposition 1.9.
Assume that $`E\widehat{}F`$ is globally generated and that $`h^i(t_x^{}EF)=0`$ for all $`xX`$ and all $`i>0`$, where $`t_x`$ is the translation by $`x`$. Then for all $`xX`$ the multiplication map
$$H^0(t_x^{}E)H^0(F)H^0(t_x^{}EF)$$
is surjective and in particular (1) is surjective.
###### Proof.
The fact that $`h^i(t_x^{}EF)=0`$ for all $`i>0`$ and all $`xX`$ implies by base change that $`E\widehat{}F`$ is locally free with fiber
$$E\widehat{}F(x)H^0(t_x^{}EF).$$
We also have a natural isomorphism
$$\phi :H^0(E\widehat{}F)\stackrel{}{}H^0(E)H^0(F)$$
obtained as follows. One one hand by Leray we naturally have
$$H^0(p_{1}^{}{}_{}{}^{}((p_1+p_2)^{}Ep_2^{}F))H^0((p_1+p_2)^{}Ep_2^{}F)$$
and on the other hand the automorphism $`(p_1+p_2,p_2)`$ of $`X\times X`$ induces an isomorphism
$$H^0((p_1+p_2)^{}Ep_2^{}F)H^0(E)H^0(F),$$
so $`\phi `$ is obtained by composition. If we identify $`H^0(E)H^0(F)`$ with both $`H^0(E\widehat{}F)`$ (via $`\phi `$) and $`H^0(t_x^{}E)H^0(F)`$ (via $`t_x^{}\times id`$), then it is easily seen that the multiplication map
$$H^0(t_x^{}E)H^0(F)H^0(t_x^{}EF)$$
coincides with the evaluation map
$$H^0(E\widehat{}F)\stackrel{ev_x}{}E\widehat{}F(x)$$
and this proves the assertion. ∎
###### Remark 1.10.
There is a clear relationship between the skew Pontrjagin product and the usual Pontrjagin product as defined in 1.6(2). In fact (see (1.2)):
$$E\widehat{}FE(1_X)^{}F.$$
If $`F`$ is symmetric the two notions coincide and one can hope to apply results like 1.6(2). This whole circle of ideas will be used in Section 5 below.
Finally we extract another result on multiplication maps that will be useful in the sequel. It is a particular case of (3.8), implicit in Kempf’s work on syzygies of abelian varieties.
###### Proposition 1.11.
Let $`E`$ be a vector bundle on $`X`$, $`L`$ an ample line bundle and $`m2`$ an integer. Assume that
$$h^i(EL^k\alpha )=0,i>0,k2\mathrm{and}\alpha \mathrm{Pic}^0(X).$$
Then the multiplication map
$$H^0(L^m)H^0(E)H^0(L^mE)$$
is surjective.
## 2. The Verlinde bundles $`E_{r,k}`$
In the present section we define the main objects of this paper and study their first properties. These are vector bundles on the Jacobian of $`X`$ which are naturally associated to generalized theta line bundles on the moduli spaces $`U_X(r,0)`$ and play a key role in what follows.
###### Definition 2.1.
For every positive integers $`r`$ and $`k`$ and every $`N`$Pic$`{}_{}{}^{g1}(X)`$, the $`(r,k)`$ -$`\mathrm{𝑉𝑒𝑟𝑙𝑖𝑛𝑑𝑒}`$ bundle on $`J(X)`$ is the vector bundle
$$E_{r,k}(=E_{r,k}^N):=\mathrm{det}_{}𝒪(k\mathrm{\Theta }_N),$$
where $`\mathrm{det}:U_X(r,0)J(X)`$ is the determinant map. The dependence of the definition on the choice of $`N`$ is implicitly assumed, but not emphasized by the notation. Recall that by convention 2.2 we will (and it is enough) to assume that $`N`$ is a theta characteristic. Also, most of the time the rank $`r`$ is fixed and we refer to $`E_{r,k}`$ as the $`\mathrm{𝑙𝑒𝑣𝑒𝑙}k`$ Verlinde bundle.
###### Remark 2.2.
The fibers of the determinant map are the moduli spaces of vector bundles of fixed determinant $`SU_X(r,L)`$ with $`L`$Pic$`{}_{}{}^{0}(X)`$. The restriction of $`𝒪(k\mathrm{\Theta }_N)`$ to such a fiber is exactly $`^k`$, where $``$ is the determinant bundle. Since on $`SU_X(r)`$ the dualizing sheaf is isomorphic to $`^{2r}`$, by the rational singularities version of the Kodaira vanishing theorem these have no higher cohomology and it is clear that $`E_{r,k}`$ is a vector bundle of rank $`s_{r,k}:=h^0(SU_X(r),^k)`$. The fiber of $`E_{r,k}`$ at a point $`L`$ is naturally isomorphic to the Verlinde vector space $`H^0(SU_X(r,L),^k)`$ and this justifies the terminology ”Verlinde bundle” that we are using.
Much of the study of the vector bundles $`E_{r,k}`$ is governed by the fact that they decompose very nicely when pulled back via multiplication by $`r`$. To see this, first recall from , §2 and §4, that there is a cartesian diagram:
where $`\tau `$ is the tensor product of vector bundles, $`p_2`$ is the projection on the second factor and $`r_J`$ is multiplication by $`r`$. The top and bottom maps are étale covers of degree $`r^{2g}`$ and one finds in the formula:
(2)
$$\tau ^{}𝒪(\mathrm{\Theta }_N)𝒪_J(r\mathrm{\Theta }_N).$$
Using the notation $`V_{r,k}:=H^0(SU_X(r),^k)`$, we have the following simple but very important
###### Lemma 2.3.
$`r_J^{}E_{r,k}V_{r,k}𝒪_J(kr\mathrm{\Theta }_N).`$
###### Proof.
By the push-pull formula (see , III.9.3) and (2) we obtain:
$$r_J^{}E_{r,k}r_J^{}\mathrm{det}_{}𝒪(k\mathrm{\Theta }_N)p_{2}^{}{}_{}{}^{}\tau ^{}𝒪(k\mathrm{\Theta }_N)$$
$$p_{2}^{}{}_{}{}^{}(^k𝒪_J(kr\mathrm{\Theta }_N))V_{r,k}𝒪_J(kr\mathrm{\Theta }_N).$$
An immediate consequence of this property is the following (recall that a vector bundle is called *polystable* if it decomposes as a direct sum of stable bundles of the same slope):
###### Corollary 2.4.
$`E_{r,k}`$ is an ample vector bundle, polystable with respect to any polarization on $`J(X)`$.
###### Proof.
Both properties can be checked up to finite covers (see e.g. §3.2) and they are obvious for $`r_J^{}E_{r,k}`$. ∎
For future reference it is also necessary to study how the bundles $`E_{r,k}`$ behave under the Fourier-Mukai transform (recall the definitions from 2.3).
###### Lemma 2.5.
$`E_{r,k}`$ satisfies I.T. with index 0 and so $`\widehat{E_{r,k}}`$ is a vector bundle of rank $`h^0(U_X(r,0),𝒪(k\mathrm{\Theta }_N))`$ satisfying I.T. with index $`g`$.
###### Proof.
Using the usual identification between Pic$`{}_{}{}^{0}(J(X))`$ and Pic$`{}_{}{}^{0}(X)`$, we have to show that $`H^i(E_{r,k}P)=0`$ for all $`P`$Pic$`{}_{}{}^{0}(X)`$ and all $`i>0`$. But $`H^i(E_{r,k}P)`$ is a direct summand in $`H^i(r_{J}^{}{}_{}{}^{}r_J^{}(E_{r,k}P))`$ and so it is enough to have the vanishing of $`H^i(r_J^{}(E_{r,k}P))`$. This is obvious by the formula in 2.3. By 1.6(5) we have:
$$rk(\widehat{E_{r,k}})=h^0(E_{r,k})=h^0(U_X(r,0),𝒪(k\mathrm{\Theta }_N)).$$
The last statement follows also from 1.6(5). ∎
By Serre duality, one obtains in a similar way the statement:
###### Lemma 2.6.
$`E_{r,k}^{}`$ satisfies I.T. with index $`g`$ and so $`\widehat{E_{r,k}^{}}`$ is a vector bundle of rank $`h^0(U_X(r,0),𝒪(k\mathrm{\Theta }_N))`$ satisfying I.T. with index 0.
The Verlinde bundles are particularly easy to compute when $`k`$ is a multiple of the rank $`r`$, but note that for general $`k`$ the decomposition of $`E_{r,k}`$ into stable factors is not so easy to describe.
###### Proposition 2.7.
For all $`m1`$
$$E_{r,mr}\underset{s_{r,mr}}{}𝒪_J(m\mathrm{\Theta }_N)(V_{r,mr}𝒪_J(m\mathrm{\Theta }_N)).$$
###### Proof.
By 2.3 we have
$$r_J^{}E_{r,mr}V_{r,mr}𝒪_J(mr^2\mathrm{\Theta }_N).$$
Since $`N`$ is a theta characteristic, $`\mathrm{\Theta }_N`$ is symmetric and so $`r_J^{}𝒪(m\mathrm{\Theta }_N)𝒪_J(mr^2\mathrm{\Theta }_N)`$. We get
$$r_J^{}E_{r,mr}r_J^{}(V_{r,mr}𝒪(m\mathrm{\Theta }_N)).$$
Note now that the diagram giving 2.3 is equivariant with respect to the $`X_r`$, the group of $`r`$-torsion points of $`J(X)`$, if we let $`X_r`$ act on $`SU_X(r)\times J(X)`$ by $`(\xi ,(F,L))(F,L\xi )`$, by twisting on $`U_X(r,0)`$, by translation on the left $`J(X)`$ and trivially on the right $`J(X)`$. Also the action of $`X_r`$ on $`\tau ^{}𝒪(mr\mathrm{\Theta }_N)^{mr}𝒪_J(mr^2\mathrm{\Theta }_N)`$ is on $`𝒪_J(mr^2\mathrm{\Theta }_N)`$ the same as the natural (pullback) action. By chasing the diagram we see then that the vector bundle isomorphism above is equivariant with respect to the natural $`X_r`$ action on both sides. Since we have an an induced isomorphism:
$$r_{J}^{}{}_{}{}^{}r_J^{}E_{r,mr}r_{J}^{}{}_{}{}^{}r_J^{}(V_{r,mr}𝒪(m\mathrm{\Theta }_N))$$
and $`E_{r,mr}`$ and $`V_{r,mr}𝒪_J(m\mathrm{\Theta }_N)`$ are both eigenbundles with respect to the trivial character, the lemma follows. ∎
The rest of the section will be devoted to a further study of these bundles in the case $`k=1`$, where more tools are available. The main result is that $`E_{r,1}`$ is a simple vector bundle and this fact will be exploited in the next section. We show this after proving a very simple lemma, to the effect that twisting by $`r`$-torsion line bundles does not change $`E_{r,1}`$.
###### Lemma 2.8.
$`E_{r,1}P_\xi E_{r,1}`$ for any $`r`$-torsion line bundle $`P_\xi `$ on $`J(X)`$ corresponding to an $`r`$-torsion $`\xi \mathrm{Pic}^0(X)`$ by the usual identification.
###### Proof.
By definition and the projection formula one has
$$E_{r,1}P_\xi \mathrm{det}_{}(𝒪(\mathrm{\Theta }_N)\mathrm{det}^{}P_\xi )\mathrm{det}_{}𝒪(\mathrm{\Theta }_{N\xi }),$$
where the last isomorphism is an application of 1.3. Now the following commutative diagram:
shows that $`\mathrm{det}_{}𝒪(\mathrm{\Theta }_N)\mathrm{det}_{}𝒪(\mathrm{\Theta }_{N\xi })`$, which is exactly the statement of the lemma. ∎
###### Proposition 2.9.
$`E_{r,1}`$ is a simple vector bundle.
###### Proof.
This follows from a direct computation of the number of endomorphisms of $`E_{r,1}`$. By lemma 2.3 and the Verlinde formula at level $`1`$, $`r_J^{}E_{r,1}\underset{r^g}{}𝒪_J(r\mathrm{\Theta }_N)`$. Then:
$$h^0(r_{J}^{}{}_{}{}^{}r_J^{}(E_{r,1}^{}E_{r,1}))=h^0(r_J^{}(E_{r,1}^{}E_{r,1}))$$
$$=h^0((\underset{r^g}{}𝒪_J(r\mathrm{\Theta }_N))(\underset{r^g}{}𝒪_J(r\mathrm{\Theta }_N)))=r^{2g}.$$
On the other hand, since $`r_J`$ is a Galois cover with Galois group $`X_r`$, we have the formula $`r_{J}^{}{}_{}{}^{}𝒪_J\underset{\xi X_r}{}P_\xi `$. Combined with lemma 2.8, this gives
$$r_{J}^{}{}_{}{}^{}r_J^{}(E_{r,1}^{}E_{r,1})\underset{\xi X_r}{}E_{r,1}^{}E_{r,1}P_\xi \underset{r^{2g}}{}E_{r,1}^{}E_{r,1}.$$
The two relations imply that $`h^0(E_{r,1}^{}E_{r,1})=1`$, so $`E_{r,1}`$ is simple. ∎
###### Remark 2.10.
Using an argument similar to the one given above, plus the Verlinde formula, it is not hard to see that the bundles $`E_{r,k}`$ are not simple for $`k2`$. In the case when $`k`$ is a multiple of $`r`$ they even decompose as direct sums of line bundles, as we have already seen in 2.7. This shows why some special results that we will obtain in the case $`k=1`$ do not admit straightforward extensions to higher $`k`$’s.
## 3. Stability of Fourier transforms and duality for generalized theta functions
One of the main features of the vector bundles $`E_{r,1}`$, already observed in the previous section, is that they are simple. We will see below that in fact they are even stable (with respect to any polarization on $`J(X)`$). This fact, combined with the fact that $`\widehat{𝒪_J(r\mathrm{\Theta }_N)}`$ is also stable (see Proposition 3.1 below), gives simple proofs of some results of Beauville-Narasimhan-Ramanan and Donagi-Tu . Note that for the proofs of these applications it is enough to use the simpleness of the vector bundles mentioned above.
###### Proposition 3.1.
$`\widehat{E_{1,r}}=\widehat{𝒪_J(r\mathrm{\Theta }_N)}`$ is stable with respect to any polarization on $`J(X)`$.
This is a consequence of a more general fact of independent interest, saying that this is indeed true for an arbitrary nondegenerate line bundle on an abelian variety. Recall that a line bundle $`A`$ on the abelian variety $`X`$ is called $`\mathrm{𝑛𝑜𝑛𝑑𝑒𝑔𝑒𝑛𝑒𝑟𝑎𝑡𝑒}`$ if $`\chi (A)0`$. By §16 this implies that there is a unique $`i`$ (the $`\mathrm{𝑖𝑛𝑑𝑒𝑥}`$ of $`A`$) such that $`H^i(A)0`$.
###### Proposition 3.2.
Let $`A`$ be a nondegenerate line bundle on an abelian variety $`X`$. The Fourier-Mukai transform $`\widehat{A}`$ is stable with respect to any polarization on $`X`$.
###### Proof.
Let’s begin by fixing a polarization on $`\widehat{X}`$, so that stability will be understood with respect to this polarization. Consider the isogeny defined by $`A`$
$$\begin{array}{cccc}\varphi _A:& X& & \mathrm{Pic}^0(X)\widehat{X}\\ & x& & t_x^{}AA^1\end{array}$$
If $`i`$ is the index of $`A`$, it follows from 1.6(3) that $`\varphi _A^{}\widehat{A}VA^1`$, where $`V:=H^i(A)`$. As we already mentioned in the proof of 2.4, by §3.2 this already implies that $`\widehat{A}`$ is polystable. On the other hand, by 1.6(5) the Fourier transform of any line bundle is simple, so $`\widehat{A}`$ must be stable. ∎
###### Remark 3.3.
It is worth noting that one can avoid the use of §3.2 quoted above and use only the easier fact that semistability is preserved by finite covers. More precisely, $`\widehat{A}`$ has to be semistable, but assume that it is not stable. Then we can choose a maximal destabilizing subbundle $`F\widehat{A}`$, which must obviously be semistable and satisfy $`\mu (F)=\mu (\widehat{A})`$. Again $`\varphi _A^{}F`$ must be semistable, with respect to the pull-back polarization. But $`\varphi _A^{}F\varphi _A^{}\widehat{A}VA^1`$ and by semistability this implies that $`\varphi _A^{}FV^{}A^1`$ with $`V^{}V`$. This situation is overruled by the presence of the action of Mumford’s theta-group $`𝒢(A)`$. Recall from that $`A`$ is endowed with a natural $`𝒢(A)`$-linearization of weight $`1`$. On the other hand, $`\varphi _A^{}F`$ has a natural $`K(A)`$-linearization which can be seen as a $`𝒢(A)`$-linearization of weight 0, where $`K(A)`$ is the kernel of the isogeny $`\varphi _A`$ above. By tensoring we obtain a weight $`1`$ $`𝒢(A)`$-linearization on $`V^{}𝒪_X\varphi _A^{}FA`$ and thus an induced weight 1 representation on $`V^{}`$. It is known though from that $`V`$ is the unique irreducible representation of $`𝒢(A)`$ up to isomorphism, so we get a contradiction.
We now return to the study of the relationship between the bundles $`E_{r,k}`$ and their Fourier transforms. First note that throughout the rest of the paper we will use the fact that $`J(X)`$ is canonically isomorphic to its dual $`\mathrm{Pic}^0(J(X))`$ and consequently we will use the same notation for both. Thus all the Fourier transforms should be thought of as coming from the dual via this identification, although since there is no danger of confusion this will not be visible in the notation. The fiber of $`E_{r,k}`$ over a point $`\xi J(X)`$ is $`H^0(SU_X(r,\xi ),^k)`$, while the fiber of $`\widehat{E_{k,r}}`$ over the same point is canonically isomorphic to $`H^0(J(X),E_{k,r}P_\xi )`$. By 1.3 the latter is isomorphic to $`H^0(U_X(k,0),𝒪(r\mathrm{\Theta }_{N\eta }))`$, where $`\eta ^r\xi `$ (it is easy to see that this does not depend on the choice of $`\eta `$). The strange duality conjecture (see §8 and §5) says that there is a canonical isomorphism:
$$H^0(SU_X(r,\xi ),^k)^{}H^0(U_X(k,0),𝒪(r\mathrm{\Theta }_{N\eta }))$$
which will be described more precisely later in this section. This suggests then that one should relate somehow the vector bundles $`E_{r,k}`$ and $`\widehat{E_{k,r}}`$ via the diagram:
The first proposition treats the case $`k=1`$ and establishes the fact that the dual of $`E_{r,1}`$ is nothing else but the Fourier transform of $`𝒪_J(r\mathrm{\Theta }_N)`$.
###### Proposition 3.4.
$`E_{r,1}^{}\widehat{E_{1,r}}.`$
###### Proof.
By the duality theorem 1.4, it is enough to show that $`\widehat{E_{r,1}^{}}(1_J)^{}E_{1,r}`$. But $`E_{1,r}`$ is just $`𝒪_J(r\mathrm{\Theta }_N)`$, which is symmetric since $`N`$ is a theta characteristic. So what we have to prove is
$$\widehat{E_{r,1}^{}}𝒪_J(r\mathrm{\Theta }_N).$$
Since $`r_J`$ is Galois with Galois group $`X_r`$, we have (see e.g (2.1)):
$$r_J^{}r_{J}^{}{}_{}{}^{}\widehat{E_{r,1}^{}}\underset{\xi X_r}{}t_\xi ^{}\widehat{E_{r,1}^{}}.$$
But translates commute with tensor products via the Fourier transform (cf. 1.6(4)) and so
$$t_\xi ^{}\widehat{E_{r,1}^{}}\widehat{E_{r,1}^{}P_\xi }E_{r,1}^{},$$
where as usual $`P_\xi `$ is the line bundle in Pic$`{}_{}{}^{0}(J(X))`$ that corresponds to $`\xi `$ and the last isomorphism follows from 2.8. Thus we get
$$r_J^{}r_{J}^{}{}_{}{}^{}\widehat{E_{r,1}^{}}\underset{r^{2g}}{}\widehat{E_{r,1}^{}}$$
and the idea is to compute this bundle in a different way, by using the behavior of the Fourier transform under isogenies. More precisely, by applying 1.6(1) we get the isomorphisms:
$$r_J^{}r_{J}^{}{}_{}{}^{}\widehat{E_{r,1}^{}}r_J^{}\widehat{r_J^{}E_{r,1}^{}}r_J^{}(V_{r,1}^{}𝒪_J(r\mathrm{\Theta }_N))^\widehat{}$$
$$\underset{r^g}{}r_J^{}\widehat{𝒪_J(r\mathrm{\Theta }_N)}\underset{r^{2g}}{}𝒪_J(r\mathrm{\Theta }_N).$$
The second isomorphism follows from 2.3, while the fourth follows from 1.6(3) and the Verlinde formula at level $`1`$. The outcome is the isomorphism
$$\underset{r^{2g}}{}\widehat{E_{r,1}^{}}\underset{r^{2g}}{}𝒪_J(r\mathrm{\Theta }_N).$$
Now $`\widehat{E_{r,1}^{}}`$ is simple since $`E_{r,1}`$ is simple, so the previous isomorphism implies the stronger fact that
$$\widehat{E_{r,1}^{}}𝒪_J(r\mathrm{\Theta }_N).$$
###### Example 3.5.
The relationship between the Chern character of a sheaf and that of its Fourier transform established in (1.18) allows us to easily compute the first Chern class of $`E_{r,1}`$ as a consequence of the previous proposition. More precisely, if $`\theta `$ is the class of a theta divisor on $`J(X)`$, we have
$$c_1(\widehat{𝒪_J(r\mathrm{\Theta }_N)})=\mathrm{ch}_1(\widehat{𝒪_J(r\mathrm{\Theta }_N)})=(1)PD_{2g2}(\mathrm{ch}_{g1}(𝒪_J(r\mathrm{\Theta }_N)))=$$
$$=(1)PD_{2g2}(r^{g1}\theta ^{g1}/(g1)!)=r^{g1}\theta .$$
where $`PD_{2g2}:H^{2g2}(J(X),𝐙)H^2(J(X),𝐙)`$ is the Poincaré duality. We get:
$$c_1(E_{r,1})=r^{g1}\theta .$$
We are able to prove the analogous fact for higher $`k`$’s only modulo multiplication by $`r`$:
###### Proposition 3.6.
$`r_J^{}E_{r,k}^{}r_J^{}\widehat{E_{k,r}}.`$
###### Proof.
We know that $`k_J^{}E_{k,r}\underset{s_{k,r}}{}𝒪_J(kr\mathrm{\Theta }_N)`$, so by 1.6(1) we obtain
$$k_{J}^{}{}_{}{}^{}\widehat{E_{k,r}}\underset{s_{k,r}}{}\widehat{𝒪_J(kr\mathrm{\Theta }_N)}.$$
As in the previous proposition, since the Galois group of $`k_J`$ is $`X_k`$, we have
$$k_J^{}k_{J}^{}{}_{}{}^{}\widehat{E_{k,r}}\underset{\xi X_k}{}t_\xi ^{}\widehat{E_{k,r}}\mathrm{and}\mathrm{so}(kr)_J^{}k_{J}^{}{}_{}{}^{}\widehat{E_{k,r}}\underset{\xi X_k}{}r_J^{}t_\xi ^{}\widehat{E_{k,r}}.$$
Moreover the isomorphisms above show as before that $`\widehat{E_{k,r}}`$ has to be semistable with respect to an arbitrary polarization. On the other by 1.6(3) we have
$$(kr)_J^{}k_{J}^{}{}_{}{}^{}\widehat{E_{r,k}}\underset{s_{k,r}}{}(kr)_J^{}\widehat{𝒪_J(kr\mathrm{\Theta }_N)}\underset{k^gr^gs_{k,r}}{}𝒪_J(kr\mathrm{\Theta }_N).$$
This gives us the isomorphism
$$\underset{\xi X_k}{}r_J^{}t_\xi ^{}\widehat{E_{k,r}}\underset{k^gr^gs_{k,r}}{}𝒪_J(kr\mathrm{\Theta }_N)$$
which in particular implies (recall semistability) that
$$r_J^{}\widehat{E_{k,r}}\underset{s_{r,k}}{}𝒪_J(kr\mathrm{\Theta }_N)r_J^{}E_{r,k}^{}.$$
The important fact that the last index of summation is $`s_{r,k}`$ follows from the well-known “symmetry” of the Verlinde formula, which is:
$$r^gs_{k,r}=k^gs_{r,k}.$$
The propositions above allow us to give alternative proofs of some results in and concerning duality between spaces of generalized theta functions. We first show how one can recapture a theorem of Donagi-Tu in the present context. The full version of the theorem (i.e. for arbitrary degree) can be obtained by the same method (cf. Section 6).
###### Theorem 3.7.
(, Theorem 1) For any $`L\mathrm{Pic}^0(X)`$ and any $`N\mathrm{Pic}^{g1}(X)`$, we have:
$$h^0(SU_X(k,L),^r)r^g=h^0(U_X(k,0),𝒪(r\mathrm{\Theta }_N))k^g.$$
###### Proof.
We will actually prove the following equality:
$$h^0(SU_X(r,L),^k)=h^0(U_X(k,0),𝒪(r\mathrm{\Theta }_N)).$$
The statement will then follow from the same symmetry of the Verlinde formula $`r^gs_{k,r}=k^gs_{r,k}`$ mentioned in the proof of 3.6. To this end we can use Proposition 3.6 to obtain $`\mathrm{rk}(E_{r,k}^{})=\mathrm{rk}(\widehat{E_{k,r}})`$. But on one hand
$$\mathrm{rk}(E_{r,k}^{})=h^0(SU_X(r,L),^k)$$
while on the other hand by 1.6(5)
$$\mathrm{rk}(\widehat{E_{k,r}})=h^0(J(X),E_{k,r})=h^0(U_X(k,0),𝒪(r\mathrm{\Theta }_N))$$
as required. ∎
It is worth mentioning that since we are assuming the Verlinde formula all throughout, an important particular case of the theorem above is:
###### Corollary 3.8.
(, Theorem 2) $`h^0(U_X(r,0),𝒪(\mathrm{\Theta }_N))=1.`$
As it is very well known, this fact is essential in setting up the strange duality. Turning to an application in this direction, the results above also lead to a simple proof of the strange duality at level $`1`$, which has been first given a proof in , Theorem 3. It is important though to emphasize again that here, unlike in the quoted paper, the Verlinde formula is granted. So the purpose of the next application is to show how, with the knowledge of the Verlinde numbers, the strange duality at level $`1`$ can simply be seen as the solution of a stability problem for vector bundles on $`J(X)`$. This may also provide a global method for understanding the conjecture for higher levels. A few facts in this direction will be mentioned at the end of this section (cf. Remark 3.13).
Before turning to the proof, we need the following general result, which is a globalization of §3 in .
###### Proposition 3.9.
Consider the tensor product map
$$\tau :U_X(r,0)\times U_X(k,0)U_X(kr,0)$$
and the map
$$\varphi :=\mathrm{det}\times \mathrm{det}:U_X(r,0)\times U_X(k,0)J(X)\times J(X).$$
Then
$$\tau ^{}𝒪(\mathrm{\Theta }_N)p_1^{}𝒪_J(k\mathrm{\Theta }_N)p_2^{}𝒪_J(r\mathrm{\Theta }_N)\varphi ^{}𝒫,$$
where $`𝒫`$ is a Poincaré line bundle on $`J(X)\times J(X)`$, normalized such that $`𝒫_{|\{0\}\times J(X)}𝒪_J`$.
###### Proof.
We will compare the restrictions of the two line bundles to fibers of the projections. First fix $`FU_X(k,0)`$. We have
$$\tau ^{}𝒪(\mathrm{\Theta }_N)|_{U_X(r,0)\times \{F\}}\tau _F^{}𝒪(\mathrm{\Theta }_N)𝒪(\mathrm{\Theta }_{FN}),$$
where $`\tau _F`$ is the map given by twisting with $`F`$. But by 1.3 one has
$$𝒪(\mathrm{\Theta }_{FN})𝒪(k\mathrm{\Theta }_N)\mathrm{det}^{}(\mathrm{detF})(𝒪(k\mathrm{\Theta }_N)\mathrm{det}^{}(𝒫_{|J(X)\times \{\mathrm{det}F\}})).$$
On the other hand obviously
$$p_1^{}𝒪_J(k\mathrm{\Theta }_N)p_2^{}𝒪_J(r\mathrm{\Theta }_N)\varphi ^{}𝒫|_{U_X(r,0)\times \{F\}}𝒪(k\mathrm{\Theta }_N)\mathrm{det}^{}(\mathrm{det}F).$$
Let’s now fix $`EU_X(r,0)`$ such that $`\mathrm{det}E𝒪_X`$. Using the same 1.3 we get
$$\tau ^{}𝒪(\mathrm{\Theta }_N)|_{\{E\}\times U_X(k,0)}𝒪(\mathrm{\Theta }_{EN})𝒪(r\mathrm{\Theta }_N)$$
and we also have
$$p_1^{}𝒪_J(k\mathrm{\Theta }_N)p_2^{}𝒪_J(r\mathrm{\Theta }_N)\varphi ^{}𝒫|_{\{E\}\times U_X(k,0)}$$
$$𝒪(r\mathrm{\Theta }_N)\mathrm{det}^{}(𝒫_{|\{𝒪_X\}\times J(X)})𝒪(r\mathrm{\Theta }_N).$$
The desired isomorphism follows now from the see-saw principle(see e.g. , I.5.6). ∎
Theorem 3.7 tells us that there is essentially a unique nonzero section
$$sH^0(U_X(kr,0),𝒪(\mathrm{\Theta }_N)),$$
which induces via $`\tau `$ a nonzero section
$$tH^0(U_X(r,0)\times U_X(k,0),\tau ^{}𝒪(\mathrm{\Theta }_N)$$
$$H^0(U_X(r,0)\times U_X(k,0),p_1^{}𝒪_J(k\mathrm{\Theta }_N)p_2^{}𝒪_J(r\mathrm{\Theta }_N)\varphi ^{}𝒫).$$
But notice that from the projection formula we get
$$\varphi _{}(p_1^{}𝒪_J(k\mathrm{\Theta }_N)p_2^{}𝒪_J(r\mathrm{\Theta }_N)\varphi ^{}𝒫)$$
$$\varphi _{}(p_1^{}𝒪_J(k\mathrm{\Theta }_N)p_2^{}𝒪_J(r\mathrm{\Theta }_N))𝒫p_1^{}E_{r,k}p_2^{}E_{k,r}𝒫,$$
so $`t`$ induces a section (denoted also by $`t`$):
$$0tH^0(J(X)\times J(X),p_1^{}E_{r,k}p_2^{}E_{k,r}𝒫)H^0(E_{r,k}\widehat{E_{k,r}}).$$
This is nothing else but a globalization of the section defining the strange duality morphism, as explained in §5 (simply because for any $`\xi `$Pic$`{}_{}{}^{0}(X)`$ the restriction of $`\tau `$ in 3.9 to $`SU_X(r,\xi )\times U_X(k,0)`$ is again the tensor product map
$$\tau :SU_X(r,\xi )\times U_X(k,0)U_X(kr,0)$$
and $`\tau ^{}𝒪(\mathrm{\Theta }_N)^k𝒪(r\mathrm{\Theta }_N)`$). In other words, $`t`$ corresponds to a morphism of vector bundles
$$SD:E_{r,k}^{}\widehat{E_{k,r}}$$
which fiberwise is exactly the strange duality morphism
$$H^0(SU_X(r,\xi ),^k)^{}\stackrel{SD}{}H^0(U_X(k,0),𝒪(r\mathrm{\Theta }_{N\eta })),$$
where $`\eta ^r\xi `$. Since this global morphism collects together all the strange duality morphisms as we vary $`\xi `$, the strange duality conjecture is equivalent to $`SD`$ being an isomorphism.
###### Conjecture 3.10.
$`SD:E_{r,k}^{}\widehat{E_{k,r}}`$ is an isomorphism of vector bundles.
In this context Proposition 3.6 can be seen as a weak form of “global” evidence for the conjecture in the case $`k2`$. For $`k=1`$ it can now be easily proved.
###### Theorem 3.11.
$`SD:E_{r,1}^{}\widehat{E_{1,r}}`$ is an isomorphism.
###### Corollary 3.12.
(cf. , Theorem 3) The level 1 strange duality morphism
$$H^0(SU_X(r),)^{}\stackrel{SD}{}H^0(J(X),𝒪(r\mathrm{\Theta }_N))$$
is an isomorphism.
###### Proof.
(of 3.11) All the ingredients necessary for proving this have been discussed above: by 3.1 and 3.4, the bundles $`E_{r,1}^{}`$ and $`\widehat{E_{1,r}}`$ are isomorphic and stable. This means that $`SD`$ is essentially the unique nonzero morphism between them, and it must be an isomorphism. ∎
###### Remark 3.13.
One step towards 3.10 is a better understanding of the properties of the kernel $`F`$ of $`SD`$. A couple of interesting remarks in this direction can already be made. Since $`E_{r,k}`$ and $`\widehat{E_{k,r}}`$ are polystable of the same slope, the same will be true about $`F`$. On the other hand some simple calculus involving 1.3 and 1.6(4) shows that $`F`$ gets multiplied by a line bundle in $`\mathrm{Pic}^0(J(X))`$ when we translate it, so in the language of Mukai (e.g. §3) it is a semi-homogeneous vector bundle (although clearly not homogeneous, i.e. not fixed by all translations).
## 4. A generalization of Raynaud’s examples
In this section we would like to discuss a generalization of the examples of base points of the theta linear system $`||`$ on $`SU_X(m)`$ constructed by Raynaud in . For a survey of this circle of ideas the reader can consult , §2. Let us recall here (see §2) only that for a semistable vector bundle $`E`$ of rank $`m`$ to induce a base point of $`||`$ it is sufficient that it satisfies the property
$$0\mu (E)g1\mathrm{and}h^0(EL)0\mathrm{for}L\mathrm{Pic}^0(X)\mathrm{generic}.$$
The examples of Raynaud are essentially the restrictions of $`\widehat{E_{1,r}^{}}=\widehat{𝒪_J(r\mathrm{\Theta }_N)}`$ to some embedding of the curve $`X`$ in the Jacobian. We will generalize this by considering the Fourier transform of higher level Verlinde bundles $`\widehat{E_{k,r}^{}}`$ with $`k2`$.
For simplicity, let’s fix $`k`$ and $`r`$ and denote $`F:=\widehat{E_{k,r}^{}}`$. This is a vector bundle by lemma 2.6. Consider also an arbitrary embedding
$$j:XJ(X)$$
and denote by $`E`$ the restriction $`F_{|X}`$.
###### Proposition 4.1.
E is a semistable vector bundle.
###### Proof.
This follows basically from the proof of Proposition 3.6. One can see in a completely analogous way that:
$$r_J^{}Fr_J^{}E_{r,k}\underset{s_{r,k}}{}𝒪_J(kr\mathrm{\Theta }_N).$$
Now if we consider $`Y`$ to be the preimage of $`X`$ by $`r_J`$, this shows that $`r_J^{}F_{|Y}`$ is semistable and so $`F_{|X}`$ is semistable. ∎
###### Proposition 4.2.
There is an embedding of $`X`$ in $`J(X)`$ such that $`E=F_{|X}`$ satisfies $`H^0(EL)0`$ for $`L\mathrm{Pic}^0(X)`$ generic.
###### Proof.
The proof goes like in (3.1) and we repeat it here for convenience: choose $`U`$ a nonempty open subset of $`J(X)`$ on which $`(1_J)^{}E_{k,r}`$ is trivial. By Mukai’s duality theorem (2.2) we know that $`\widehat{F}(1_J)^{}E_{k,r}^{}`$, so there exists a nonzero section $`s\mathrm{\Gamma }(p_{2}^{}{}_{}{}^{1}(U),p_1^{}F𝒫)`$. Choose now an $`x_0J(X)`$ such that $`s_{|\{x_0\}\times U}0`$ and consider an embedding of $`X`$ in $`J(X)`$ passing through $`x_0`$. The image of $`s`$ in $`\mathrm{\Gamma }(X\times U,p_1^{}F𝒫)`$ is nonzero and this implies that $`H^0(EL)0`$ for $`L`$ generic. ∎
We are only left with computing the invariants of $`E`$.
###### Proposition 4.3.
The rank and the slope of $`E`$ are given by:
$$\mathrm{rk}(E)=s_{r,k}=h^0(SU_X(r),^k)\mathrm{and}\mu (E)=gk/r.$$
###### Proof.
By 1.6(5) we have
$$\mathrm{rk}(F)=\mathrm{rk}(\widehat{E_{k,r}^{}})=h^0(E_{k,r})=h^0(U_X(k,0),𝒪(r\mathrm{\Theta }_N)).$$
But from the proof of 3.7 we know that $`h^0(U_X(k,0),𝒪(r\mathrm{\Theta }_N))=h^0(SU_X(r,0),^k)`$ and so $`\mathrm{rk}(F)=s_{r,k}`$. To compute the slope of $`E`$, first notice that by the proof of Proposition 4.1 we know that
$$r^{2g}\mu (E)=\mathrm{deg}(𝒪_J(kr\mathrm{\Theta }_N)_{|Y}).$$
But $`r_J^{}𝒪_J(kr\mathrm{\Theta }_N)𝒪_J(kr^3\mathrm{\Theta }_N)`$, so
$$r^2\mathrm{deg}(𝒪_J(kr\mathrm{\Theta }_N)_{|Y})=\mathrm{deg}(r_J^{}𝒪_J(kr\mathrm{\Theta }_N)_{|Y})=r^{2g}\mathrm{deg}(𝒪_J(kr\mathrm{\Theta }_N)_{|X})=r^{2g+1}kg.$$
Combining the two equalities we get
$$\mu (E)=\mathrm{deg}(𝒪_J(kr\mathrm{\Theta }_N)_{|Y})/r^{2g}=gk/r.$$
In conclusion, for each $`r`$ and $`k`$ we obtain a semistable vector bundle $`E`$ on $`X`$ of rank equal to the Verlinde number $`s_{r,k}`$ and of slope $`gk/r`$, satisfying the property that $`H^0(EL)0`$ for $`L`$ generic in Pic$`{}_{}{}^{0}(X)`$. So as long as $`k<r`$ and $`r`$ divides $`gk`$ we obtain new examples of base points for $`||`$ on the moduli spaces $`SU_X(s_{r,k})`$, as explained above. Raynaud’s examples correspond to the case $`k=1`$. See also for a study of a different kind of examples of such base points and for bounds on the dimension of the base locus of $`||`$.
## 5. Linear series on $`U_X(r,0)`$
The main application of the Verlinde vector bundles concerns the global generation and normal generation of line bundles on $`U_X(r,0)`$. The specific goal is to give effective bounds for multiples of the generalized theta line bundles that satisfy the properties mentioned above. In this direction analogous results for $`SU_X(r)`$ will be used. The starting point is the following general result:
###### Proposition 5.1.
Let $`f:XY`$ be a flat morphism of projective schemes, with reduced fibers. Let $`L`$ be a line bundle on $`X`$ and denote $`E:=f_{}L`$. Assume that if $`X_y`$ denotes the fiber of $`f`$ over $`yY`$ the following conditions hold:
(i) $`h^1(L)=0`$
(ii) $`h^i(L_{|X_y})=0,yY,i>0`$.
Then $`L`$ is globally generated as long as $`L_{|X_y}`$ is globally generated for all $`yY`$ and $`E`$ is globally generated.
###### Proof.
Start with $`xX`$ and consider $`y=f(x)`$. By $`(i)`$ we have the exact sequence on $`X`$:
$$0H^0(L_{X_y})H^0(L)H^0(L_{|X_y})H^1(L_{X_y})0.$$
The global generation of $`L_{|X_y}`$ implies that there exists a section $`sH^0(L_{|X_y})`$ such that $`s(x)0`$. We would like to lift $`s`$ to some $`\overline{s}`$, so it is enough to prove that $`H^1(L_{X_y})=0`$. The fibers of $`f`$ are reduced, so $`_{X_y}f^{}_{\{y\}}`$. Condition (ii) implies, by the base change theorem, that $`R^if_{}L=0`$ for all $`i>0`$ and so by the projection formula we also get $`R^if_{}(L_{X_y})=0`$ for all $`i>0`$. The Leray spectral sequence then gives $`H^i(E)H^i(L)`$ and $`H^i(E_{\{y\}})H^i(L_{Xy})`$ for all $`i>0`$. The first isomorphism implies that there is an exact sequence:
$$0H^0(E_{\{y\}})H^0(E)\stackrel{ev_y}{}H^0(E_y)H^1(E_{\{y\}})0.$$
But $`E`$ is globally generated, which means that $`ev_y`$ is surjective. This implies the vanishing of $`H^1(E_{\{y\}})`$, which by the second isomorphism is equivalent to $`H^1(L_{X_y})=0`$. ∎
The idea is to apply this result to the situation when the map is $`\mathrm{det}:U_X(r,0)J(X)`$, $`L=𝒪(k\mathrm{\Theta }_N)`$ and $`E=E_{r,k}`$. Modulo detecting for what $`k`$ global generation is attained, the conditions of the proposition are satisfied. The fiberwise global generation problem (i.e. the $`SU_X(r)`$ case) has been given some effective solutions in the literature. The most recent is the author’s result in , improving earlier bounds of Le Potier and Hein . We show in §4 that $`^k`$ on $`SU_X(r)`$ is globally generated if $`k\frac{(r+1)^2}{4}`$. We now turn to the effective statement for $`E_{r,k}`$ and in this direction we make essential use of Pareschi’s cohomological criterion described in Section 1.
###### Proposition 5.2.
$`E_{r,k}`$ is globally generated if and only if $`kr+1`$.
###### Proof.
The trick is to write $`E_{r,k}`$ as $`E_{r,k}𝒪_J(\mathrm{\Theta }_N)𝒪_J(\mathrm{\Theta }_N)`$. Denote $`E_{r,k}𝒪_J(\mathrm{\Theta }_N)`$ by $`F`$. Pareschi’s criterion 1.7 says in our case that $`E_{r,k}`$ is globally generated as long as the condition
$$h^i(F\alpha )=0,\alpha \mathrm{Pic}^0(X),i>0$$
is satisfied (in fact under this assumption $`FA`$ will be globally generated for every ample line bundle $`A`$). Arguing as usual, $`h^i(F\alpha )=0`$ is implied by $`h^i(r_J^{}(F\alpha ))=0`$. We have
$$r_J^{}Fr_J^{}E_{r,k}r_J^{}𝒪_J(\mathrm{\Theta }_N)𝒪_J((krr^2)\mathrm{\Theta }_N)$$
and this easily gives the desired vanishing for $`kr+1`$. On the other hand from 2.7 we know that $`E_{r,r}\underset{s_{r,r}}{}𝒪_J(\mathrm{\Theta }_N)`$, which is clearly not globally generated. This shows that the bound is optimal. ∎
Combining all these we obtain effective bounds for global generation on $`U_X(r,0)`$ in terms of the analogous bounds on $`SU_X(r)`$. We prefer to state the general result in a non-effective form though, in order to emphasize that it applies algorithmically (but see the corollaries for effective statements):
###### Theorem 5.3.
$`𝒪(k\mathrm{\Theta }_N)`$ is globally generated on $`U_X(r,0)`$ as long as $`kr+1`$ and $`^k`$ is globally generated on $`SU_X(r)`$. Moreover, $`𝒪(r\mathrm{\Theta }_N)`$ is not globally generated.
###### Proof.
The first part follows by puting together 5.2 and 5.1 in our particular setting. To prove that $`𝒪(r\mathrm{\Theta }_N)`$ is not globally generated, let us begin by assuming the contrary. Then the restriction $`^r`$ of $`𝒪(r\mathrm{\Theta }_N)`$ to any of the fibers $`SU_X(r,L)`$ is also globally generated.
Choose in particular a line bundle $`L`$ in the support of $`\mathrm{\Theta }_N`$ on $`J(X)`$. Restriction to the fiber gives the following long exact sequence on cohomology:
$$0H^0(𝒪(r\mathrm{\Theta }_N)_{SU_X(r,L)})H^0(𝒪(r\mathrm{\Theta }_N))\stackrel{𝛼}{}$$
$$H^0(^r)H^1(𝒪(r\mathrm{\Theta }_N)_{SU_X(r,L)})0.$$
As in the proof of 5.1, this sequence can be written in terms of the cohomology of $`E_{r,r}`$:
$$0H^0(E_{r,r}_{\{L\}})H^0(𝒪(r\mathrm{\Theta }_N))\stackrel{𝛼}{}H^0(^r)H^1(E_{r,r}_{\{L\}})0.$$
The assumption on $`𝒪(r\mathrm{\Theta }_N)`$ ensures the fact that the map $`\alpha `$ in the sequence above is nonzero, and as a result
$$h^1(E_{r,r}_{\{L\}})<h^0(^r)=s_{r,r}.$$
On the other hand we use again the fact 2.7 that $`E_{r,r}`$ is isomorphic to $`\underset{s_{r,r}}{}𝒪_J(\mathrm{\Theta }_N)`$. The additional hypothesis that $`L\mathrm{\Theta }_N`$ says then that
$$h^1(E_{r,r}_{\{L\}})=s_{r,r}$$
which is a contradiction. ∎
As suggested above, by combining this with the effective result on $`SU_X(r)`$ given in §4, we get:
###### Corollary 5.4.
(cf. (5.3)) $`𝒪(k\mathrm{\Theta }_N)`$ is globally generated on $`U_X(r,0)`$ for
$$k\frac{(r+1)^2}{4}.$$
It is important, as noted in the introduction, to emphasize the fact that the $`SU_X(r)`$ bound may still allow for improvement. Thus the content and formulation of the theorem certainly go beyond this corollary. For moduli spaces of vector bundles of rank $`2`$ and $`3`$ though, in view of the second part of 5.3, we actually have optimal results (see also (5.4)).
###### Corollary 5.5.
(i) $`𝒪(3\mathrm{\Theta }_N)`$ is globally generated on $`U_X(2,0)`$.
(ii) $`𝒪(4\mathrm{\Theta }_N)`$ is globally generated on $`U_X(3,0)`$.
These are natural extensions of the fact that $`𝒪(2\mathrm{\Theta }_N)`$ is globally generated on $`J(X)U_X(1,0)`$ (see e.g. , Theorem 2, p.317). In §5 we also state some questions and conjectures about optimal bounds in general.
###### Remark 5.6.
A similar technique can be applied to study the base point freeness of more general linear series on $`U_X(r,0)`$. This is done in (5.9).
###### Remark 5.7.
A result analogous to 5.1, combined with a more careful study of the cohomological properties of $`E_{r,k}`$, gives information about effective separation of points by the linear series $`|k\mathrm{\Theta }_N|`$. We will not insist on this aspect in the present article.
In the same spirit of studying properties of linear series on $`U_X(r,0)`$ via vector bundle techniques, one can look at multiplication maps on spaces of sections and normal generation. The Verlinde bundles are again an essential tool. The underlying theme is the study of surjectivity of the multiplication map
$$H^0(𝒪(k\mathrm{\Theta }_N))H^0(𝒪(k\mathrm{\Theta }_N))\stackrel{\mu _k}{}H^0(𝒪(2k\mathrm{\Theta }_N)).$$
To this respect we have to start by assuming that $`k`$ is already chosen such that $`^k`$ and $`E_{r,k}`$ are globaly generated (in particular $`kr+1`$). As proved in Theorem 5.3, this also induces the global generation of $`𝒪(k\mathrm{\Theta }_N)`$. The method will be to look at the kernels of various multiplication maps–in the spirit of for example–and study their cohomology vanishing properties. Let $`M_k`$ on $`U_X(r,0)`$ and $`M_{r,k}`$ on $`J(X)`$ be the vector bundles defined by the sequences:
(3)
$$0M_kH^0(𝒪(k\mathrm{\Theta }_N))𝒪𝒪(k\mathrm{\Theta }_N)0$$
and
(4)
$$0M_{r,k}H^0(E_{r,k})𝒪_JE_{r,k}0.$$
By twisting (3) with $`𝒪(k\mathrm{\Theta }_N)`$ and taking cohomology, it is clear that the surjectivity of $`\mu _k`$ is equivalent to $`H^1(M_k𝒪(k\mathrm{\Theta }_N))=0`$.
On the other hand, the global generation of $`𝒪(k\mathrm{\Theta }_N)`$ implies that the natural map $`\mathrm{det}^{}E_{r,k}𝒪(k\mathrm{\Theta }_N)`$ is surjective, so we can consider the vector bundle $`K`$ defined by the following sequence:
(5)
$$0K\mathrm{det}^{}E_{r,k}𝒪(k\mathrm{\Theta }_N)0.$$
###### Remark 5.8.
Fixing $`L`$ Pic$`{}_{}{}^{0}(X)`$, we can also look at the evaluation sequence for $`^k`$ on $`SU_X(r,L)`$:
$$0M_^kH^0(^k)𝒪_{SU_X}^k0.$$
The sequence (5) should be interpreted as globalizing this picture. It induces the above sequence when restricted to the fiber of the determinant map over each $`L`$.
The study of vanishing for $`M_{r,k}`$ and $`K`$ will be the key to obtaining the required vanishing for $`M_k`$. This is reflected in the top exact sequence in the following commutative diagram, obtained from (3),(4) and (5) as an application of the snake lemma:
We again state our result in part (b) of the following theorem in a form that allows algorithmic applications. The main ingredient is an effective normal generation bound for $`E_{r,k}`$, which is the content of part (a).
###### Theorem 5.9.
(a) The multiplication map
$$H^0(E_{r,k})H^0(E_{r,k})H^0(E_{r,k}^2)$$
is surjective for $`k2r+1`$.
(b) Under the global generation assumptions formulated above, the multiplication map
$$\mu _k:H^0(𝒪(k\mathrm{\Theta }_N))H^0(𝒪(k\mathrm{\Theta }_N))H^0(𝒪(2k\mathrm{\Theta }_N))$$
is surjective as long as the multiplication map $`H^0(^k)H^0(^k)H^0(^{2k})`$ on $`SU_X(r)`$ is surjective and $`k2r+1`$.
###### Proof.
(a) This is not hard to deal with when $`k`$ is a multiple of $`r`$, since we know from 2.7 that $`E_{r,k}`$ decomposes in a particularly nice way. To tackle the general case though, we have to appeal to 1.9. Concretely, we have to see precisely when the skew Pontrjagin product
$$E_{r,k}\widehat{}E_{r,k}E_{r,k}(1_J)^{}E_{r,k}$$
is globally generated, and since by the initial choice of a theta characteristic $`E_{r,k}`$ is symmetric, this is the same as the global generation of the usual Pontrjagin product $`E_{r,k}E_{r,k}`$. As an aside, recall from 1.9 that this would imply the surjectivity of all the multiplication maps
$$H^0(t_x^{}E_{r,k})H^0(E_{r,k})H^0(t_x^{}E_{r,k}E_{r,k})$$
for all $`xJ(X)`$.
The global generation of this Pontrjagin product is in turn another application of the general cohomological criterion 1.7 for vector bundles on abelian varieties. We first prove that $`E_{r,k}E_{r,k}`$ also has a nice decomposition when pulled back by an isogeny, namely this time by multiplication by $`2r`$. Denote by $`F`$ the Fourier transform $`\widehat{E_{r,k}}`$, so that $`\widehat{F}(1_J)^{}E_{r,k}E_{r,k}`$. Then we have the following isomorphisms:
$$E_{r,k}E_{r,k}\widehat{F}\widehat{F}\widehat{FF},$$
where the second one is obtained by the correspondence between the Potrjagin product and the tensor product via the Fourier-Mukai transform, as in 1.6(2). Next, as in the previous sections, we look at the behaviour of our bundle when pulled back via certain isogenies (cf. 1.6(1)):
(6)
$$k_{J}^{}{}_{}{}^{}(E_{r,k}E_{r,k})\widehat{k_J^{}Fk_J^{}F}.$$
In 3.6 we proved that
$$k_J^{}Fk_J^{}\widehat{E_{r,k}}k_J^{}E_{k,r}^{}\underset{s_{k,r}}{}𝒪_J(kr\mathrm{\Theta }_N).$$
and by plugging this into (6) we obtain
$$k_{J}^{}{}_{}{}^{}(E_{r,k}E_{r,k})(\underset{s_{k,r}}{}𝒪_J(kr\mathrm{\Theta }_N)\underset{s_{k,r}}{}𝒪_J(kr\mathrm{\Theta }_N))^\widehat{}$$
$$\widehat{\underset{s_{k,r}^2}{}𝒪_J(2kr\mathrm{\Theta }_N)}.$$
Finally we apply $`(2r)_J^{}k_J^{}`$ to both sides of the isomorphism above and use the behavior of the Fourier transform of a line bundle when pulled back via the isogeny that it determines (see 1.6(3)). Since $`E_{r,k}E_{r,k}`$ is a direct summand in $`k_J^{}k_{J}^{}{}_{}{}^{}(E_{r,k}E_{r,k})`$, we obtain the desired decomposition:
(7)
$$(2r)_J^{}(E_{r,k}E_{r,k})𝒪_J(2kr\mathrm{\Theta }_N).$$
This allows us to apply a trick analogous to the one used in the proof of 5.2. Namely (7) implies that
$$(2r)_J^{}(E_{r,k}E_{r,k}𝒪_J(\mathrm{\Theta }_N))𝒪_J((2kr4r^2)\mathrm{\Theta }_N).$$
Thus if we denote by $`U_{r,k}`$ the vector bundle $`E_{r,k}E_{r,k}𝒪_J(\mathrm{\Theta }_N)`$ we clearly have:
$$h^i(U_{r,k}\alpha )=0,\alpha \mathrm{Pic}^0(X),i>0\mathrm{and}k2r+1.$$
Pareschi’s criterion 1.7 immediately gives then that $`E_{r,k}E_{r,k}`$ is globally generated for $`k2r+1`$, since
$$E_{r,k}E_{r,k}U_{r,k}𝒪_J(\mathrm{\Theta }_N).$$
(b) We will show the vanishing of $`H^1(M_k𝒪(k\mathrm{\Theta }_N))`$. By the top sequence in the diagram preceding the theorem, it is enough to prove that
$$H^1(K𝒪(k\mathrm{\Theta }_N))=0\mathrm{and}H^1(\mathrm{det}^{}M_{r,k}𝒪(k\mathrm{\Theta }_N))=0.$$
First we prove the vanishing of $`H^1(K𝒪(k\mathrm{\Theta }_N))`$. The key point is to identify the pull-back of $`K`$ by the étale cover $`\tau `$ in the diagram
described in section 3. In the pull-back sequence
$$0\tau ^{}K\tau ^{}\mathrm{det}^{}E_{r,k}\tau ^{}𝒪(k\mathrm{\Theta }_N)0$$
we can identify $`\tau ^{}\mathrm{det}^{}E_{r,k}`$ with $`p_2^{}r_{J}^{}{}_{}{}^{}E_{r,k}`$ and $`\tau ^{}𝒪(k\mathrm{\Theta }_N)`$ with $`^k𝒪_J(kr\mathrm{\Theta }_N)`$. In other words we have the exact sequence
$$0\tau ^{}KH^0(^k)𝒪_J(kr\mathrm{\Theta }_N)^k𝒪_J(kr\mathrm{\Theta }_N)0,$$
which shows that the following isomorphism holds (cf. 5.8):
$$\tau ^{}KM_^k𝒪_J(kr\mathrm{\Theta }_N).$$
Finally we obtain the isomorphism
$$\tau ^{}(K𝒪(k\mathrm{\Theta }_N))(M_^k^k)𝒪_J(2kr\mathrm{\Theta }_N).$$
Certainly by the argument mentioned earlier the surjectivity of the multiplication map $`H^0(^k)H^0(^k)H^0(^{2k})`$ is also equivalent to $`H^1(M_^k^k)=0`$. The required vanishing is then an easy application of the Künneth formula.
The next step is to prove the vanishing of $`H^1(\mathrm{det}^{}M_{r,k}𝒪(k\mathrm{\Theta }_N))`$. From the projection formula we know that
$$R^i\mathrm{det}_{}(\mathrm{det}^{}M_{r,k}𝒪(k\mathrm{\Theta }_N))M_{r,k}R^i\mathrm{det}_{}𝒪(k\mathrm{\Theta }_N)=0\mathrm{for}\mathrm{all}i>0$$
since obviously $`R^i\mathrm{det}_{}𝒪(k\mathrm{\Theta }_N)=0`$ for all $`i>0`$. The Leray spectral sequence reduces then our problem to proving the vanishing $`H^1(M_{r,k}E_{r,k})=0`$, which is basically equivalent to the surjectivity of the multiplication map
$$H^0(E_{r,k})H^0(E_{r,k})H^0(E_{r,k}^2).$$
This is the content of part (a). ∎
###### Corollary 5.10.
For $`k`$ as in 5.9, $`𝒪(k\mathrm{\Theta }_N)`$ is very ample.
###### Proof.
Since $`\mathrm{\Theta }_N`$ is ample, by a standard argument the assertion is true if the multiplication maps
$$H^0(𝒪(k\mathrm{\Theta }_N)H^0(𝒪(kl\mathrm{\Theta }_N))H^0(𝒪(k(l+1)\mathrm{\Theta }_N))$$
are surjective for $`l1`$. For $`l=1`$ this is proved in the theorem and the case $`l2`$ is similar but easier. ∎
###### Remark 5.11.
The case of line bundles in the theorem above (i.e. $`r=1`$) is the statement for Jacobians of a well known theorem of Koizumi (see and ). Applied to that particular case, the metod of proof in Step 2 above is of course implicit in Pareschi’s paper .
For an effective bound implied by the previous theorem we have to restrict ourselves to the case of rank $`2`$ vector bundles, since to the best of our knowledge nothing is known about multiplication maps on $`SU_X(r)`$ for $`r3`$.
###### Corollary 5.12.
For a generic curve $`X`$ the multiplication map
$$H^0(𝒪(k\mathrm{\Theta }_N))H^0(𝒪(k\mathrm{\Theta }_N))\stackrel{\mu _k}{}H^0(𝒪(2k\mathrm{\Theta }_N))$$
on $`U_X(2,0)`$ is surjective for $`k\mathrm{max}\{5,g2\}`$ and so $`𝒪(k\mathrm{\Theta }_N)`$ is very ample for such $`k`$.
###### Proof.
This follows by a theorem of Laszlo , which says that on a generic curve the multiplication map
$$S^kH^0(^2)H^0(^{2k})$$
on $`SU_X(2)`$ is surjective for $`kg2`$. See also §4 for a survey of results in this direction. ∎
###### Remark 5.13.
A more refined study along these lines gives analogous results in the extended setting of higher syzygies and $`N_p`$ properties. We hope to come back to this somewhere else.
We would like to end this section with another application to multiplication maps. Although probably not of the same significance as the previous results, it still brings some new insight through the use of methods characteristic to abelian varieties. Recall from 1.2 that the Picard group of $`U_X(r,0)`$ is generated by $`𝒪(\mathrm{\Theta }_N)`$ and the preimages of line bundles on $`J(X)`$. We want to study “mixed” multiplication maps of the form:
(8)
$$H^0(𝒪(k\mathrm{\Theta }_N)H^0(\mathrm{det}^{}𝒪_J(m\mathrm{\Theta }_N))\stackrel{𝛼}{}H^0(𝒪(k\mathrm{\Theta }_N)\mathrm{det}^{}𝒪_J(m\mathrm{\Theta }_N)).$$
###### Proposition 5.14.
The multiplication map $`\alpha `$ in (8) is surjective if $`m2`$ and $`k2r+1`$.
###### Proof.
By repeated use of the projection formula, from the commutative diagram
we see that it is enough to prove the surjectivity of the multiplication map $`\beta `$ on $`J(X)`$.
This is an application of the cohomological criterion 1.11 going back to Kempf . What we need to check is that
$$h^i(E_{r,k}𝒪_J(l\mathrm{\Theta }_N)\alpha )=0,i>0,l2\mathrm{and}\alpha \mathrm{Pic}^0(J(X)).$$
It is again enough to prove this after pulling back by multiplication by $`r`$. Now:
$$r_J^{}(E_{r,k}𝒪_J(l\mathrm{\Theta }_N)\alpha )\underset{s_{r,k}}{}𝒪_J((kr+lr^2)\mathrm{\Theta }_N)r_J^{}\alpha $$
hence the required vanishings are obvious as long as $`l2`$ and $`k2r+1`$. ∎
## 6. Variants for arbitrary degree
As it is natural to expect, some of the facts discussed in the previous sections for moduli spaces of vector bundles of degree $`0`$ can be extended to arbitrary degree. On the other hand, as the reader might have already observed, there are results that do not admit (at least straightforward) such extensions. In this last paragraph we would like to emphasize what can and what cannot be generalized using the present techniques.
Fix $`r`$ and $`d`$ arbitrary positive integers. Then we can look at the moduli space $`U_X(r,d)`$ of semistable vector bundles of rank $`r`$ and degree $`d`$. For $`A\mathrm{Pic}^d(X)`$, denote also by $`SU_X(r,A)`$ the moduli space of rank $`r`$ bundles with fixed determinant $`A`$. We will write $`SU_X(r,d)`$ when it is not important what determinant is involved. On these moduli spaces one can construct generalized theta divisors as in the degree $`0`$ case. More precisely, denote
$$h=\mathrm{gcd}(r,d),r_1=r/h\mathrm{and}d_1=d/h.$$
Then for any vector bundle $`F`$ of rank $`r_1`$ and degree $`d_1`$, we can consider $`\mathrm{\Theta }_F`$ to be the closure in $`U_X(r,d)`$ of the locus
$$\mathrm{\Theta }_F^s=\{E|h^0(EF)0\}U_X^s(r,d).$$
This does not always have to be a proper subset, but it is so for generic $`F`$ (see ) and in that case $`\mathrm{\Theta }_F`$ is a divisor. We can of course do the same thing with vector bundles of rank $`kr_1`$ and degree $`kd_1`$ for any $`k1`$ and a formula analogous to 1.3 holds. On $`SU_X(r,d)`$ there are similar divisors $`\mathrm{\Theta }_F`$ and they all determine the same determinant line bundle. As before this generates $`\mathrm{Pic}(SU_X(r,d))`$ and is denoted by $``$.
To study the linear series determined by these divisors we define Verlinde type bundles as in Section 2. They will now depend on two parameters (again not emphasized by the notation). Namely fix $`FU_X(r_1,r_1(g1)d_1)`$ and $`L\mathrm{Pic}^d(X)`$. Consider the composition
$$\pi _L:U_X(r,d)\stackrel{det}{}\mathrm{Pic}^d(X)\stackrel{L^1}{}J(X)$$
and define:
$$E_{r,d,k}(=E_{r,d,k}^{F,L}):=\pi _{L}^{}{}_{}{}^{}𝒪(k\mathrm{\Theta }_F).$$
This is a vector bundle on $`J(X)`$ of rank $`s_{r,d,k}:=h^0(SU_X(r,d),^k)`$. There is again a fiber diagram
where $`\tau `$ is given by tensor product and the top and bottom maps are Galois with Galois group $`X_r`$. By §3 one has the formula
$$\tau ^{}𝒪(\mathrm{\Theta }_F)𝒪_J(krr_1\mathrm{\Theta }_N),$$
where $`N\mathrm{Pic}^{g1}(X)`$ is a line bundle such that $`N^rL(\mathrm{det}F)^h`$. As in 2.3 we obtain the decomposition:
$$r_J^{}E_{r,d,k}\underset{s_{r,d,k}}{}𝒪_J(krr_1\mathrm{\Theta }_N).$$
The basic duality setup via Fourier-Mukai transform presented in Section 3 can be extended with a little extra care to this general setting. The purpose is to relate linear series on the complementary moduli spaces $`SU_X(r,d)`$ and $`U_X(kr_1,kr_1(g1)kd_1)`$ (cf. §5) and this is realized via a diagram of the form:
where $`L\mathrm{Pic}^d(X)`$, $`M\mathrm{Pic}^{kr_1(g1)kd_1}(X)`$ and $`\pi _L`$ and $`\pi _M`$ are defined as above. One can also choose a vector bundle $`GU_X(r_1,d_1)`$ and consider the Verlinde bundle $`E_{kr_1,kr_1(g1)kd_1,h}`$ associated to $`G`$ and $`M`$. As before, there is an obvious tensor product map
$$U_X(r,d)\times U_X(kr_1,kr_1(g1)kd_1)\stackrel{𝜏}{}U_X(krr_1,krr_1(g1)).$$
With the extra (harmless) assumption on our choices that $`L(\mathrm{det}G)^h`$ and $`M(\mathrm{det}F)^k`$ we can show exactly as in 3.9 that
$$\tau ^{}𝒪(\mathrm{\Theta })𝒪(k\mathrm{\Theta }_F)𝒪(h\mathrm{\Theta }_G)(\pi _L\times \pi _M)^{}𝒫,$$
where $`𝒫`$ is a normalized Poincaré line bundle on $`J(X)\times J(X)`$. Note that this is slightly different from 3.9 in the sense that we are twisting up to slope $`g1`$ and on $`U_X(krr_1,krr_1(g1))`$, $`\mathrm{\Theta }`$ represents the canonical theta divisor. The two formulations are of course equivalent.
The unique nonzero section of $`𝒪(\mathrm{\Theta })`$ induces then a nonzero map:
$$SD:E_{r,d,k}^{}(E_{kr_1,kr_1(g1)kd_1,h})^\widehat{}$$
and the global formulation of the full strange duality conjecture is:
###### Conjecture 6.1.
(cf. §5) $`SD`$ is an isomorphism.
On the positive side, the properties of the kernel of this map described in 3.13 still hold. On the other hand, the method of proof of Theorem 3.11 cannot be used for arbitrary degree even in the level $`1`$ situation. The point is that these new Verlinde type bundles may always fail to be simple. Probably the most suggestive example is the case of $`r`$ and $`d`$ coprime, when the other extreme is attained for any $`k`$:
###### Example 6.2.
If $`\mathrm{gcd}(r,d)=1`$, then $`E_{r,d,k}`$ decomposes as a direct sum of line bundles for any $`k`$. More precisely:
$$E_{r,d,k}=\underset{s_{r,d,k}}{}𝒪(k\mathrm{\Theta }_N),k1.$$
This can be seen by imitating the proof of 2.7.
On a more modest note, the main result of can be naturally integrated into these global arguments on $`J(X)`$. It is obtained by calculus with Fourier transforms in the spirit of Section 3 and we do not repeat the argument here:
###### Proposition 6.3.
(, Theorem 1) $`h^0(U_X(r,d),𝒪(k\mathrm{\Theta }_F))=\frac{k^g}{h^g}s_{r,d,k}`$.
Turning to effective global generation and normal generation $`U_X(r,d)`$, the picture described in Section 5 completely extends, with the appropriate modifications, to the general case. All the effective bounds turn out to depend on the number $`h=\mathrm{gcd}(r,d)`$. The global generation result analogous to 5.3 is formulated as follows:
###### Theorem 6.4.
$`𝒪(k\mathrm{\Theta }_F)`$ is globally generated on $`U_X(r,d)`$ as long as $`kh+1`$ and $`^k`$ is globally generated on $`SU_X(r,d)`$. Moreover, $`𝒪(k\mathrm{\Theta }_F)`$ is not globally generated for $`kh`$.
In §4 it is proved that $`^k`$ is globally generated on $`SU_X(r,d)`$ for $`k\mathrm{max}\{\frac{(r+1)^2}{4r}h,`$ $`\frac{r^2}{4s}h\}`$, where $`s`$ is an invariant of the moduli space that we will not define here, but satisfying $`sh`$ so that in particular $`k\frac{(r+1)^2}{4}`$ always works. This implies then:
###### Corollary 6.5.
$`𝒪(k\mathrm{\Theta }_F)`$ is globally generated on $`U_X(r,d)`$ for $`k\mathrm{max}\{\frac{(r+1)^2}{4r}h,\frac{r^2}{4s}h\}`$.
This again produces optimal results in the case of rank $`2`$ and rank $`3`$ vector bundles (see (5.7)):
###### Corollary 6.6.
$`𝒪(2\mathrm{\Theta }_F)`$ is globally generated on $`U_X(2,1)`$ and $`U_X(3,\pm 1)`$.
In the same vein, the normal generation result 5.9 can be generalized to:
###### Theorem 6.7.
The multiplication map
$$\mu _k:H^0(𝒪(k\mathrm{\Theta }_F))H^0(𝒪(k\mathrm{\Theta }_F))H^0(𝒪(2k\mathrm{\Theta }_F))$$
on $`U_X(r,d)`$ is surjective as long as the multiplication map $`H^0(^k)H^0(^k)H^0(^{2k})`$ on $`SU_X(r,d)`$ is surjective and $`k2h+1`$. For such $`k`$, $`𝒪(k\mathrm{\Theta }_F)`$ is very ample.
###### Corollary 6.8.
For $`X`$ generic $`\mu _k`$ is surjective on $`U_X(2,1)`$ if $`k\mathrm{max}\{3,\frac{g2}{2}\}`$.
###### Proof.
This is a consequence of 6.7 and , where it is proved that on $`SU_X(2,1)`$
$$S^kH^0()H^0(^k)$$
is surjective for $`kg2`$. ∎
It is certainly not hard to formulate further results corresponding to 5.14. We leave this to the interested reader. |
warning/0002/hep-th0002098.html | ar5iv | text | # 1 Introduction
## 1 Introduction
For a long time, it was believed that no regular particle-like stable solutions (solitons) with finite mass can exist in self gravitating systems unless the stability is guaranteed topologically. The Einstein theory in vacuum and the Einstein-Maxwell system do not admit solitons. It came as quite a surprise when Bartnik and McKinnon (BK) found globally regular solutions to the $`SU(2)`$ Einstein Yang-Mills (EYM) theory without scalar fields. It was unexpected to find that self gravitating Yang-Mills systems produced solitons. Unfortunately, the BK solutions were shown to be unstable against linear perturbations . Later, other fields such as Higgs scalar fields and dilaton fields were included in the EYM action, but with the exception of the Skyrmions, all turned out to be unstable (see for a review).
Interest in the BK solutions was renewed with the discovery of black hole solutions to the EYM equations . These non-Abelian black holes aparently violate the no-hair conjecture . But these non-Abelian black hole solutions are also unstable, and again other fields were added in the hope of achieving stability without success (see Ref.’s ) for a review).
We stress that it is a surprise that there are static solutions to the Einstein Yang-Mills equations at all. There are no static solutions to the Yang-Mills equations in 4 dimensional flat space. We can see this with a simple argument given by Deser . The conservation of the canonical energy momentum tensor, $`_\nu T^{\mu \nu }=0`$, implies that for a static field configuration $`_jT_{i}^{}{}_{}{}^{j}=0`$. The total divergence of the quantity $`x^iT_{i}^{}{}_{}{}^{j}`$ must vanish to maintain finite energy and regularity, $`d^{d1}x_j(x^iT_{i}^{}{}_{}{}^{j})=0`$. But $`_j(x^iT_{i}^{}{}_{}{}^{j})=T_{j}^{}{}_{}{}^{j}+x^i_jT_{i}^{}{}_{}{}^{j}=T_{j}^{}{}_{}{}^{j}`$ so that
$$d^{d1}xT_i^i=d^{d1}x\left[\frac{1}{2}(5d)F_{ij}^2+(d3)F_{0i}^2\right]=0.$$
(1)
Since the integrand above is positive definite for $`d=4`$, $`F_{ij}`$ and $`F_{0i}`$ must vanish. Thus there are no regular static solutions.
The argument above cannot be extended to curved spacetime. The conservation law $`T_{}^{\mu \nu }{}_{;\nu }{}^{}=0`$ leads to
$$d^{d1}x\sqrt{g}T_{j}^{}{}_{}{}^{j}=d^{d1}x\sqrt{g}x^k\mathrm{\Gamma }_{k\mu \nu }T^{\mu \nu }0.$$
(2)
The failure of Deser’s simple argument in curved space implies the possibility of having static solutions in curved space. Gravity supplies the attractive force needed to balance the repulsive force of Yang-Mills gauge interactions. Indeed, any solution to $`SU(2)`$ EYM equations in asymptotically Minkowski space which is regular asymptotically is also regular for all $`r>0`$ .
The particle-like and black hole solutions were later studied in a cosmological context. The behavior of static solutions to the Einstein Yang-Mills equations depends considerably on the sign of the cosmological constant. The solutions can be separated into two families; $`\mathrm{\Lambda }0`$ and $`\mathrm{\Lambda }<0`$. The solutions where $`\mathrm{\Lambda }=0`$ are the BK solutions. Their asymptotically de Sitter analogs ($`\mathrm{\Lambda }>0`$) were discovered independently by Volkov et. al. and Torii et. al. . The BK solutions and the cosmological extensions all share similar behavior, and are unstable . (See Ref. for a review). Recently, asymptotically anti-de Sitter black hole solutions and soliton solutions were found which are strikingly different from the BK type solutions. In particular, the asymptotically anti-de Sitter AdS EYM equations have solutions where the field strengths are non-zero everywhere. These solutions were also shown to be stable against spherically symmetric linear perturbations. These solutions are the only EYM solutions solutions that are stable. This discovery would be very important to cosmology if the universe was ever in a phase where the cosmological constant is negative.
Another new feature of the EYM theory in AdS is the existence of dyon solutions. If $`\mathrm{\Lambda }0`$ the electric part of the gauge fields is forbidden if the ADM mass is to remain finite. Scalar fields must be added to the theory in order for the boundary conditions at infinity to permit the electric fields and maintain a finite ADM mass .
Recently a tremendous amount of interest has evolved in field theories in AdS space. There is the AdS/CFT correspondence . Conformal field theories in $`d`$ dimensions ($`R_d`$) are described in terms of supergravity or string theory on the product space of AdS<sub>d+1</sub> and a compact manifold. There are intimate relations between data on the boundary $`R_d`$ of AdS<sub>d+1</sub> and data in the bulk AdS<sub>d+1</sub>. In the present paper we are examining the Einstein-Yang-Mills theory in asymptotically AdS space. The boundary in space must be playing a crucial role for the existence of stable monopole and dyon solutions, more detailed analysis of which is, however, left for future investigation. We also note that in the three-dimensional AdS space there exist nontrivial black holes and monopole/instanton solutions .
When the value of the cosmological constant $`\mathrm{\Lambda }`$ is varied, the space of monopole and dyon solutions, the moduli space, also changes. With a finite negative $`\mathrm{\Lambda }`$, solutions exist in continuum. They are classified in a finite number of families, or branches. With a vanishing or positive $`\mathrm{\Lambda }`$ solution exists only in a discrete set, but there are infinitely many. One natural question emerging is how these finite number of branches of solutions in continuum become infinitely many discrete points as $`\mathrm{\Lambda }<0`$ approaches 0. There is a surprising hidden feature in this limit. We shall find a fractal structure in the moduli space, which seems to explain the transition.
In the next section the general formalism is given and the equations of motion are derived with a spherically symmetric ansatz. Conserved charges in the Yang-Mills theory is defined in section 3. Some general no-go theorems are derived from sum rules in Section 4. New soliton solution in asymptotically anti-de Sitter space are explained in section 5. The critical spacetime which have universality near the edge of the space is also examined. Black hole solutions which have both magnetic and electric non-Abelian charges are presented in Section 6. The dependence of the moduli space on the cosmological constant $`\mathrm{\Lambda }`$ is investigated in Section 7 where the fractal structure is revealed when $`\mathrm{\Lambda }`$ approaches zero from the negative side. The detailed analysis of the stability of the monopole solutions is presented in Section 8. The subtle boundary condition in the problem requires elaboration of the previous argument presented in the $`\mathrm{\Lambda }=0`$ and $`\lambda >0`$ cases.
## 2 General Formalism
In non-Abelian gauge theory, the field equations have solutions which exhibit a magnetic charge. In the ’t Hooft-Polyakov monopole solution
$$\begin{array}{ccc}\mathrm{\Phi }_a=\frac{x_a}{er^2}H(er),\hfill & A_a^0=0,& \hfill A_a^i=ϵ_{aij}\frac{x_j}{er^2}(1K(er)).\end{array}$$
(3)
where $`\mathrm{\Phi }_a`$ is a triplet Higgs scalar field. Its stability is guaranteed by the topology of the triplet Higgs scalar field. The $`U(1)`$ magnetic charge takes a quantized value, $`4\pi /e`$. Dyon solutions were obtained starting with the above ansatz (3) but with a non-zero value for $`A_a^0`$, (i.e. $`A_a^0=(x_a/er^2)J(er)`$).
In this paper we look for monopole and dyon solutions in the Einstein-Yang-Mills theory without scalar fields;
$$S=d^4x\sqrt{g}\left[\frac{1}{16\pi G}(R2\mathrm{\Lambda })\frac{1}{4}F^{a\mu \nu }F_{}^{a}{}_{\mu \nu }{}^{}\right].$$
(4)
The Einstein and Yang-Mills equations are given by
$`R^{\mu \nu }{\displaystyle \frac{1}{2}}g^{\mu \nu }(R2\mathrm{\Lambda })=8\pi GT^{\mu \nu }`$ (5)
$`F_{}^{\mu \nu }{}_{;\mu }{}^{}+e[A_\mu ,F^{\mu \nu }]=0`$ (6)
We suspect that the gravity provides attractive force to balance the equation.
We look for spherically symmetric solutions. The metric takes the form
$$ds^2=\frac{H}{p^2}dt^2+\frac{dr^2}{H}+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)$$
(7)
whereas Yang-Mills fields are given , in the regular gauge, by
$`A^{(0)}`$ $`=`$ $`{\displaystyle \frac{\tau ^j}{2e}}\{A_0{\displaystyle \frac{x_j}{r}}dt+A_1{\displaystyle \frac{x_jx_k}{r^2}}dx_k`$ (9)
$`+{\displaystyle \frac{\varphi _1}{r}}(\delta _{jk}{\displaystyle \frac{x_jx_k}{r^2}})dx_kϵ_{jkl}{\displaystyle \frac{1\varphi _2}{r^2}}x_kdx_l\}`$
Here the Cartesian coordinate $`x^k`$’s are related to the polar coordinates $`(r,\theta ,\varphi )`$ as in the flat space. $`H`$, $`p`$, $`A_0`$, $`A_1`$, $`\varphi _1`$ and $`\varphi _2`$ are functions of $`r`$ for monopole or dyon solutions. In the discussion of the stability of the solutions they depend on both $`t`$ and $`r`$. The regularity of solutions at the origin demands that $`H`$, $`p`$ are finite, whereas $`A_0`$, $`A_1`$, and $`\varphi _10`$ and $`\varphi _21`$ at $`r=0`$.
### 2.1 Simplification of the static gauge field ansatz
Let $`A=A_\mu dx^\mu =\frac{1}{2}\tau ^aA_\mu ^adx^\mu `$, where $`\tau ^a`$ are the usual Pauli matrices. In terms of the basis in spherical coordinates $`(\tau _r,\tau _\theta ,\tau _\varphi )=(\stackrel{}{n}_r,\stackrel{}{n}_\theta ,\stackrel{}{n}_\varphi )\stackrel{}{\tau }`$ which satisfies $`[\tau _i,\tau _j]=2iϵ_{ijk}\tau _k`$ ($`i=r,\theta ,\varphi `$), the ansatz (9) is written as
$`A^{(0)}={\displaystyle \frac{1}{2e}}[A_0\tau _rdt+A_1\tau _rdr+(\varphi _1\tau _\theta +(\varphi _21)\tau _\varphi )d\theta `$ (10)
$`+((\varphi _21)\tau _\theta +\varphi _1\tau _\varphi )\mathrm{sin}\theta d\varphi ]`$ (11)
Note that there are no singularities in this gauge. Next make a gauge transformation $`A=SA^{(0)}S^1(i/e)dSS^1`$ where
$$S=\left(\begin{array}{cc}+e^{i(\varphi +\mathrm{\Omega })/2}\mathrm{cos}\frac{\theta }{2}& +e^{i(\varphi \mathrm{\Omega })/2}\mathrm{sin}\frac{\theta }{2}\\ e^{i(\varphi \mathrm{\Omega })/2}\mathrm{sin}\frac{\theta }{2}& +e^{i(\varphi +\mathrm{\Omega })/2}\mathrm{cos}\frac{\theta }{2}\end{array}\right),\mathrm{\Omega }=\mathrm{\Omega }(t,r).$$
(12)
Useful identities are
$`S\tau _rS^1`$ $`=`$ $`\tau _3`$ (13)
$`S\tau _\theta S^1`$ $`=`$ $`\mathrm{cos}\mathrm{\Omega }\tau _1\mathrm{sin}\mathrm{\Omega }\tau _2`$ (14)
$`S\tau _\varphi S^1`$ $`=`$ $`\mathrm{sin}\mathrm{\Omega }\tau _1+\mathrm{cos}\mathrm{\Omega }\tau _2`$ (15)
$`2idSS^1`$ $`=`$ $`(\mathrm{\Omega }^{}dr+\dot{\mathrm{\Omega }}dt)\tau _3d\theta (\mathrm{sin}\mathrm{\Omega }\tau _1+\mathrm{cos}\mathrm{\Omega }\tau _2)`$ (17)
$`+d\varphi (\mathrm{sin}\theta \mathrm{cos}\mathrm{\Omega }\tau _1\mathrm{sin}\theta \mathrm{sin}\mathrm{\Omega }\tau _2\mathrm{cos}\theta \tau _3)`$
The new gauge potential is
$$A=\frac{1}{2e}\left\{u\tau _3dt+\nu \tau _3dr+(w\tau _1+\stackrel{~}{w}\tau _2)d\theta +(\mathrm{cot}\theta \tau _3+w\tau _2\stackrel{~}{w}\tau _1)\mathrm{sin}\theta d\varphi \right\}.$$
(18)
where
$`u`$ $`=`$ $`A_0+\dot{\mathrm{\Omega }}`$ (19)
$`\nu `$ $`=`$ $`A_1+\mathrm{\Omega }^{}`$ (20)
$`w`$ $`=`$ $`+\varphi _1\mathrm{cos}\mathrm{\Omega }+\varphi _2\mathrm{sin}\mathrm{\Omega }`$ (21)
$`\stackrel{~}{w}`$ $`=`$ $`\varphi _1\mathrm{sin}\mathrm{\Omega }+\varphi _2\mathrm{cos}\mathrm{\Omega }.`$ (22)
Note that the gauge transformation (12) is singular at $`\theta =0`$ and $`\pi `$. Eq. (22) is the gauge potential in the singular gauge. It has a Dirac string. One can always choose $`\mathrm{\Omega }(t,r=0)=\pi /2`$ with which the boundary conditions at $`r=0`$ are $`u=\nu =\stackrel{~}{w}=0`$ and $`w=1`$. With appropriate $`\mathrm{\Omega }(t,r)`$ one can set $`\nu (t,r)=0`$ or $`u(t,r)=0`$.
A straightforward calculation leads to the field strength $`F=dAieAA`$:
$`F`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\{(\dot{\nu }u^{})\tau _3dtdr+[(\dot{w}u\stackrel{~}{w})\tau _1+(\dot{\stackrel{~}{w}}+uw)\tau _2]dtd\theta `$ (27)
$`\left[(uw+\dot{\stackrel{~}{w}})\tau _1+(u\stackrel{~}{w}\dot{w})\tau _2\right]dt\mathrm{sin}\theta d\varphi `$
$`+\left[(w^{}\nu \stackrel{~}{w})\tau _1+(\stackrel{~}{w}^{}+w\nu )\tau _2\right]drd\theta `$
$`+\left[(w^{}\nu \stackrel{~}{w})\tau _2+(\stackrel{~}{w}^{}\nu w)\tau _1\right]dr\mathrm{sin}\theta d\varphi `$
$`(1w^2\stackrel{~}{w}^2)\tau _3d\theta \mathrm{sin}\theta d\varphi \}.`$
The configurations where $`\nu =u=0`$, $`w=\stackrel{~}{w}=`$constant, and $`w^2+\stackrel{~}{w}^2=1`$ are pure gauge.
### 2.2 Equations of motion
In the general spherically symmetric metric (7) tetrads are
$$e_0=\frac{\sqrt{H}}{p}dt,e_1=\frac{1}{\sqrt{H}}dr,e_2=rd\theta ,e_3=r\mathrm{sin}\theta d\varphi .$$
(28)
In the tetrad basis $`F_{ab}=(e_a)_\mu (e_b)_\nu F^{\mu \nu }`$ and the energy-momentum tensors are $`T_{ab}=F_{ac}^{(i)}F_b^{c(i)}\frac{1}{4}\eta _{ab}F_{de}^{(i)}F^{de(i)}`$. The nonvanishing components of the Yang-Mills equations (6) are
$`\left(pr^2(u^{}\dot{\nu })\right)^{}2{\displaystyle \frac{p}{H}}\left\{w(uw+\dot{\stackrel{~}{w}})+\stackrel{~}{w}(u\stackrel{~}{w}\dot{w})\right\}=0`$ (29)
$`\left(pr^2(u^{}\dot{\nu })\right)_{,t}2{\displaystyle \frac{H}{p}}\left\{\stackrel{~}{w}w^{}+\stackrel{~}{w}^{}w+\nu (w^2+\stackrel{~}{w}^2)\right\}=0`$ (30)
$`\left({\displaystyle \frac{H}{p}}(w^{}\stackrel{~}{w}\nu )\right)^{}\left({\displaystyle \frac{p}{H}}(\dot{w}u\stackrel{~}{w})\right)_{,t}`$ (31)
$`+{\displaystyle \frac{p}{H}}u(uw+\dot{\stackrel{~}{w}})+{\displaystyle \frac{w(1w^2\stackrel{~}{w}^2)}{pr^2}}{\displaystyle \frac{H}{p}}\nu (\stackrel{~}{w}^{}+w\nu )=0`$ (32)
$`\left({\displaystyle \frac{H}{p}}(\stackrel{~}{w}^{}+w\nu )\right)^{}\left({\displaystyle \frac{p}{H}}(\dot{\stackrel{~}{w}}+uw)\right)_{,t}`$ (33)
$`+{\displaystyle \frac{p}{H}}u(u\stackrel{~}{w}\dot{w})+{\displaystyle \frac{\stackrel{~}{w}(1w^2\stackrel{~}{w}^2)}{pr^2}}+{\displaystyle \frac{H}{p}}\nu (w^{}\stackrel{~}{w}\nu )=0.`$ (34)
The nonvanishing components of the energy-momentum tensor are given by
$`T_{00}`$ $`=`$ $`{\displaystyle \frac{1}{e^2}}(A+B)`$ (35)
$`T_{11}`$ $`=`$ $`{\displaystyle \frac{1}{e^2}}(A+B)`$ (36)
$`T_{22}`$ $`=`$ $`T_{33}={\displaystyle \frac{1}{e^2}}A`$ (37)
$`T_{01}`$ $`=`$ $`{\displaystyle \frac{1}{e^2}}C`$ (38)
where
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2}}p^2(\dot{\nu }u^{})^2+{\displaystyle \frac{1}{2r^4}}(1w^2\stackrel{~}{w}^2)^2`$ (39)
$`B`$ $`=`$ $`{\displaystyle \frac{p^2}{r^2H}}\left\{(uw+\dot{\stackrel{~}{w}})^2+(\dot{w}u\stackrel{~}{w})^2\right\}+{\displaystyle \frac{H}{r^2}}\left\{(w^{}\nu \stackrel{~}{w})^2+(\stackrel{~}{w}^{}+\nu w)^2\right\}`$ (40)
$`C`$ $`=`$ $`{\displaystyle \frac{2p}{r^2}}\left\{(\stackrel{~}{w}^{}+\nu w)(uw+\dot{\stackrel{~}{w}})+(w^{}\nu \stackrel{~}{w})(\dot{w}u\stackrel{~}{w})\right\}.`$ (41)
The Einstein equations reduce to
$`{\displaystyle \frac{p^{}}{p}}={\displaystyle \frac{8\pi G}{e^2}}{\displaystyle \frac{rB}{H}}`$ (42)
$`{\displaystyle \frac{H^{}}{r}}+{\displaystyle \frac{1H}{r^2}}={\displaystyle \frac{8\pi G}{e^2}}(A+B)+\mathrm{\Lambda }`$ (43)
$`{\displaystyle \frac{p}{2}}\left\{\left({\displaystyle \frac{p\dot{H}}{H^2}}\right)_{,t}+\left({\displaystyle \frac{pH^{}2p^{}H}{p^2}}\right)^{}\right\}+{\displaystyle \frac{1H}{r^2}}={\displaystyle \frac{16\pi G}{e^2}}A`$ (44)
$`{\displaystyle \frac{p\dot{H}}{rH}}={\displaystyle \frac{8\pi G}{e^2}}C`$ (45)
It is convenient to introduce $`m(r)`$ defined by
$$H(r)=1\frac{2m(r)}{r}\frac{\mathrm{\Lambda }r^2}{3}.$$
(46)
$`m(r)`$ is the mass contained inside the radius $`r`$. $`p(r)=`$constant and $`m(r)=0`$ corresponds to the Minkowski, de Sitter, or anti-de Sitter space for $`\mathrm{\Lambda }=0,>0,`$ or $`<0`$, respectively. Then the second equation in (45) becomes
$$m^{}=\frac{4\pi G}{e^2}r^2(A+B).$$
(47)
The system of the Einstein-Yang-Mills equations contains one redundant equation. The third equation in (45) follows from (34) and the rest of (45).
### 2.3 Static configurations
It is most convenient to take the $`\nu =0`$ gauge for static configurations. The second equation in (34) then yields $`w\stackrel{~}{w}^{}w^{}\stackrel{~}{w}=0`$, which leads to $`\stackrel{~}{w}(r)=Cw(r)`$. By a further global rotation $`\mathrm{\Omega }=`$constant in (22) one can set $`\stackrel{~}{w}=0`$. As a result
$`A`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\left\{u\tau _3dt+w\tau _1d\theta +(\mathrm{cot}\theta \tau _3+w\tau _2)\mathrm{sin}\theta d\varphi \right\}`$ (48)
$`F`$ $`=`$ $`{\displaystyle \frac{1}{2e}}\{u^{}\tau _3dtdr+uwdt(\tau _2d\theta \tau _1\mathrm{sin}\theta d\varphi )`$ (50)
$`+w^{}dr(\tau _1d\theta +\tau _2\mathrm{sin}\theta d\varphi )(1w^2)\tau _3d\theta \mathrm{sin}\theta d\varphi \}.`$
Then the Einstein-Yang-Mills equations are
$`\left({\displaystyle \frac{H}{p}}w^{}\right)^{}`$ $`=`$ $`{\displaystyle \frac{p}{H}}u^2w{\displaystyle \frac{w}{p}}{\displaystyle \frac{(1w^2)}{r^2}}`$ (51)
$`\left(r^2pu^{}\right)^{}`$ $`=`$ $`{\displaystyle \frac{2p}{H}}w^2u`$ (52)
$`p^{}`$ $`=`$ $`{\displaystyle \frac{2v}{r}}p\left[(w^{})^2+{\displaystyle \frac{u^2w^2p^2}{H^2}}\right]`$ (53)
$`m^{}`$ $`=`$ $`v\left[{\displaystyle \frac{(w^21)^2}{2r^2}}+{\displaystyle \frac{1}{2}}r^2p^2(u^{})^2+H(w^{})^2+{\displaystyle \frac{u^2w^2p^2}{H}}\right]`$ (54)
where $`v=4\pi G/e^2`$.
These equations are solved with the given boundary conditions. Near the origin solutions must be regular so that
$`u(r)`$ $`=`$ $`ar+{\displaystyle \frac{a}{5}}\left\{2b+{\displaystyle \frac{1}{3}}\mathrm{\Lambda }+2v(a^2+4b^2)\right\}r^3`$ (55)
$`w(r)`$ $`=`$ $`1br^2`$ (56)
$`m(r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}v(a^2+4b^2)r^3`$ (57)
$`p(r)`$ $`=`$ $`1v(a^2+4b^2)r^2`$ (58)
where $`a`$ and $`b`$ are arbitrary constants. The boundary conditions at the origin of the EYM equations are completely determined by the values of the constants $`a`$ and $`b`$.
At space infinity the energy-momentum tensors $`T_{ab}`$ in (38) must approach zero sufficiently fast. Further we expect that the metric must asymptotically (anti-) de Sitter space, depending on the value of $`\mathrm{\Lambda }`$. This, with the equations of motion, leads to the asymptotic expansion at large $`r`$;
$`u=u_0+u_1{\displaystyle \frac{1}{r}}+\mathrm{},w=w_0+w_1{\displaystyle \frac{1}{r}}+\mathrm{}`$
$`m=M+m_1{\displaystyle \frac{1}{r}}+\mathrm{},p=p_0+p_4{\displaystyle \frac{1}{r^4}}+\mathrm{}`$ (59)
where $`u_0`$, $`u_1`$, $`w_0`$, $`w_1`$, $`m_1`$, $`p_0`$ and $`p_4`$ are constants to be determined and $`M`$ is the ADM mass, $`M=m(\mathrm{})m(0)`$.
## 3 Conserved charges
Solutions to eq.’s (51) to (54) are classified by the ADM mass, $`M=m(\mathrm{})m(0)`$, electric and magnetic charges, $`Q_E`$ and $`Q_M`$. From the Gauss flux theorem
$$\left(\begin{array}{c}Q_E\\ Q_M\end{array}\right)=\frac{e}{4\pi }𝑑S_k\sqrt{g}\left(\begin{array}{c}F^{k0}\\ \stackrel{~}{F}^{k0}\end{array}\right)$$
(60)
are conserved. With the ansatz in the singular gauge (18) and the asympotitic behavior (59), the charges are given by
$$\left(\begin{array}{c}Q_E\\ Q_M\end{array}\right)=\left(\begin{array}{c}u_1p_0\\ 1w_0^2\end{array}\right)\frac{\tau _3}{2}$$
(61)
Notice that the electric charge $`Q_E`$ is determined by $`u_1`$, whereas the magnetic charge $`Q_M`$ by $`w_0`$. If $`(u,w,m,p)`$ is a solution, then $`(u,w,m,p)`$ is also a solution. Dyon solutions come in a pair with $`(\pm Q_E,Q_M,M)`$.
The charges (60) are not gauge invariant, however. Under a local gauge transformation $`AUAU^1(i/e)dUU^1`$, $`Q_E`$ and $`Q_M`$ are transformed to
$$\left(\begin{array}{c}Q_{E}^{}{}_{}{}^{U}\\ Q_{M}^{}{}_{}{}^{U}\end{array}\right)=\frac{e}{4\pi }𝑑S_k\sqrt{g}U(x)\left(\begin{array}{c}F^{k0}\\ \stackrel{~}{F}^{k0}\end{array}\right)U^1(x)$$
(62)
In non-Abelian gauge theory a set of charges $`\{Q_{E}^{}{}_{}{}^{U},Q_{M}^{}{}_{}{}^{U}\}`$ are conserved. In the rest of the paper we use the charges, (61), defined in the singular gauge.
The effective charge $`Q_{\mathrm{eff}}`$ is defined by the asymptotic behavior of $`H(r)`$;
$$H(r)=1\frac{2M}{r}+\frac{Q_{\mathrm{eff}}^2}{r^2}\frac{1}{3}\mathrm{\Lambda }r^2.$$
(63)
In terms of the coefficients in (59), $`Q_{\mathrm{eff}}^2=2m_1`$. This requires that $`m_1<0`$ which indeed is the case. After inserting eq. (59) into eq.’s (51) and (54) we find the relation
$$Q_{\mathrm{eff}}^2=2v\mathrm{Tr}(Q_E^2+Q_M^2)\frac{4\mathrm{\Lambda }}{3}\frac{p_4}{p_0}.$$
(64)
Eq. (53) implies that $`p(r)`$ is a monotonically decreasing positive function so that $`p_0>0`$ and $`p_4>0`$. The effective charge is smaller (larger) than $`2v\mathrm{Tr}(Q_E^2+Q_M^2)`$ for $`\mathrm{\Lambda }>0`$ ($`<0`$). The relation (64) incidentally implies that the charges defined in the singular gauge have physical, gauge invariant meaning.
## 4 Sum rules
Sum rules are obtained from the equations of motion. First, multiply both sides of (52) by $`u`$ and integrate in part.
$$pr^2uu^{}|_{r_1}^{r_2}=_{r_1}^{r_2}𝑑r\left\{r^2p(u^{})^2+2\frac{p}{H}u^2w^2\right\}.$$
(65)
Secondly, multiply both sides of (51) by $`w`$ and integrate in part:
$$\frac{H}{p}ww^{}|_{r_1}^{r_2}=_{r_1}^{r_2}𝑑r\left\{\frac{H}{p}(w^{})^2\frac{p}{H}u^2w^2\frac{1}{pr^2}w^2(1w^2)\right\}.$$
(66)
Thirdly, divide both sides of (51) by $`w`$ and integrate in part:
$$\frac{Hw^{}}{pw}|_{r_1}^{r_2}=_{r_1}^{r_2}𝑑r\left\{\frac{p}{H}u^2+\frac{H}{p}\left(\frac{w^{}}{w}\right)^2+\frac{1}{pr^2}(1w^2)\right\}.$$
(67)
These relations are valid, provided the integrals on the right hand sides are defined. Several important conclusions follow from (65) - (67).
### 4.1 In asymptotically flat space
Consider (65) with $`r_1=0`$ and $`r_2=\mathrm{}`$. For regular solutions $`u(0)=0`$. Both $`p`$ and $`H`$ approach constant as $`r\mathrm{}`$. The finiteness of the ADM mass requires that $`uw|_{r=\mathrm{}}=0`$. In the expansion (59), $`u_0w_0=0`$. On the other hand, if $`u_00`$, $`w_0=0`$ and Eq. (52) implies $`w_10`$ so that Eq. (51) cannot be satisfied. Hence $`u(\mathrm{})=0`$. Then the left hand side of (65) vanishes, implying that $`u(r)`$ must vanish identically. There is no regular electrically charged solution. Furthermore, (51) can be solved only if $`(w_0)^2=1`$ as $`H(\mathrm{})=1`$, therefore the magnetic charge $`Q_M`$ vanishes.
Suppose that $`w(r)`$ never vanishes and $`w^21`$ for $`0<r<\mathrm{}`$. Consider (67) with $`r_1=0`$ and $`r_2=\mathrm{}`$. The l.h.s. vanishes, but the integrand on the r.h.s. is positive definite except for the pure gauge configuration $`w(r)=\pm 1`$. This implies that non-trivial solutions with $`w^21`$ must vanish at least once.
We also note that the singular solution $`w(r)=0`$, $`u^{}(r)=r^2`$, and $`p(r)=1`$ is nothing but the Reissner-Nordström solution.
### 4.2 In asymptotically de Sitter space
In asymptotically de Sitter space $`H(r)\mathrm{\Lambda }r^2/3`$ as $`r\mathrm{}`$. In this case the finiteness of the ADM mass does not forbid non-vanishing $`uw`$ at $`r=\mathrm{}`$. However, there arises a cosmological horizon at $`r=r_h`$ where $`H(r_h)=0`$.
It follows from (53) and (54) that $`u`$ or $`w`$ must vanish at $`r=r_h`$. Now consider (67). Suppose that $`w(r_h)0`$, or equivalently $`w0`$ for $`r_1rr_2=r_h`$ for some $`r_1>0`$, the left hand side is finite so that $`u(r_h)=0`$. However, Eq. (52) implies that $`u(r)`$ is a monotonically increasing or decreasing function. With the boundary condition $`u(0)=0`$, the only possibility available is $`u(r)=0`$. This conclusion remains valid even if $`w(r_h)=0`$. In this case the first and second terms on the r.h.s. give positively divergent contributions near $`r_h`$, whereas the l.h.s. remains finite. To summarize, there is no solution with nonvanishing $`u(r)`$.
If $`w>0`$ for $`0rr_h`$, then the left hand side vanishes in the $`r_10`$ and $`r_2r_h`$ limit. If one further assumes that $`w^21`$ in the interval, the integrand on the right hand side is positive definite so that the only solution is $`w(r)=\pm 1`$, which is a pure gauge.
This argument also shows that a nontrivial solution $`w(r)`$, which satisfies $`w^21`$ for $`0rr_h`$, must vanish at least once in this interval. We have numerically looked for solutions in which $`w^21`$ for small $`r`$, but have found that no such solution exists.
### 4.3 In asymptotically anti-de Sitter space
In asymptotically anti-de Sitter space there exist solutions in which $`H(r)>0`$ everywhere. As $`H(r)|\mathrm{\Lambda }|r^2/3`$ for large $`r`$, the condition for the finiteness of the ADM mass requires only that $`uw`$ constant. In other words, both $`u`$ and $`w`$ may approach nonvanishing values as $`r\mathrm{}`$. Consequently, with the expansion (59) the l.h.s. of (65) with $`r_1=0`$ and $`r_2=\mathrm{}`$ is $`p_0u_0u_1`$ and can be non-vanishing. Solutions with $`u(r)0`$ are allowed.
In critical cases a cosmological horizon appears and $`H(r_h)=0`$. In this case the argument above for the asymptotically de Sitter case applies and either $`u`$ or $`w`$ must vanish at $`r=r_h`$. We shall find, indeed, $`w(r_h)=0`$ below.
## 5 Soliton solutions
### 5.1 The BK solution
Particle like solutions of the EYM equations in asymptotically Minkowski space were first found by Bartnik and McKinnon (BK) in 1988. When $`u=0`$, we can solve the equations (51) to (54) numerically with the boundary condition $`a=0`$.
As already discussed, $`w(r)=1`$ and $`H(r)=1`$ corresponds to a pure gauge configuration. Similarly, if $`w=0`$ as $`r\mathrm{}`$ we are left with the RN solution which is singular at $`r=0`$. The Yang-Mills charge (61) vanishes in asymptotically Minkowski space since $`w_0=\pm 1`$. The effective charge $`Q_{\mathrm{eff}}=0`$ in all BK solutions, since $`Q_M=\mathrm{\Lambda }=0`$. All solutions possess a metric that is asymptotically Schwarzschild.
There is a discrete set of BK solutions, labeled by the number of nodes in $`w`$, $`n[1,\mathrm{})`$, and the free shooting parameter $`b`$. Since all BK solutions have at least one node, $`n1`$, the solutions are unstable against spherically symmetric perturbations .
### 5.2 Solutions with a cosmological horizon
Solutions in asymptotically de Sitter space were obtained by adding a cosmological constant ($`\mathrm{\Lambda }`$) term to the Einstein equations. Solutions display the same basic properties as the BK solutions. Just as in the BK solution, these equations are solved numerically, using the shooting method and requiring $`a=0`$.
In asymptotically de Sitter space, a cosmological horizon, where $`H=0`$, develops at $`r=r_h`$. At the horizon $`H^{}(r_h)0`$. Near the horizon
$`w(r)`$ $`=`$ $`w_0+w_1x+w_2x^2+\mathrm{}`$ (68)
$`p(r)`$ $`=`$ $`p_0+p_1x+p_2x^2+\mathrm{}`$ (69)
$`m(r)`$ $`=`$ $`m_0+m_1x+m_2x^2+\mathrm{}`$ (70)
$`H(r)`$ $`=`$ $`h_1x+h_2x^2+\mathrm{}`$ (71)
where $`x=rr_h`$. $`m_0`$ and $`r_h`$ are related by $`1(2m_0/r_h)(\mathrm{\Lambda }r_h^2/3)=0`$. With given $`w_0`$, $`r_h`$, and $`p_0`$, the equations (51), (53), and (54) (with $`u=0`$) determine all other coefficients, provided $`w_00,\pm 1`$ and $`p_00`$;
$`m_1`$ $`=`$ $`{\displaystyle \frac{v}{2r_h^2}}(w_0^21)^2`$ (72)
$`h_1`$ $`=`$ $`{\displaystyle \frac{2m_1}{r_h}}+{\displaystyle \frac{2m_0}{r_h^2}}{\displaystyle \frac{2}{3}}\mathrm{\Lambda }r_h`$ (73)
$`w_1`$ $`=`$ $`{\displaystyle \frac{w_0(1w_0^2)}{h_1r_h^2}}`$ (74)
$`p_1`$ $`=`$ $`{\displaystyle \frac{2vw_1^2p_0}{r_h}}`$ (75)
and so on. This expansion is valid independent of the value of $`\mathrm{\Lambda }`$. The critical case $`h_1=H^{}(r_h)=0`$ requires a special treatment, and will be analyzed in Section (5.3). Künzle and Masood-ul-Alam have argued that $`w(r)`$ has $`\sqrt{|rr_h|}`$ singularity at $`r=r_h`$. However, we have found that the regular expansion (71) is valid.
Just like the BK case, there are a discrete set of solutions labeled by the number of nodes in $`w(r)`$, $`n`$ and the parameter $`b`$. Solutions in $`w`$ and $`m`$, have the same form as the BK solutions except that $`w(r)`$ no longer approaches 1 at infinity. Eq. (64) implies that there is a non-vanishing charge $`Q_M`$ (or $`w1`$) if $`\mathrm{\Lambda }0`$ and that the solutions are all asymptotically Reissner-Nordström type. $`|w(\mathrm{})|`$ is slightly greater than 1 for the $`n=1`$ solutions. When $`n2`$, $`1>w(\mathrm{})>0`$ where $`w(\mathrm{})0`$ as $`n\mathrm{}`$.
The mass $`m(r)`$ also stays finite. For all indices $`n`$, it is small near the origin and does not grow until it approaches the horizon where it quickly climbs to a value near 1. After the horizon it stays almost constant.
The position of the horizon depends on $`\mathrm{\Lambda }`$. As long $`\mathrm{\Lambda }`$ is below some critical value, the geometry approaches Reissner-Nordström-de Sitter space in the asymptotic region as indicated by Eq. (64). Above some critical value the topology changes and a singularity appears. Since the position of the horizon depends on $`\mathrm{\Lambda }`$, there is a value for $`\mathrm{\Lambda }`$ where this singularity and the horizon meet, in which the topology becomes a completely regular manifold. Above this value for $`\mathrm{\Lambda }`$, the solutions are no longer regular. More details of the topology dependence on $`\mathrm{\Lambda }`$ can be found in . All the solutions are unstable .
### 5.3 Solutions in asymptotically anti-de Sitter space
As already discussed, there are no boundary conditions that forbid a solution to the EYM equations in asymptotically anti-de Sitter space with a non-zero electric component, $`u(r)`$, to the Yang-Mills fields. Solutions to Eqs. (51) to (54) are determined with the cosmological constant $`\mathrm{\Lambda }`$ fixed at some negative value.
(a) Monopole solutions
Monopole solutions are obtained by setting $`a=0`$ ($`u=0`$). By varying the initial condition parameter $`b`$, a continuum of monopole solutions are found which are regular in the entire space. Just as in the BK and dS solutions, $`w`$ crosses the axis an arbitrary number of times depending on the value of the adjustable shooting parameter $`b`$. In contrast to the $`\mathrm{\Lambda }=0`$ and $`\mathrm{\Lambda }>0`$ cases which have a discrete set of solutions in $`b`$ and $`n`$, there is a continuum of solutions in $`b`$ for each $`n`$. Typical solutions are displayed in fig. 1.
The behavior of $`m`$ and $`p`$ is similar to that of the asymptotically dS solutions . In contrast, as shown in fig. 1, there exist solutions where $`w`$ has no nodes. These solutions are of particular interest because they are shown to be stable against linear perturbations.
(b) Dyon solutions
Dyon solutions to the EYM equations are determined if the adjustable shooting parameter $`a`$ is chosen to be non-zero for a given negative $`\mathrm{\Lambda }`$. Just as in the monopole solutions, we find a continuum of solutions where $`w`$ crosses the axis an arbitrary number of times depending on $`a`$ and $`b`$. Also similar to the monopole solutions is the existence of solutions where $`w`$ does not cross the axis. As shown in Fig. 2, the electric component, $`u`$, of the EYM equations starts at zero and monotonically increases to some finite value. The behavior of $`w`$, $`m`$ , $`H`$, and $`p`$ is similar to that in the monopole solutions.
Just as for the monopole case, dyon solutions are found for a continuous set of parameters, $`a`$ and $`b`$. For some values of $`a`$ and $`b`$, solutions blow up, or the function $`H(r)`$ crosses the axis and becomes negative.
(c) Critical solutions
As the parameter $`b`$ is increased, the minimum of $`H(r)`$ hits zero from above, i.e. $`H(r_h)=H^{}(r_h)=0`$. This constitutes a special case and needs careful examination. Numerical studies indicate that this happens in a finite range of the parameter $`a`$. The critical solution exists for both $`u(r)=0`$ and $`u(r)0`$ cases. One example of solutions near the critical value \[$`(a,b)=(0.01,0.69)`$\] is displayed in fig. 3. $`H(r)`$ becomes very close to zero at $`r1`$. It has $`(Q_E,Q_M,M)(0.015,0.998,0.995)`$.
When $`b=b_c`$, $`w`$ and $`p`$ vanish at $`r=r_h`$ as well. As $`p(r_h)=0`$, $`p(r)=0`$ for $`rr_h`$. The space ends at $`r=r_h`$. There is universality at the critical point.
The numerical integration of the differential equations indicates that $`m(r)`$ and $`u(r)`$ are regular at $`r=r_h`$. The appropriate ansatz for the critical solutions with $`H(r_h)=H^{}(r_h)=0`$ is, for $`y=r_hr0`$,
$`u`$ $`=`$ $`u_0+u_1y+u_2y^2+\mathrm{}`$ (76)
$`w`$ $`=`$ $`y^\alpha \left\{w_0+w_1y+w_2y^2+\mathrm{}\right\}`$ (77)
$`p`$ $`=`$ $`y^\beta \left\{p_0+p_1y+p_2y^2+\mathrm{}\right\}`$ (78)
$`H`$ $`=`$ $`y^\gamma \left\{h_0+h_1y+h_2y^2+\mathrm{}\right\}`$ (79)
$`m`$ $`=`$ $`m_0+m_1y+m_2y^2+\mathrm{}`$ (80)
Eq. (51) implies that $`\gamma =2`$. When $`u(r)0`$, Eq. (52) leads to $`\alpha =\frac{1}{2}`$. When $`u(r)=0`$, Eq. (53), instead, implies that $`\alpha =\frac{1}{2}`$. Other relations obtained from eqs. (51) - (54) are
$`\left(\beta {\displaystyle \frac{3}{2}}\right)h_0`$ $`=`$ $`{\displaystyle \frac{2}{r_h^2}}`$ (81)
$`\beta r_h^2u_1`$ $`=`$ $`{\displaystyle \frac{2w_0^2u_0}{h_0}}`$ (82)
$`\beta `$ $`=`$ $`{\displaystyle \frac{vw_0^2}{2r_h}}`$ (83)
$`m_1`$ $`=`$ $`{\displaystyle \frac{v}{2r_h^2}}.`$ (84)
The value of the index $`\beta `$ is unconstrained when $`u=0`$. However, if $`u0`$, the consistency of eq. (51), for instance, demands that $`2\beta `$ be an integer. The smallest value for $`\beta `$ which satisfies the first relation in (84) is $`\beta =2`$, as $`h_0>0`$. We have confirmed this by numerical studies. The relation (46) further implies that
$`{\displaystyle \frac{2m_0}{r_h}}`$ $`=`$ $`1{\displaystyle \frac{\mathrm{\Lambda }}{3}}r_h^2`$ (85)
$`m_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Lambda }r_h^21).`$ (86)
From the two relations for $`m_1`$, one in (84) and the other in (86), $`r_h`$ is determined as a function of $`v`$ and $`\mathrm{\Lambda }<0`$:
$$r_h^2=\frac{1}{2|\mathrm{\Lambda }|}\left(\sqrt{1+4v|\mathrm{\Lambda }|}1\right).$$
(87)
To summarize, the indices in (80) are given by
$$\alpha =\frac{1}{2},\beta =\gamma =2.$$
(88)
The coefficients $`m_0`$, $`m_1`$, and
$$w_0^2=\frac{4r_h}{v},h_0=\frac{4}{r_h^2}$$
(89)
are all determined by $`v`$ and $`\mathrm{\Lambda }`$ only. We are observing the universality in the behavior of the critical solutions. The coefficients $`u_0`$ and $`p_0`$ depend on $`a`$ or $`b`$ as well.
For small $`v|\mathrm{\Lambda }|1`$
$`r_h\sqrt{v},m_0{\displaystyle \frac{1}{2}}\sqrt{v},m_1{\displaystyle \frac{1}{2}}`$ (90)
$`w_0{\displaystyle \frac{2}{v^{1/4}}},h_0{\displaystyle \frac{4}{v}}`$ (91)
They are all determined by $`v`$ only. This universal behavior is clearly seen in the solution in fig. 3 which is very close to the critical one.
The meaning of the critical spacetime is yet to be clarified. The space ends at $`r=r_h`$. It defines a spacetime with a boundary.
(d) Spectrum of monopole and dyon solutions
Monopole and dyon solutions permit non-vanishing charges $`Q_M`$ and $`Q_E`$, although there are solutions where $`Q_M=0`$ and $`Q_E0`$ or where $`Q_E=0`$ but $`Q_M0`$. Non-zero charges $`Q_M`$ or $`Q_E`$ ensures that $`Q_{\mathrm{eff}}0`$ (see Eq. (64) ) so that solutions are asymptotically of the AdS Reissner-Nordström type.
In fig. 4 the mass $`M`$ is plotted as a function of $`Q_M`$ for monopole solutions at $`\mathrm{\Lambda }=0.01`$ and $`v=1`$. The behavior of the solutions near t$`b=b_c=0.7104`$ needs more careful analysis.
Dyon solutions are found in a good portion of the $`Q_E`$-$`Q_M`$ plane. There are solutions with $`Q_M=0`$ but $`Q_E0`$. Although $`Q_M=0`$, i.e. $`w(\mathrm{})=\pm 1`$, $`w(r)0`$. In the shooting parameter space $`(a,b)`$, these solutions correspond not exactly, but almost to a universal value for $`b0.0054`$. See fig. 5. More surprising is the fact that $`Q_M`$ takes a quantized value $`(4\pi )^{1/2}`$ at $`b=0.0061`$ independent of the value of $`a`$ within numerical errors. We have not understood why it should be so.
Solutions with no node in $`w(r)`$ have special importance, as they are stable against small fluctuations. (See section 8.) In fig. 6 the spectrum of nodeless dyon solutions are presented in the parameter space $`(a,b)`$. Notice that $`a`$ must be small enough ($`a<0.005`$) even for $`b<0`$.
(e) Dependence of the coupling $`v`$ on the solutions
The ADM mass $`M`$ depends on the value of the coupling $`v=4\pi G/e^2`$. As shown in fig. 7, $`M`$ increases as $`v`$ gets larger and decreases when $`v`$ gets smaller. With fixed $`(a,b)`$, roughly $`Mv`$. The Yang-Mills fields $`u`$ and $`w`$ are roughly independent of $`v`$.
## 6 Black hole solutions
Not long after the BK solutions were discovered, black hole solutions were also found to be contained in the EYM equations if different boundary conditions were used. These solutions generated a large amount of further study, as they apparently violate the no-hair conjecture . Later, EYM black holes were studied in a cosmological context by including a positive cosmological constant. These black hole solutions share most of the properties as the soliton solutions, including their instability.
Recently, purely magnetic black hole solutions were found in asymptotically anti-de Sitter space . These solutions are drastically different from their asymptotically Minkowski or de Sitter counterparts. There are a continuum of solutions in terms of the adjustable shooting parameter that specifies the initial conditions at the horizon. Furthermore, there exist solutions that have no node in $`w`$ and are in turn stable against spherically symmetric linear perturbations. Here, we discuss the solutions found by Winstanley and also present new dyon black hole solutions. We also discuss the apparent shrinking of the moduli space when the magnitude of $`\mathrm{\Lambda }`$ is decreased. Similar to the particle-like solutions already discussed, the moduli space becomes discrete in the $`\mathrm{\Lambda }0`$ limit.
### 6.1 Boundary conditions at the horizon
Black hole solutions are obtained numerically by specifying the boundary conditions at the horizon and shooting for regular solutions $`w`$, $`u`$, $`m`$ and $`p`$ for $`r_hr<\mathrm{}`$. The location of the horizon, $`r_h`$, and the value $`p(r_h)>0`$ can be arbitrarily chosen by scaling of $`t`$ and $`r`$. We look for solutions in which $`H(r)>0`$ for $`r>r_h`$. As $`H(r_h)=0`$ but $`p(r_h)0`$, Eqs. (51) - (54) require that either $`u`$ or $`w`$ vanishes at the horizon. A stronger condition is obtained from the sum rule (67) with $`r_1=r_h`$ and $`r_2=\mathrm{}`$. Its l.h.s. is finite so that $`u(r_h)=0`$ on its r.h.s. Hence we are led to the expansion
$`w`$ $`=`$ $`w_0+w_1x+\mathrm{}`$ (92)
$`u`$ $`=`$ $`u_1x+\mathrm{}`$ (93)
$`p`$ $`=`$ $`1+p_1x+\mathrm{}`$ (94)
$`H`$ $`=`$ $`h_1x+\mathrm{}`$ (95)
$`m`$ $`=`$ $`m_0+m_1x+\mathrm{}`$ (96)
where $`x=rr_h`$. We have chosen $`p(r_h)=1`$ without loss of generality.
There are two adjustable shooting parameters, $`(a,b)=(u_1,w_0)`$. After inserting the ansatz into eq.’s (46) to (54) we find
$`m_0`$ $`=`$ $`{\displaystyle \frac{r_h}{2}}{\displaystyle \frac{\mathrm{\Lambda }r_h^3}{6}}`$ (97)
$`m_1`$ $`=`$ $`{\displaystyle \frac{v}{2}}\left\{{\displaystyle \frac{(1w_0^2)^2}{r_h^2}}+u_1^2\right\}`$ (98)
$`h_1`$ $`=`$ $`{\displaystyle \frac{1}{r_h}}(1\mathrm{\Lambda }r_h^22m_1)`$ (99)
$`w_1`$ $`=`$ $`{\displaystyle \frac{w_0(1w_0^2)}{r_h^2h_1}}`$ (100)
$`p_1`$ $`=`$ $`{\displaystyle \frac{2v}{r_h}}\left\{(w_1)^2+{\displaystyle \frac{w_0^2u_1^2}{h_1^2}}\right\}`$ (101)
$`u_2`$ $`=`$ $`\left\{{\displaystyle \frac{w_0^2}{r_h^2h_1}}+{\displaystyle \frac{1}{r_h}}+{\displaystyle \frac{p_1}{2}}\right\}.`$ (102)
The asymptotic expansion at large $`r`$ is the same as in (59). The ADM mass is given by $`M=m(\mathrm{})`$.
### 6.2 New electrically and magnetically charged black hole solutions
Just as the soliton solutions, purely magnetically charged black hole solutions are obtained by setting the adjustable parameter $`a`$ to zero. The behavior of the solutions are similar to that of the solitons (see Ref. for more information). The number of nodes $`n`$ in $`w`$ can be 0, 1, 2, $`\mathrm{}`$. The black hole monopole spectrum of mass versus charge is displayed in fig. 8. It shows the spectrum for the $`n=0`$ and $`n=1`$ arms.
Solutions with both magnetic and electric charge are obtained by giving $`a`$ a finite value. Dyon black hole solutions are similar to the monopole solutions except that $`u`$ is nonzero. At the horizon $`u`$ starts at zero and monotonically increases asymptotically to a finite value. $`H`$ starts at one and quickly diverges. $`p`$ starts at one and remains almost constant. Typical black hole dyon solutions are shown in Fig. 9.
Again black hole dyon solutions with no node in $`w(r)`$ are stable against small spherically symmetric perturbations. The spectrum of those nodeless black hole dyons in the parameter space $`(a,b)`$ is plotted in fig. 10. Notice the similarity between fig. 6 and fig. 10. The nodeless solutions exist only for small $`a=u_1<0.0055`$. $`b=w_0`$ must be around 1.
## 7 Dependence on $`\mathrm{\Lambda }`$ – fractal structure
The soliton and black hole solutions depend non-trivially on the value of the cosmological constant $`\mathrm{\Lambda }`$. It has not been well understood why the continuum of solutions for negative $`\mathrm{\Lambda }`$ become a discrete set of solutions in the $`\mathrm{\Lambda }0`$ limit, and remain discrete for all $`\mathrm{\Lambda }>0`$. Just as fig. 8 shows for the black hole solutions, fig. 4 and fig. 11 shows the spectrum in mass vs. magnetic charge $`Q_M`$ plane for a give $`\mathrm{\Lambda }`$. The width of each branch for a given $`\mathrm{\Lambda }`$ gets smaller as $`\mathrm{\Lambda }`$ approaches zero. Fig. 11 indicates that as $`\mathrm{\Lambda }0`$, the branches collapse to one point, the BK solution, as the continuum of solutions vanishes. It is still unknown mathematically why and how this occurs.
We would like to point out that there is a fractal structure in the moduli space of the solutions. This is most clearly seen in the parameter $`b`$ v.s. mass $`M`$ plot as displayed in fig. 12. As $`\mathrm{\Lambda }`$ becomes smaller, a new branch appears. The shape of branches has approximate self-similarity. Similarly, in fig. 13 the magnetic charge $`Q_M`$ is plotted against $`b`$. Delicate structure is observed near the critical $`b=b_c`$ which signifies the critical solution discussed in Section 5.3 (c). There may be some connection between the limiting point in the monopole spectrum and the critical solution.
## 8 Stability
It has been shown that the soliton and black hole solutions in asymptotically Minkowski and de Sitter space, which necessarily have at least one node in $`w(r)`$, are unstable . In contrast, the monopole and black hole solutions in the asymptotically anti-de Sitter space with no node in $`w(r)`$ are stable for $`u=0`$. One expects the presence of the electric field not to change the stability of the solitons and black hole configurations.
In this section we give a detailed discussion for establishing the stability. We shall find that in asymptotically anti-de Sitter space the boundary condition for the resultant Schrödinger problem becomes subtle, and that the previous argument given in asymptotically Minkowski space needs elaboration.
### 8.1 Perturbation equations
We consider small time-dependent perturbations to the static solutions to the coupled EYM equations. In the static solutions $`\nu (r)=\stackrel{~}{w}(r)=0`$. In the general ansatz, (18) and (28) we set
$`u(r,t)`$ $`=`$ $`u(r)+\delta u(r,t)`$ (103)
$`w(r,t)`$ $`=`$ $`w(r)+\delta w(r,t)`$ (104)
$`\stackrel{~}{w}(r,t)`$ $`=`$ $`\delta \stackrel{~}{w}(r,t)`$ (105)
$`\nu (r,t)`$ $`=`$ $`\delta \nu (r,t)`$ (106)
$`p(r,t)`$ $`=`$ $`p(r)+\delta p(r,t)`$ (107)
$`H(r,t)`$ $`=`$ $`H(r)+\delta H(r,t)`$ (108)
and $`m(r,t)=m(r)+\delta m(r,t)`$. Substituting (108) into the Yang-Mills equation (34) and retaining only terms linear in perturbations, one finds
$`\left[r^2u^{}\delta p+r^2p(\delta u^{}\delta \dot{\nu })\right]^{}{\displaystyle \frac{2uw^2p}{H}}\left({\displaystyle \frac{\delta p}{p}}{\displaystyle \frac{\delta H}{H}}\right){\displaystyle \frac{2pw}{H}}(w\delta u+2u\delta w+\delta \dot{\stackrel{~}{w}})=0`$ (109)
$`r^2p(\delta \dot{u}^{}\delta \ddot{\nu })+r^2u^{}\delta \dot{p}{\displaystyle \frac{2H}{p}}\left[w\delta \stackrel{~}{w}^{}w^{}\delta \stackrel{~}{w}+w^2\delta \nu \right]=0`$ (110)
$`\left[{\displaystyle \frac{w^{}H}{p}}\left({\displaystyle \frac{\delta H}{H}}{\displaystyle \frac{\delta p}{p}}\right)+{\displaystyle \frac{H}{p}}\delta w^{}\right]^{}+{\displaystyle \frac{pu}{H}}(u\delta w+2w\delta u+2\delta \dot{\stackrel{~}{w}}){\displaystyle \frac{p}{H}}\delta \ddot{w}`$ (111)
$`+{\displaystyle \frac{wu^2p}{H}}\left({\displaystyle \frac{\delta p}{p}}{\displaystyle \frac{\delta H}{H}}\right){\displaystyle \frac{w(1w^2)}{r^2p^2}}\delta p+{\displaystyle \frac{13w^2}{r^2p}}\delta w=0`$ (112)
$`\left[{\displaystyle \frac{H}{p}}(\delta \stackrel{~}{w}^{}+w\delta \nu )\right]^{}+{\displaystyle \frac{p}{H}}(u^2\delta \stackrel{~}{w}2u\delta \dot{w}w\delta \dot{u}\delta \ddot{\stackrel{~}{w}})`$ (113)
$`{\displaystyle \frac{uwp}{H}}\left({\displaystyle \frac{\delta \dot{p}}{p}}{\displaystyle \frac{\delta \dot{H}}{H}}\right)+{\displaystyle \frac{(1w^2)}{r^2p}}\delta \stackrel{~}{w}+{\displaystyle \frac{H}{p}}w^{}\delta \nu =0.`$ (114)
The Einstein equations (45) and (47) yield
$`\delta m^{}=v\{r^2p^2u^{}(\delta u^{}\delta \dot{\nu })+(r^2u^2+{\displaystyle \frac{2u^2w^2}{H}})p\delta p+(w^2{\displaystyle \frac{p^2u^2w^2}{H^2}})\delta H`$ (115)
$`+2Hw^{}\delta w^{}+{\displaystyle \frac{2uwp^2}{H}}(w\delta u+u\delta w+\delta \dot{\stackrel{~}{w}}){\displaystyle \frac{2w(1w^2)}{r^2}}\delta w\}`$ (116)
$`\left({\displaystyle \frac{\delta p}{p}}\right)^{}={\displaystyle \frac{4v}{r}}\left\{w^{}\delta w^{}+{\displaystyle \frac{p^2u^2w^2}{H^2}}\left({\displaystyle \frac{\delta p}{p}}{\displaystyle \frac{\delta H}{H}}\right)+{\displaystyle \frac{p^2uw}{H^2}}(w\delta u+u\delta w+\delta \dot{\stackrel{~}{w}})\right\}`$ (117)
$`\delta \dot{H}={\displaystyle \frac{4vH}{r}}\left\{w^{}\delta \dot{w}u(w^{}\delta \stackrel{~}{w}w\delta \stackrel{~}{w}^{})+uw^2\delta \nu \right\}.`$ (118)
There is residual gauge invariance specified by a gauge function $`\mathrm{\Omega }(r,t)`$ in (12) and (22). Making use of this freedom, one can always set either $`\delta u(r,t)=0`$ or $`\delta \nu (r,t)=0`$.
### 8.2 Stability analysis
In examining time-dependent fluctuations around monopole solutions for which $`u(r)=0`$, it is convenient to work in the $`\delta u(t,r)=0`$ gauge. Eqs. (109), (110), and (114) become
$`(r^2p\delta \nu )^{}={\displaystyle \frac{2pw}{H}}\delta \stackrel{~}{w}`$ (119)
$`r^2p\delta \ddot{\nu }+{\displaystyle \frac{2H}{p}}(w\delta \stackrel{~}{w}^{}w^{}\delta \stackrel{~}{w}+w^2\delta \nu )=0`$ (120)
$`\left[{\displaystyle \frac{H}{p}}(\delta \stackrel{~}{w}^{}+w\delta \nu )\right]^{}{\displaystyle \frac{p}{H}}\delta \ddot{\stackrel{~}{w}}+{\displaystyle \frac{1w^2}{r^2p}}\delta \stackrel{~}{w}+{\displaystyle \frac{H}{p}}w^{}\delta \nu =0,`$ (121)
whereas Eqs. (112), (116), (117), and (118) become
$`\left[{\displaystyle \frac{H}{p}}\left\{w^{}\left({\displaystyle \frac{\delta H}{H}}{\displaystyle \frac{\delta p}{p}}\right)+\delta w^{}\right\}\right]^{}{\displaystyle \frac{p}{H}}\delta \ddot{w}{\displaystyle \frac{w(1w^2)}{r^2p^2}}\delta p+{\displaystyle \frac{13w^2}{r^2p}}\delta w=0`$ (122)
$`\delta m^{}=v\left\{w^2\delta H+2Hw^{}\delta w^{}{\displaystyle \frac{2w(1w^2)}{r^2}}\delta w\right\}`$ (123)
$`\left({\displaystyle \frac{\delta p}{p}}\right)^{}={\displaystyle \frac{4v}{r}}w^{}\delta w^{}`$ (124)
$`\delta H={\displaystyle \frac{4v}{r}}Hw^{}\delta w.`$ (125)
Notice that Eqs. (119) - (121) involve only $`\delta \nu `$ and $`\delta \stackrel{~}{w}`$, defining the odd parity group, whereas Eqs. (122) - (125) involve only $`\delta w`$, $`\delta H`$, and $`\delta p`$, defining the even parity group. The number of the equations is larger than the number of the unknown functions. Indeed, one equation in each group follows from the others.
To derive the equation for each unknown function in a closed form, we introduce the tortoise radial coordinate $`\rho `$ by
$$\frac{d\rho }{dr}=\frac{p}{H}$$
(126)
with which the equations for $`w`$, $`p`$, and $`m`$ become
$`{\displaystyle \frac{d^2w}{d\rho ^2}}`$ $`=`$ $`{\displaystyle \frac{H}{r^2p^2}}w(1w^2)`$ (127)
$`{\displaystyle \frac{dp}{d\rho }}`$ $`=`$ $`{\displaystyle \frac{2vp^2}{rH}}\left({\displaystyle \frac{dw}{d\rho }}\right)^2`$ (128)
$`{\displaystyle \frac{dm}{d\rho }}`$ $`=`$ $`v\left\{p\left({\displaystyle \frac{dw}{d\rho }}\right)^2+{\displaystyle \frac{H}{2r^2p}}(1w^2)^2\right\}.`$ (129)
The range of $`\rho `$ is finite, $`0\rho \rho _{\mathrm{max}}`$, since $`pp_0`$ and $`H|\mathrm{\Lambda }|r^2/3`$ as $`r\mathrm{}`$:
$$\rho =\{\begin{array}{cc}r\hfill & \text{for }r0\hfill \\ \multicolumn{2}{c}{}\\ \rho _{\mathrm{max}}\frac{3p_0}{\left|\mathrm{\Lambda }\right|r}\hfill & \text{for }r\mathrm{}\hfill \end{array}$$
(130)
where $`p=p_0+\mathrm{O}(1/r)`$.
In the odd parity group Eq. (119) expresses $`\delta w`$ in terms of $`\delta \nu `$. Substituting it into Eq. (120) and making use of (127), one finds
$`\left\{{\displaystyle \frac{d^2}{d\rho ^2}}+U_\beta (\rho )\right\}\beta =\omega ^2\beta `$ (131)
$`U_\beta ={\displaystyle \frac{H}{r^2p^2}}(1+w^2)+{\displaystyle \frac{2}{w^2}}\left({\displaystyle \frac{dw}{d\rho }}\right)^2,`$ (132)
$`\delta \nu ={\displaystyle \frac{w}{r^2p}}\beta ,\delta \stackrel{~}{w}={\displaystyle \frac{1}{2w}}{\displaystyle \frac{d}{d\rho }}(w\beta ).`$ (133)
Here we have supposed fluctuations to be harmonic: $`\delta \nu (r,t)=e^{i\omega t}\delta \nu (\rho )`$ and $`\delta \stackrel{~}{w}(r,t)=e^{i\omega t}\delta \stackrel{~}{w}(\rho )`$. Eq. (121) follows from (119), (120), and (127).
In the even parity group, (124) and (125) express $`\delta p`$ and $`\delta H`$ ($`\delta m`$) in terms of $`\delta w`$. Eq. (123) automatically follows from (124), (125), (127) and (128). Eq. (122) becomes, with the use of (124), (125) and (127),
$`\left\{{\displaystyle \frac{d^2}{d\rho ^2}}+U_w(\rho )\right\}\delta w=\omega ^2\delta w,`$ (134)
$`U_w={\displaystyle \frac{H}{r^2p^2}}(3w^21)+4v{\displaystyle \frac{d}{d\rho }}\left[{\displaystyle \frac{p}{rH}}\left({\displaystyle \frac{dw}{d\rho }}\right)^2\right]`$ (135)
Again harmonic fluctuations $`\delta w(r,t)=e^{i\omega t}\delta w(\rho )`$ are supposed.
Eqs. (133) and (135) have the same form as the Schrödinger equation on a one-dimensional interval. Both of the potentials $`U_\beta `$ and $`U_w`$ are singular at $`\rho =0`$, behaving as $`+2/\rho ^2`$. $`U_\beta `$ has an additional singularity if $`w`$ has a zero at $`\rho _k`$; $`U_\beta +2/(\rho \rho _k)^2`$.
The integrated energy-momentum density $`T^{ab}\sqrt{g}d^3x`$ due to fluctuations must remain finite. At the origin $`r=0`$ it implies that $`\delta w=\delta \stackrel{~}{w}=0`$ whereas $`\delta \nu =\mathrm{O}(1)`$. Taking advantage of the general coordinate invariance, one can impose $`\delta p=0`$ and $`\delta H=2\delta m/r=0`$. At $`r\mathrm{}`$, $`\delta w^{},\delta \stackrel{~}{w}^{},\delta \nu =\mathrm{O}(r^2)`$. These are mild boundary conditions. One can impose more strict conditions such as the regularity at $`r=0`$ and vanishing at $`r=\mathrm{}`$. As physical perturbations we demand that all $`\delta w`$, $`\delta \stackrel{~}{w}`$, $`\delta \nu `$, $`\delta H`$, and $`\delta p`$ vanish at $`r=\mathrm{}`$.
In Eq. (133) the potential $`U_\beta (\rho )`$ is positive definite. However, this does not necessarily mean that the eigenvalue $`\omega ^2`$ is positive definite. It depends on the boundary condition. Clearly $`\beta (0)=0`$. At $`\rho =\rho _{\mathrm{max}}`$, $`\delta \nu =\delta \stackrel{~}{w}=0`$ so that $`\beta ^{}+h\beta =0`$ where $`h=w^{}/w`$. Note that
$$h=\frac{1}{w}\frac{dw}{d\rho }|_{\rho _{\mathrm{max}}}=\frac{|\mathrm{\Lambda }|}{3p_0}\frac{r^2}{w}\frac{dw}{dr}|_{r=\mathrm{}}=\frac{|\mathrm{\Lambda }|w_1}{3p_0w_0}$$
(136)
where $`w_j`$’s are defined in (59). For the monopole configurations with no nodes in $`w`$, $`h<0`$ ($`h>0`$) when $`w`$ is monotonically decreasing (increasing).
Following Courant and Hilbert , we define
$`𝒟(\phi ;h)={\displaystyle _0^{\rho _{\mathrm{max}}}}𝑑\rho \left\{\phi ^{}(\rho )^2+U_\beta (\rho )\phi (\rho )^2\right\}+h\phi (\rho _{\mathrm{max}})^2`$ (137)
$`𝒩(\phi )={\displaystyle _0^{\rho _{\mathrm{max}}}}𝑑\rho \phi (\rho )^2.`$ (138)
If $`w(r)`$ is nodeless, then $`U_\beta (\rho )`$ is regular on the interval except at $`\rho =0`$. The equation implies that $`\beta =\mathrm{O}(\rho ^2)`$ near the origin. In this case, for an eigenfunction $`\beta (\rho )`$ in (133) satisfies
$$\omega ^2=\frac{𝒟(\beta )}{𝒩(\beta )}.$$
(139)
It follows immediately that all eigenvalues $`\omega ^2`$ are positive definite if $`h0`$ so that the solution is stable against small odd-parity perturbations.
For $`h<0`$ more careful analysis is necessary. The lowest eigenvalue $`\omega ^2\lambda _1`$ in the eigenvalue equation (133) is exactly the lower bound of the set of values assumed by the functional $`𝒟(\phi ,h)`$, where $`\phi `$ is any function continuous on the interval $`[0,\rho _{\mathrm{max}}]`$ with piecewise continuous derivatives satisfying $`\phi (0)=1`$ and $`𝒩(\phi )=1`$.
$$\lambda _1(h)=\mathrm{min}_\phi 𝒟(\phi ;h).$$
(140)
If $`\lambda _1>0`$, then the solution is stable against odd-parity perturbations. Suppose that $`\phi _1(\rho )`$ saturates the lower bound for $`h_1`$: $`\lambda _1(h_1)=𝒟(\phi _1;h_1)`$. As
$`\lambda _1(h_1)`$ $`=`$ $`𝒟(\phi _1;h_2)+(h_1h_2)\phi _1(\rho _{\mathrm{max}})^2`$ (141)
$``$ $`\lambda _1(h_2)+(h_1h_2)\phi _1(\rho _{\mathrm{max}})^2,`$ (142)
$`\lambda _1(h)`$ is a monotonically increasing function of $`h`$. Hence, if $`\lambda _1(h_1)>0`$, then $`\lambda _1(h)>0`$ for $`hh_1`$.
To establish the stability we utilize the residual gauge invariance. There is a zero-mode (with $`\omega ^2=0`$) for Eq. (133) with an appropriate boundary condition $`h_0`$. In the $`\mathrm{\Lambda }=0`$ case the existence of the zero-mode was utilized to prove the instability of the BK and black hole solutions which has at least one node in $`w(r)`$ . Consider the time-independent gauge function $`\mathrm{\Omega }(r)`$ in (12). For $`|\mathrm{\Omega }|1`$, $`\delta \nu =d\mathrm{\Omega }/dr`$ and $`\delta \stackrel{~}{w}=w\mathrm{\Omega }`$. Eq. (119) is satisfied if
$$\frac{d}{d\rho }\left(\frac{r^2p^2}{H}\frac{d\mathrm{\Omega }}{d\rho }\right)=2w^2\mathrm{\Omega }.$$
(143)
As $`\mathrm{\Omega }(0)=0`$, $`\mathrm{\Omega }a\rho +\mathrm{O}(\rho ^3)`$ for $`\rho 0`$. Hence Eq. (143) determines $`\mathrm{\Omega }(r)`$ up to an over-all constant. $`\beta _0=(r^2p^2/wH)(d\mathrm{\Omega }/d\rho )`$ is the zero mode of Eq. (133) and $`U_\beta =\beta _0^{\prime \prime }/\beta _0`$, where $`\mathrm{\Omega }^{}d\mathrm{\Omega }/d\rho `$ etc. In this case $`\mathrm{\Omega }(\rho _{\mathrm{max}})0`$ and $`h_0=\beta _0^{}/\beta _0|_{\rho _{\mathrm{max}}}`$ differs from $`h`$ in the eigenvalue problem under consideration. If $`h>h_0`$, then $`\lambda _1(h)>0`$, establishing the stability. As $`d(w\beta _0)/d\rho =2w^2\mathrm{\Omega }`$,
$$h=h_0+\frac{2w^2H\mathrm{\Omega }}{r^2p^2\mathrm{\Omega }^{}}.$$
(144)
Nodeless solutions ($`w>0`$) are stable if $`\mathrm{\Omega }/\mathrm{\Omega }^{}>0`$ at $`\rho _{\mathrm{max}}`$.
Solving (143) numerically, we have determined $`\mathrm{\Omega }(\rho )`$ to find that indeed $`h>h_0`$ for nodeless solutions. This analysis also shows that $`h`$ becomes exactly $`h_0`$ for the configuration with $`w(r=\mathrm{})=0`$. In this limiting case the zero mode is not normalizable; it diverges as $`(\rho \rho _{\mathrm{max}})^1`$.
This is a general behavior. When $`w`$ has a node at $`\rho _k<\rho _{\mathrm{max}}`$, there appears a negative $`\omega ^2`$ mode which behaves as $`(\rho \rho _k)^1`$ near $`\rho _k`$.
If $`w(r)`$ has $`n`$ nodes, i.e. $`w(r_j)=0`$ ($`j=1,\mathrm{},n`$), the potential $`U_\beta `$ develops $`(\rho \rho _j)^2`$ singularities. Volkov et al. have shown for the BK solutions in the $`\mathrm{\Lambda }=0`$ case that there appear exactly $`n`$ negative eigenmodes ($`\omega ^2<0`$) if $`w`$ has $`n`$ nodes . A similar conclusion has been obtained for black hole solutions as well . Their argument needs elaboration in the $`\mathrm{\Lambda }<0`$ case, however.
To investigate the eigenvalue spectrum of (133) in this case, it is convenient to consider the dual equation as was done in . One can write the Schrödinger equation in (133) as
$`Q_+Q_{}\phi _n=\lambda _n\phi _n`$ (145)
$`Q_\pm =\pm {\displaystyle \frac{d}{d\rho }}+{\displaystyle \frac{\beta _0^{}}{\beta _0}}`$ (146)
where $`\beta _0(\rho )`$ is the zero mode described above. $`\phi _n=\mathrm{O}(\rho ^2)`$ near $`\rho =0`$ and $`\phi _n^{}+h\phi _n=0`$ at $`\rho _{\mathrm{max}}`$. The dual equation is given by
$`Q_{}Q_+\stackrel{~}{\phi }_n=\left\{{\displaystyle \frac{d^2}{d\rho ^2}}+\stackrel{~}{U}_\beta \right\}\stackrel{~}{\phi }_n=\lambda _n\stackrel{~}{\phi }_n`$ (147)
$`\stackrel{~}{U}_\beta ={\displaystyle \frac{H}{r^2p^2}}(1+w^2){\displaystyle \frac{8ww^{}H\mathrm{\Omega }}{r^2p^2\mathrm{\Omega }^{}}}+8\left({\displaystyle \frac{w^2H\mathrm{\Omega }}{r^2p^2\mathrm{\Omega }^{}}}\right)^2.`$ (148)
$`\phi _n`$ and $`\stackrel{~}{\phi }_n`$ are related to each other by
$$\stackrel{~}{\phi }_n=\{\begin{array}{cc}Q_{}\phi _n\hfill & \text{for }\lambda _n0\hfill \\ \phi _{n}^{}{}_{}{}^{1}=\beta _{0}^{}{}_{}{}^{1}\hfill & \text{for }\lambda _n=0\hfill \end{array}$$
(149)
However, the boundary condition for $`\stackrel{~}{\phi }_n(\rho )`$ depends on $`\lambda _n`$:
$`\stackrel{~}{\phi }_n(0)=0`$ (150)
$`\stackrel{~}{\phi }_n^{}(\rho _{\mathrm{max}})+\stackrel{~}{h}_n\stackrel{~}{\phi }_n(\rho _{\mathrm{max}})=0`$ (151)
$`\stackrel{~}{h}_n=h_0{\displaystyle \frac{\lambda _n}{hh_0}}.`$ (152)
The advantage of considering the dual equation is that the dual potential $`\stackrel{~}{U}_\beta (\rho )`$ is regular except at $`\rho =0`$ where it behaves as $`+6/\rho ^2`$. However, the eigenvalue $`\lambda _n`$ has to be determined self-consistently such that the boundary condition (152) is satisfied. We have determined $`\lambda _n`$’s numerically for the monopole configurations in the lower branch in fig. 4. The first, second and third eigenvalues are displayed in fig. 14. One sees that $`\lambda _1>0`$ for the nodeless configurations, but the unstable mode develops when $`w`$ has a node.
The wave function of the unstable mode in the original equation, not in the dual equation, diverges at the zeroes of $`w(r)`$. In other words, the instability sets in around the zeroes of $`w(r)`$. The potential $`U_\beta (\rho )`$ and $`\stackrel{~}{U}_\beta (\rho )`$ for the solution at $`b=0.0025`$ are plotted in fig. 15. At the node of $`w`$, $`U_\beta `$ diverges, but $`\stackrel{~}{U}_\beta `$ remains finite. The wave function $`\stackrel{~}{\phi }_1(\rho )`$ of the lowest eigenvalue ($`\lambda _1=0.0033)`$ and the corresponding $`\phi _1(\rho )`$ also have been plotted in fig. 15.
For even-parity perturbation the potential $`U_w(\rho )`$ in Eq. (135) is not positive definite. The first term in $`U_w`$ becomes negative for $`w^2<1/3`$. The second term also can become negative when $`w^{}`$ vanishes at finite $`r`$. We have solved the Schrödinger equation (135) numerically for typical monopole solutions, and found that for the solutions with no node in $`w(r)`$, the eigenvalues $`\omega ^2`$ are always positive even if $`w(r=\mathrm{})<1/\sqrt{3}`$.
Hence we have established the stability of the monopole solutions with no node in $`w(r)`$.
## 9 Summary
New monopole, dyon, and black hole solutions to the Einstein-Yang-Mills equations have been found in asymptotically anti-de Sitter space. The solutions with no node in the non-Abelian field strengths are shown to be stable against spherically symmetric perturbations. The non-trivial boundary condition plays a crucial role in developing the instability for solutions with nodes. The stability of nodeless dyon solutions need to be established.
Though electric and magnetic charges of monopole and dyon solutions are not quantized in classical theory, they are expected to be quantized in quantum theory. If this is the case, then at least solutions with the smallest charge would become absolutely stable.
We have also found the critical spacetime solutions which end at finite $`r`$. These solutions may have connections to black hole solutions, though more detailed study is necessary.
The solutions found in the present paper may have profound consequences in the evolution of the early universe which may have gone through the anti-de Sitter phase. We hope to report on these subjects in future publications.
Acknowledgments
This work was supported in part by the U.S. Department of Energy under contracts DE-FG02-94ER-40823, DE-FG02-87ER40328 and DE-AC02-98CH10886. |
warning/0002/hep-th0002003.html | ar5iv | text | # A Generalization of the Submodel of Nonlinear 𝐂𝑃¹ Models
## 1 Introduction
Integrable theories in a space-time of two dimension have achieved remarkable development and great many integrable models exist. But in the dimensions greater than two, we do not have so many interesting models because it is difficult to extend the concepts of integrability to higher dimensions.
In these circumstances, O. Alvarez, L. A. Ferreira and J. S. Guillen proposed a new approach to higher dimensional integrable theories (see also ). Instead of higher dimensional models themselves , they considered their submodels to construct integrable models in the sense of possessing an infinite number of conserved currents. They applied their theories to nonlinear $`𝐂P^1`$ model in $`(1+2)`$-dimensions in .
Then we calculated an infinite number of conserved currents explicitly in the submodel of nonlinear $`𝐂P^1`$ model in $`(1+2)`$-dimensions , (see also ). Furthermore, we generalized the definition of submodels to nonlinear Grassmann sigma models and constructed an infinite number of conserved currents and a wide class of exact solutions ,. (later Ferreira and Leite generalized them to homogeneous-space models ). The idea of submodels was also applied in ,.
Moreover, we also generalized them to another direction in . In view of the fact that the $`𝐂P^1`$-submodel
$$^\mu _\mu u=0\text{and}^\mu u_\mu u=0,$$
$$\text{for}u:M^{1+n}𝐂$$
is equivalent to
$$\mathrm{}_2u=0\text{and}\mathrm{}_2(u^2)=0,$$
we defined a system of $`p`$-th order ($`p=2,3,\mathrm{}`$) nonlinear partial differential equations (PDE) by generalizing the $`𝐂P^1`$-submodel, which have an infinite number of conserved currents and a wide class of exact solutions;
$$\mathrm{}_p(u^k)\left(\frac{^p}{x_0^p}\underset{j=1}{\overset{n}{}}\frac{^p}{x_j^p}\right)(u^k)=0\text{for}1kp.$$
Hereafter, we call this system of PDE the $`p`$-submodel for short. In constructing conserved currents of the $`p`$-submodel, we defined differential operators by using the Bell polynomials. We investigated the reason why such a form of operators appeared. Then we found a kind of “symbol structure” for the operators and could define a wider class of PDE than the $`p`$-submodel.
In this paper, we define a new system of PDE including the $`p`$-submodel and construct a Bäcklund-like transformation of solutions and an infinite number of conserved currents by using the Bell polynomials.
## 2 Bell Polynomials
Firstly, we prepare a mathematical tool which plays an important role in our following theory.
###### Definition 2.1.
Let $`g(x)`$ be a smooth function and $`z`$ a complex parameter. Put $`g_r_x^rg(x)`$. We define the Bell polynomials of degree $`n`$ ():
$$F_n(zg)=F_n(zg_1,\mathrm{},zg_n)\text{e}^{zg(x)}_x^n\text{e}^{zg(x)}.$$
(2.1)
The generating function of Bell polynomials is
$$\mathrm{exp}\left\{z\underset{j=1}{\overset{\mathrm{}}{}}\frac{g_j}{j!}t^j\right\}=\underset{n=0}{\overset{\mathrm{}}{}}\frac{F_n(zg)}{n!}t^n.$$
(2.2)
By (2.2), we can write them explicitly as follows:
$$F_n(zg_1,\mathrm{},zg_n)=\underset{\genfrac{}{}{0pt}{}{k_1+2k_2+\mathrm{}+nk_n=n}{k_10,\mathrm{},k_n0}}{}\frac{n!}{k_1!\mathrm{}k_n!}\left(\frac{zg_1}{1!}\right)^{k_1}\left(\frac{zg_2}{2!}\right)^{k_2}\mathrm{}\left(\frac{zg_n}{n!}\right)^{k_n}.$$
(2.3)
For example,
$$F_0=1,F_1=zg_1,F_2=zg_2+z^2g_1^2.$$
These polynomials are used in the differential calculations of composite functions.
###### Lemma 2.2.
We have a recursion formura for the Bell polynomials.
$$F_{n+1}(zg)=\left\{\underset{r=1}{\overset{n}{}}g_{r+1}\frac{}{g_r}+zg_1\right\}F_n(zg).$$
(2.4)
Now, we define the Bell matrix (we have adopted notations in );
$$B_{nj}=B_{nj}[g]=B_{nj}(g_1,\mathrm{},g_{nj+1})$$
by the following equation
$$F_n(zg_1,\mathrm{},zg_n)=\underset{j0}{}z^jB_{nj}(g_1,\mathrm{},g_{nj+1}).$$
(2.5)
Note that
$$B_{n0}=\delta _{n0},B_{nj}=0(n<j).$$
(2.6)
An important formula for constructing a Bäcklund-like transformation is as follows;
###### Lemma 2.3.
$$B_{jk}[f(g(x))]=\underset{n=k}{\overset{j}{}}B_{jn}[g]B_{nk}[f]$$
(2.7)
## 3 A New System of Higher Order Equations
In this section, we generalize the equations of motion of the $`𝐂P^1`$-submodel to higher order. Hereafter, we use a notation of Minkowski summation as follows:
$$\underset{\mu }{}{}_{}{}^{}A_{\mu }^{}A_0\underset{j=1}{\overset{n}{}}A_j.$$
(3.1)
Now, given $`p=2,3,\mathrm{}`$ and $`i=0,1,\mathrm{},[(p1)/2]`$, where \[ \] means the Gauss’s symbol, we define a system of higher order nonlinear PDE as follows:
###### Definition 3.1.
$$\underset{\mu }{}{}_{}{}^{}_{\mu }^{pi}(u^k)_\mu ^i(\overline{u}^l)=0$$
(3.2)
for $`k=1,\mathrm{},pi`$, $`l=0,\mathrm{},i`$.
We call this system of PDE the $`(p,i)`$-submodel.
For example, $`(p,i)=(2,0);`$
$$\underset{\mu }{}{}_{}{}^{}_{\mu }^{2}u=0,\underset{\mu }{}{}_{}{}^{}_{\mu }^{2}(u^2)=0.$$
(3.3)
We note that (3.3) is equivalent to the $`𝐂P^1`$-submodel.
$`(p,i)=(3,0);`$
$$\underset{\mu }{}{}_{}{}^{}_{\mu }^{3}u=0,\underset{\mu }{}{}_{}{}^{}_{\mu }^{3}(u^2)=0,\underset{\mu }{}{}_{}{}^{}_{\mu }^{3}(u^3)=0,$$
(3.4)
$`(p,i)=(3,1);`$
$$\underset{\mu }{}{}_{}{}^{}_{\mu }^{2}u_\mu \overline{u}=0,\underset{\mu }{}{}_{}{}^{}_{\mu }^{2}(u^2)_\mu \overline{u}=0.$$
(3.5)
## 4 Conserved Currents for the System of Higher Order Equations
In this section, we construct conserved currents for the system of PDE (3.2).
let $`g(x)`$, $`\overline{g}(x)`$ be smooth functions and $`z`$, $`\overline{z}`$ complex parameters. Put $`𝒫_\text{B}`$ the vector space over $`𝐂`$ spaned by the products of two Bell polynomials $`F_n(zg)\overline{F}_m(\overline{z}\overline{g})`$. (We use a notation $`\overline{F}_m(\overline{z}\overline{g})`$ instead of $`F_m(\overline{z}\overline{g})`$ for convenience.) We consider a linear map
$$\mathrm{\Phi }:𝒫_\text{B}𝐂[\xi ,\overline{\xi }],$$
(4.1)
$$\mathrm{\Phi }(F_n(zg)\overline{F}_m(\overline{z}\overline{g}))=\xi ^n\overline{\xi }^m.$$
(4.2)
###### Remark 4.1.
The map $`\mathrm{\Phi }`$ is considered as the tensor product of a “symbol map”
$$F_n(zg)=\text{e}^{zg(x)}_x^n\text{e}^{zg(x)}\text{e}^{zg(x)}\xi ^n\text{e}^{zg(x)}=\xi ^n.$$
By this map, $`𝒫_\text{B}`$ is linear isomorphic to $`𝐂[\xi ,\overline{\xi }]`$. Now, we define an operator
$$\underset{r=1}{\overset{\mathrm{}}{}}\left(g_{r+1}\frac{}{g_r}+\overline{g}_{r+1}\frac{}{\overline{g}_r}\right)+zg_1+\overline{z}\overline{g}_1$$
(4.3)
This operator is well-defined on $`𝒫_\text{B}`$ and we have
$$(F_n(zg)\overline{F}_m(\overline{z}\overline{g}))=F_{n+1}(zg)\overline{F}_m(\overline{z}\overline{g})+F_n(zg)\overline{F}_{m+1}(\overline{z}\overline{g}).$$
(4.4)
In fact, because $`F_n`$ is $`n`$-variable polynomial and on account of (2.4),
$`(F_n(zg)\overline{F}_m(\overline{z}\overline{g}))`$
$`=`$ $`\left\{{\displaystyle \underset{r=1}{\overset{n}{}}}g_{r+1}{\displaystyle \frac{}{g_r}}+{\displaystyle \underset{r=1}{\overset{m}{}}}\overline{g}_{r+1}{\displaystyle \frac{}{\overline{g}_r}}+zg_1+\overline{z}\overline{g}_1\right\}F_n(zg)\overline{F}_m(\overline{z}\overline{g})`$
$`=`$ $`{\displaystyle \underset{r=1}{\overset{n}{}}}g_{r+1}{\displaystyle \frac{F_n(zg)}{g_r}}\overline{F}_m(\overline{z}\overline{g})+zg_1F_n(zg)\overline{F}_m(\overline{z}\overline{g})`$
$`+F_n(zg){\displaystyle \underset{r=1}{\overset{m}{}}}\overline{g}_{r+1}{\displaystyle \frac{\overline{F}_m(\overline{z}\overline{g})}{\overline{g}_r}}+\overline{z}\overline{g}_1F_n(zg)\overline{F}_m(\overline{z}\overline{g})`$
$`=`$ $`F_{n+1}(zg)\overline{F}_m(\overline{z}\overline{g})+F_n(zg)\overline{F}_{m+1}(\overline{z}\overline{g}).`$
Because of (4.4), we have
$$\mathrm{\Phi }\mathrm{\Phi }^1=(\xi +\overline{\xi }),$$
(4.6)
where the right hand side of (4.6) means the multiplication operator. By using the linear isomorphism $`\mathrm{\Phi }`$, we can identify
$$F_n\overline{F}_m\text{with}\xi ^n\overline{\xi }^m\text{and}\text{with}(\xi +\overline{\xi }).$$
(4.7)
Nextly, we choose an $`\mu \{0,\mathrm{},n\}`$ and put $`x=x_\mu `$,
$`g(x_\mu )=u(x_0,\mathrm{},x_\mu ,\mathrm{},x_n)`$. Then, we have $`g_r=_\mu ^ru`$. We set $`F_{n,\mu }`$ as
$`F_{n,\mu }`$ $``$ $`:F_n(zg_1,\mathrm{},zg_n)|_{z=\frac{}{u}}:`$ (4.8)
$`=`$ $`:F_n(_\mu u{\displaystyle \frac{}{u}},_\mu ^2u{\displaystyle \frac{}{u}},\mathrm{},_\mu ^nu{\displaystyle \frac{}{u}}):`$
$`=`$ $`{\displaystyle \underset{\genfrac{}{}{0pt}{}{k_1+2k_2+\mathrm{}+nk_n=n}{k_10,\mathrm{},k_n0}}{}}{\displaystyle \frac{n!}{k_1!\mathrm{}k_n!}}\left({\displaystyle \frac{_\mu u}{1!}}\right)^{k_1}\left({\displaystyle \frac{_\mu ^2u}{2!}}\right)^{k_2}\mathrm{}\left({\displaystyle \frac{_\mu ^nu}{n!}}\right)^{k_n}\left({\displaystyle \frac{}{u}}\right)^{k_1+k_2+\mathrm{}+k_n}`$
and $`\overline{F}_{n,\mu }`$ its complex conjugate of $`F_{n,\mu }`$, where $`::`$ means the normal ordering.
###### Lemma 4.2.
$$_\mu :F_{n,\mu }\overline{F}_{m,\mu }:f(u,\overline{u})=:(F_n(zg)\overline{F}_m(\overline{z}\overline{g}))|_{z=\frac{}{u}}:f(u,\overline{u}).$$
(4.11)
where $`f=f(u,\overline{u})`$ is any function in $`C^{n+m+1}`$-class.
proof: If $`_\mu `$ acts on a functional of the form
$$h(u,_\mu u,\mathrm{},_\mu ^nu;\overline{u},_\mu \overline{u},\mathrm{},_\mu ^m\overline{u}),$$
(4.12)
we can write
$$_\mu =\underset{r=1}{\overset{n}{}}_\mu ^{r+1}u\frac{}{(_\mu ^ru)}+\underset{r=1}{\overset{m}{}}_\mu ^{r+1}\overline{u}\frac{}{(_\mu ^r\overline{u})}+_\mu u\frac{}{u}+_\mu \overline{u}\frac{}{\overline{u}}.$$
(4.13)
On the other hand, in view of (4) and since $`g_r=_\mu ^ru`$, $`z=\frac{}{u}`$, we have proved the lemma. ∎
Moreover, the $`(p,i)`$-submodel is expressed by using $`F_{n,\mu }`$s, namely,
###### Lemma 4.3.
the $`(p,i)`$-submodel is equivalent to
$$\underset{\mu }{}{}_{}{}^{}:F_{pi,\mu }\overline{F}_{i,\mu }:=0.$$
(4.14)
proof: We note that
$`{\displaystyle \underset{\mu }{}}{}_{}{}^{}_{\mu }^{pi}(u^k)_\mu ^i(\overline{u}^l)`$
$`=`$ $`{\displaystyle \underset{\mu }{}}{\displaystyle {}_{}{}^{}\underset{j_1=1}{\overset{pi}{}}}B_{pi,j_1}(g_1,\mathrm{},g_{pij_1+1})\left({\displaystyle \frac{}{u}}\right)^{j_1}(u^k)`$
$`\times {\displaystyle \underset{j_2=0}{\overset{i}{}}}B_{i,j_2}(\overline{g}_1,\mathrm{},\overline{g}_{ij_2+1})\left({\displaystyle \frac{}{\overline{u}}}\right)^{j_2}(\overline{u}^l)`$
$`=`$ $`{\displaystyle \underset{j_1=1}{\overset{pi}{}}}{\displaystyle \underset{j_2=0}{\overset{i}{}}}j_1!j_2!\left(\begin{array}{cc}k& \\ j_1& \end{array}\right)\left(\begin{array}{cc}l& \\ j_2& \end{array}\right)`$
$`\times {\displaystyle \underset{\mu }{}}{}_{}{}^{}B_{pi,j_1}^{}(g_1,\mathrm{},g_{pij_1+1})B_{i,j_2}(\overline{g}_1,\mathrm{},\overline{g}_{ij_2+1})u^{kj_1}\overline{u}^{lj_2}`$
$`\text{for}k=1,\mathrm{},pi,l=0,\mathrm{},i.`$
Because of this, the $`(p,i)`$-submodel holds if and only if
$`{\displaystyle \underset{\mu }{}}{}_{}{}^{}B_{pi,j_1}^{}(g_1,\mathrm{},g_{pij_1+1})B_{i,j_2}(\overline{g}_1,\mathrm{},\overline{g}_{ij_2+1})=0`$
$`\text{for}j_1=1,\mathrm{},pi,j_2=0,\mathrm{},i,`$
namely
$$\underset{\mu }{}{}_{}{}^{}:F_{pi,\mu }\overline{F}_{i,\mu }:=0.\mathit{}$$
In view of (4.7) and the lemmas above, we can search an infinite number of conserved currents as follows;
For fixed $`(p,i)`$, we consider $`\xi ^{pi}\overline{\xi }^i`$ and $`\xi ^i\overline{\xi }^{pi}`$ in $`𝐂[\xi ,\overline{\xi }]`$ corresponding to the $`(p,i)`$-submodel. Find the polynomials $`p(\xi ,\overline{\xi })`$ such that
$$(\xi +\overline{\xi })p(\xi ,\overline{\xi })=\alpha \xi ^i\overline{\xi }^{pi}+\beta \xi ^{pi}\overline{\xi }^i\text{for some}\alpha ,\beta 𝐂.$$
(4.16)
Then we can decide it uniquely (up to constant). That is
$$p(\xi ,\overline{\xi })=\underset{k=0}{\overset{p12i}{}}(1)^k\xi ^{p1ik}\overline{\xi }^{i+k}.$$
(4.17)
Therefore, if we define the operator
$$V_{(p,i),\mu }\underset{k=0}{\overset{p12i}{}}(1)^k:F_{p1ik,\mu }\overline{F}_{i+k,\mu }:,$$
(4.18)
we obtain the next theorem.
###### Theorem 4.4.
For $`p=2,3,\mathrm{}`$ and $`i=0,1,\mathrm{},[(p1)/2]`$,
$$V_{(p,i),\mu }(f)$$
(4.19)
are conserved currents for the $`(p,i)`$-submodel, where $`f=f(u,\overline{u})`$ is any function in $`C^p`$-class.
For example, corresponding to (3.3), (3.4), and (3.5),
$`V_{(2,0),\mu }(f)`$ $`=`$ $`F_{1,\mu }(f)\overline{F}_{1,\mu }(f)`$ (4.20)
$`=`$ $`_\mu u{\displaystyle \frac{f}{u}}_\mu \overline{u}{\displaystyle \frac{f}{\overline{u}}},`$
$`V_{(3,0),\mu }(f)`$ $`=`$ $`F_{2,\mu }(f):F_{1,\mu }\overline{F}_{1,\mu }:(f)+\overline{F}_{2,\mu }(f)`$
$`=`$ $`_\mu ^2u{\displaystyle \frac{f}{u}}+(_\mu u)^2{\displaystyle \frac{^2f}{u^2}}_\mu u_\mu \overline{u}{\displaystyle \frac{^2f}{u\overline{u}}}+_\mu ^2\overline{u}{\displaystyle \frac{f}{\overline{u}}}+(_\mu \overline{u})^2{\displaystyle \frac{^2f}{\overline{u}^2}},`$
$`V_{(3,1),\mu }(f)`$ $`=`$ $`:F_{1,\mu }\overline{F}_{1,\mu }:(f)`$ (4.22)
$`=`$ $`_\mu u_\mu \overline{u}{\displaystyle \frac{^2f}{u\overline{u}}}.`$
(4.19) is a generalization of the conserved currents for the $`p`$-submodel in .
## 5 Exact Solutions
In spite of higher order equations, we find that the $`(p,i)`$-submodel has a Bäcklund-like transformation of solutions
$$v=f(u)=\underset{i=0}{\overset{\mathrm{}}{}}f_iu^if:𝐂𝐂:\text{ holomorphic}$$
(5.1)
with an infinite number of parameters $`f_i(i=0,1,\mathrm{})`$. This property is similar to the Grassmann submodel .
###### Theorem 5.1.
If $`u`$ is a solution of the $`(p,i)`$-submodel, then, for any holomorphic function $`f`$, $`v=f(u)`$ is also a solution of the $`(p,i)`$-submodel.
proof: Suppose that $`u`$ is a solution of the $`(p,i)`$-submodel. By lemma 4.3, the $`(p,i)`$-submodel is equivalent to
$$\underset{\mu }{}{}_{}{}^{}:F_{pi,\mu }\overline{F}_{i,\mu }:=0,$$
(5.2)
namely
$$\underset{\mu }{}{}_{}{}^{}B_{pi,j}^{}[u]B_{ik}[\overline{u}]=0\text{for}j=1,\mathrm{},pi,k=0,\mathrm{},i.$$
(5.3)
Therefore, by using of (2.7)
$`{\displaystyle \underset{\mu }{}}{}_{}{}^{}B_{pi,j}^{}[f(u)]B_{ik}[\overline{f(u)}]`$ $`=`$ $`{\displaystyle \underset{\mu }{}}{\displaystyle {}_{}{}^{}\underset{n=j}{\overset{pi}{}}}{\displaystyle \underset{m=k}{\overset{i}{}}}B_{pi,n}[u]B_{nj}[f]B_{im}[\overline{u}]B_{mk}[\overline{f}]`$
$`=`$ $`0.\mathit{}`$
For complex numbers $`a_\mu (\mu =0,1,\mathrm{},n)`$ with $`_\mu {}_{}{}^{}a_{\mu }^{pi}\overline{a}_\mu ^i=0`$,
$$u=a_0x_0+\underset{i=1}{\overset{n}{}}a_ix_i$$
(5.4)
are clearly solutions of (3.2). Therefore, we obtain the following corollary.
###### Corollary 5.2.
Let $`f`$ be any holomorphic function. Then
$$f(a_0x_0+\underset{i=1}{\overset{n}{}}a_ix_i)$$
(5.5)
under
$$\underset{\mu }{}{}_{}{}^{}a_{\mu }^{pi}\overline{a}_\mu ^i=0$$
(5.6)
are solutions of the $`(p,i)`$-submodel.
## 6 Discussion
In this paper, we defined a new system of PDE and constructed a Bäcklund-like transformation of solutions and an infinite number of conserved currents by using the Bell polynomials. In particular, we constructed the necessary differential operators by a natural correspondence of Bell polynomials with usual monomials.
Our result is a generalization of that of and gives a simpler method for constructing an infinite number of conserved currents. We remark that the extended Smirnov-Sobolev construction in terms of , which is a method for constructing exact solutions, is also valid for our new system of PDE.
It is important to know the relationship between the conserved currents and the exact solutions of our submodels. Therefore, we need to investigate symmetries, Poisson structures on our system of PDE. A study of these structures is in progress.
## Acknowledgements
The author is very grateful to Kazuyuki Fujii for helpful comments on an earlier draft on this paper and to Yasushi Homma for helpful suggestion in Theorem 5.1. |
warning/0002/astro-ph0002325.html | ar5iv | text | # Multi-colour 𝑃𝐿-relations of Cepheids in the hipparcos catalogue and the distance to the LMC Based on data from the ESA hipparcos astrometry satellite.
## 1 Introduction
Cepheids are important standard candles in determining the extra-galactic distance scale. The results of the hipparcos mission allow, in principle, a calibration of the period-luminosity relation based on the available parallaxes. Feast & Catchpole (1997; hereafter FC) did just that for the $`M_\mathrm{V}\mathrm{log}P`$-relation based on pre-released hipparcos data of 223 Cepheids available to them at that time. Now that the entire catalog has become available (ESA 1997) it is timely to analyse the full sample of Cepheids in it.
In a recent paper, Lanoix et al. (1999, hereafter L99) presented a study similar to ours and they derived the zero points of the $`M_\mathrm{V}\mathrm{log}P`$\- and $`M_\mathrm{I}\mathrm{log}P`$-relations, without, however, discussing the distance to the LMC. We will indicate where the two studies agree and differ.
The paper is organised as follows. In Sect. 2 the sample considered in the present paper is presented, and compared to that in FC and L99. In Sect. 3 several aspects involved in the analysis of parallax data are described, and the method of “reduced parallaxes” is outlined, together with all necessary recipes to obtain the reddening. In Sect. 4 the zero points of the $`PL`$-relations in $`V,I,K`$ and the reddening-free “Wesenheit-index” (e.g. Tanvir 1999, and Eq. (11)) are presented for different selections of the sample, which are discussed in Sect. 5. In Sect. 6 we construct and present the zero points for volume complete samples of stars. In Sect. 7 we describe numerical simulations that are first of all tuned to fit the observed properties of the Cepheids in the hipparcos catalog, and then are used to show that the method of “reduced parallaxes” introduces a bias which is of the order of 0.01 mag or less. Based on these results we discuss in Sect. 8 the distance to the LMC, and elaborate on the various uncertainties.
## 2 Sample selection
The number and some properties of the Cepheid population in the hipparcos catalog were discussed by Groenewegen (1999, hereafter G99). To summarise: by cross-correlating the general hipparcos database, the hipparcos “resolved variable catalog” and the electronic database of Fernie et al. (1995; hereafter F95), a total of 280 Cepheids was identified. Then, 9 Type ii Cepheids, 1 factual RR Lyrae variable, 1 CH-like carbon-rich Cepheid, 10 double-mode Cepheids, 7 Cepheids with an unreliable hipparcos solution and 4 Cepheids with no or unreliable optical photometry were excluded. Note that for RY Sco and Y Lac we use the new determinations for the parallax and its error from Falin & Mignard (1999). This leaves 248 stars, of which 32 are classified as overtone pulsators by F95, Antonello et al. (1990) or Sachkov (1997). Luri et al. (1999) have classified Cepheids in fundamental and overtone pulsators using the hipparcos lightcurves, but did not yet publish the results for individual stars. For the sample, G99 calculated intensity-mean $`I`$ (on the Cousins system) and $`VI`$ magnitudes for 189 stars, and collected magnitude-mean colours for additional 14 stars, and provided $`JHK`$ intensity-mean magnitudes on the Carter system for 69 stars.
By comparison, the FC dataset consisted of 223 stars of which 3 were discarded. These were DP Vel for lack of photometric data, and AW Per and AX Cir because they are in binaries and the photometry might be affected by the companions. However, in the F95 catalog there are many more stars which are flagged for this reason. From the sample in G99 are therefore excluded the stars flagged “O C” (RX Cam, AW Per, T Mon, SS CMa, S Mus, AX Cir, W Sgr, V350 Sgr, U Aql, SU Cyg, V1334 Cyg) and “O: C” (VY Car) in F95. This leaves 32 overtone pulsators and 204 fundamental mode pulsators in the sample considered here. They are listed for completeness in Appendix A together with some adopted parameters.
Recently, L99 also studied the Cepheids in the hipparcos catalog. They selected stars listed as “DCEP” or “DCEPS” from the hipparcos catalog. Interestingly, they state that they selected 247 stars, while in reality there are 250 such stars (G99). After removing 9 Cepheids for which there is no photometry listed in F95, their final sample consists of of 238 stars (including 31 overtones). They did not eliminate double-mode Cepheids, Cepheids where the photometry is (likely) contaminated by a binary companion or Cepheids with unreliable hipparcos solutions. Furthermore, they assumed all Cepheids classified as “DCEPS” to be overtone pulsators, which is not the case. Their sample therefore includes stars they consider overtones which we and FC consider fundamental mode pulsators (e.g., SZ Cas, Y Oph, V496 Aql, V924 Cyg, V532 Cyg), and stars they consider fundamental mode pulsators which are overtone pulsators (e.g., V465 Mon, DK Vel, V950 Sco, see Mantegazza & Poretti 1992). In addition, V473 Lyr is considered by them to be a first overtone pulsator, while it probably is a second overtone pulsator (Van Hoolst & Waelkens 1995; Andrievsky et al. 1998; also see below).
Of the 236 stars in our sample there are 198 in common with the sample of FC. In other words, 22 stars in the FC sample would not have made it through the selection process outlined in G99 and here. Specifically, their sample contains 7 stars with unreliable hipparcos solutions (we use the improved parallax values for RY Sco and Y Lac from Falin & Mignard (1999), information which was not available to FC), 9 binaries (in addition to AW Per and AX Cir) where the photometry may be contaminated by the companions and 8 double-mode Cepheids. Another difference is that 7 more stars than in FC are flagged as overtone pulsators following F95 and Sachkov (1997).
For the 198 stars in common we have compared the intensity-mean $`V`$ and $`BV`$. FC mention that they use F95, except when the data were in Laney & Stobie (1993). Our photometry was at first instance taken solely from F95, except for RW Cas (see discussion in G99). In $`V`$, the photometry is identical for 165 stars. For 18 stars $`V`$ differs by $`>0.01`$mag, for 8 by $`>0.02`$mag and for 5 by $`>0.03`$mag. The latter cases were inspected individually.
For RW Cas the difference in the sense “Fernie et al $``$ FC” is 0.101 mag. As discussed in G99, there may be a typographical error in F95.
For X Pup the difference is $`0.047`$. The F95 entry of $`V`$ = 8.460 is close to the value derived in G99 for the dataset of Moffett & Barnes (1984). The other dataset considered in G99 gives a value of $`V=8.536`$. The value by FC of $`V=8.507`$ is intermediate. We have kept the F95 value.
For AQ Pup the difference is 0.122 mag. G99 calculated $`V,I`$ for 3 data sets. The entry in F95 ($`V=8.791`$) is identical to the Moffett & Barnes (1984) dataset. The FC value of $`V=8.669`$ is close to the 2 other datasets (8.676, 8.686 mag), and one of these was the preferred one in G99 regarding the $`V,I`$ photometry. For AQ Pup we use the $`V`$ and $`BV`$ from FC.
For RS Pup the difference is $`0.081`$. The value in F95 ($`V`$ = 6.947) is clearly off from both values in G99 (6.999 and 7.020 mag) which both are in agreement with the value of $`V`$ = 7.028 used by FC. For RS Pup we use $`V`$ and $`BV`$ from FC.
For S Nor the difference is $`0.032`$. The value used by FC is extremely close to the value derived in G99, and for S Nor we use the $`V`$ and $`BV`$ from FC.
After these changes in $`V`$ and $`BV`$, there are 164 with identical values for $`BV`$, for 12 stars $`BV`$ differs by $`>0.014`$ mag, for 9 by $`>0.028`$ mag, and for 5 by $`>0.042`$ mag (RW Cas, X Pup, U Nor, RY Sco, RU Sct). These cases have not been considered separately. It merely indicates that errors in $`V`$ and $`BV`$ (and hence reddening) can contribute to the uncertainties in the derivation of the zero point of the $`PL`$-relation, as will be discussed below.
For V1162 Aql, we discovered an error in the $`(BV)`$ value listed in F95. From the data in Berdnikov & Turner (1995) that was used in G99 to calculate $`(VI)`$ we derive and adopt an intensity weighted mean of $`(BV)`$ = 0.879 magnitudes.
In G99, $`V`$ and $`I`$ photometry was presented for many Cepheids in the hipparcos catalog, based on a literature search. When possible, intensity-mean magnitudes were calculated based on the original published datasets. Some magnitude-means were also presented. The $`VI`$ magnitudes are taken from G99. When there are multiple entries the first one was taken following the considerations in G99. In total there is $`I`$-band data for 191 stars (or 81% of the sample), 178 of which are intensity-mean magnitudes. In G99 it was shown that there is no significant difference between the intensity-mean and magnitude-mean, and so the remaining 13 magnitude-means are used without correction.
G99 also presented intensity-mean $`JHK`$ magnitudes in, or transformed to, the Carter system. When there are multiple entries the first one is taken, following the considerations in G99. In total there is $`JHK`$-band data for 63 stars, or 27% of the sample.
## 3 Analysing parallax data
Analysing parallax data is not a trivial exercise, and has led to some confusion in the literature. Part of this confusion is probably related to the fact that parallax data suffer from many types of biases (see Brown et al. 1997). For example, the most conspicuous bias, the Lutz-Kelker (LK; Lutz & Kelker 1973) bias, although known for a long time, could not be empirically investigated due to a lack of data on which it could be tested.
One can visualize the LK effect by analogy with the well-known Malmquist bias. Malmquist bias occurs on samples of objects because due to a particular magnitude cut, objects that are brighter than the mean will be included in a sample, while objects that are fainter than the mean will be excluded. The net effect of this type of bias is that the intrinsic magnitude for a certain sample will be too bright if no corrections are applied. To investigate the presence of such biases, one can look at so-called Spaenhauer diagrams, where the observed magnitude is plotted as a function of, for example, distance. In such a way it is relatively easy to find the point where the Malmquist bias starts to dominate (see e.g. Sandage 1994).
A similar, but opposite effect, is due to the LK bias. For a given parallax cut (or any selection based on the parallax, or its associated error), on average, more objects that are located far away will scatter into the sample, than stars will scatter out of the sample, as the sampled volume is larger for the former objects. Because the distances to the objects scattered into the sample are then underestimated, this effectively results in a too faint mean ‘intrinsic’ magnitude of the entire sample. As with the Malmquist bias, the LK-bias depends on the space distribution of the sample of stars under consideration. Contrary also to the Malmquist bias is that although the error-bars on the parallax are symmetric, those on the derived distances and intrinsic magnitudes are a-symmetric, which seriously affects the analysis of the data. In addition, the relative error-bars on the parallax increase with distance.
The LK bias was shown to exist empirically by Oudmaijer et al. (1998). Figure 1 gives an illustrative plot, showing the magnitudes of the Cepheids under consideration, derived from their photometry and observed parallax, as a function of parallax in the upper panel, and as function of relative error in the lower panel ($`\rho =V_0+5\mathrm{log}\pi +5\delta \mathrm{log}P`$, for $`\delta =2.81`$ – see next Section for details). As the absolute error in the hipparcos data is more or less constant, the two figures are equivalent. The inferred magnitude of the objects in the Cepheid sample is not a random distribution around the mean, but shows a clear trend as function of parallax. For large parallaxes (and small relative errors) the inferred magnitudes are too faint, and then, for larger distances, the objects become too bright. It was only in Oudmaijer et al. (1998), that this was shown empirically for the first time.
Indeed, it is clear from this type of graph, which is a version of a Spaenhauer diagram, that a selection on parallax (or relative error in parallax) will not return the correct answer. In most cases the result will be biased (but see Oudmaijer et al. 1999, for a counter-example).
Lutz & Kelker (1973) were the first to investigate this effect, and found that the resulting bias can be substantial, increasing with increasing relative error. For a relative error of 17.5% in the parallax, the bias is 0.43 magnitudes – yet as Koen (1992) later showed, the confidence intervals around these corrections are large. The 17.5% limit was taken literally by subsequent authors, and often stars with worse determinations were deleted or individual stars were corrected with the values LK derived.
LK claimed that individual stars suffer from the same bias, a point they emphasize strongly. Actually, it is quite a surprising result that LK concluded that individual stars are biased. LK start their calculations assuming that the observed parallax equals the true parallax with a random error, a standard procedure when investigating measurement errors and the main assumption for Monte-Carlo simulations as performed in e.g. Sect. 7 of this paper. It seems rather contradictory then that LK infer from their results that individual stars are biased if their initial conditions assume otherwise.
Let us first comment on whether the conjecture that the observed parallax is equal to the true parallax (plus a random observational error) is correct. It is actually quite hard to investigate this in a ‘bias-free’ manner, but the evidence provided by the hipparcos mission swings the balance towards answering this question with a ‘yes’. Measurements of stars beyond the detection limit of hipparcos, such as those in the Magellanic Clouds (Van Leeuwen 1997, Arenou et al. 1995) consistently give average observed parallaxes close to zero mas with a scatter of order the observational error. This indicates that the measurements are not biased, even though the relative errors on the parallax are substantial – and the results for individual stars can be widely off the true distance, but crucially, in the mean they appear to give the right answer within the errors rather than show a conspicuous bias in one direction. In addition, analyses of Open Clusters, for which photometric distances are known, do not show a deviation from the mean distance for increasing distances and thus lower quality data (Arenou et al. 1995, Robichon et al. 1999). Based on this, it appears that individual stellar determinations are most likely not biased.
There is one exception, mentioned briefly by Koen & Laney (1998). These authors argue that if a single star is the result of a selection on parallax, its measurement is biased. Although this sounds counter-intuitive, it is a formal implication from the fact that a selection on parallax gives rise to a bias.
A final comment concerns the use of parallaxes of a sample of which one does not know that the objects should have the same intrinsic magnitude, while their distances are not known. The trend observed in Fig. 1 that, for a given sample, large parallaxes return in principle too faint intrinsic magnitudes, while small parallaxes return too bright magnitudes, can have serious implications in the interpretation of the data.
### 3.1 How to deal with the Lutz-Kelker bias?
The remaining question is how one should deal with the effects of the bias. Should one, as often the case with Malmquist bias, introduce a (sample-dependent) correction factor, or are there ways to circumvent the problem?
In a paper dealing with this problem, Turon Lacarrieu & Crézé (1977) discuss two different methods: <sup>1</sup><sup>1</sup>1Note that Turon Lacarrieu & Crézé (1977) repeat the phrase by LK that ‘individual stars’ suffer from the bias, without specifying whether this would mean all stars (which contradicts their own assumption that the observed parallax is equal to the true parallax with an associated error-bar), or individual stars that are the result of a selection on parallax.
Accepting that a sample selected on parallax is inevitably biased, they first considered using the better parallaxes, and investigate the resulting mean magnitudes, and provide corrections along similar lines as LK (also see Smith 1987a+b+c, Koen 1992).
Secondly they considered a full sample, and, avoiding transformations from parallaxes to distances and magnitudes, they derive the mean parallax first and use the resulting mean parallax to derive the mean magnitude of the sample. Since negative parallaxes can not be incorporated into a mean magnitude, one in principle has to discard these data, in effect selecting on parallax and thus biasing the sample. Therefore, Turon Lacarrieu & Crézé (1977) introduced the so-called “reduced parallax” (10$`{}_{}{}^{0.2M_\mathrm{V}}\pi `$), which can take into account negative parallaxes, and hence the (weighted) mean of the reduced parallax can be converted into a mean magnitude. This method, as will be discussed in Sect. 3.2, indeed appears to be “bias-free”, mainly because a weighting scheme puts less weight on the larger deviations around the mean from the lower quality data, whilst not being a formal selection on parallax.
The Cepheids in hipparcos were investigated previously by FC who used the second, reduced parallax method, and Oudmaijer et al. (1998), who used a scheme based on the first method. Their results were equal within the error-bars. Although FC used the reduced parallax method, which formally does not suffer from LK bias (Koen & Laney 1998), they still suggested in a footnote that a LK correction of 0.02 mag should be applied on their final result. This is not necessary; the result only has to be corrected for Malmquist bias, as FC also pointed out, but did not actually apply, as they estimated it would essentially counteract their proposed LK correction of 0.02 mag.
These, and other issues will be investigated in the remainder of this paper. First, we will outline the method for the reduced parallaxes again.
### 3.2 The “reduced parallax” method
The method of “reduced parallax” (discussed by Turon Lacarrieu & Crézé 1977) was the one used by FC in analysing Cepheid data.
Consider a Period-Luminosity relation of the form:
$$M_\mathrm{V}=\delta \mathrm{log}P+\rho ,$$
(1)
where $`P`$ is the fundamental period in days. If $`V`$ is the intensity-mean visual magnitude and $`V_0`$ its reddening corrected value, then one can write:
$$10^{0.2\rho }=\pi \times \mathrm{0.01\hspace{0.17em}\hspace{0.33em}10}^{0.2(V_0\delta \mathrm{log}P)}\pi \times \mathrm{RHS},$$
(2)
which defines the expression rhs and where $`\pi `$ is the parallax in milli-arcseconds. This method has the advantage that negative parallaxes can be used in the analyses as well. A weighted-mean, with error, of the quantity 10<sup>0.2ρ</sup> is calculated, with the weight (weight = $`\frac{1}{\sigma ^2}`$) for the individual stars derived from:
$$\sigma ^2=\left(\sigma _\pi \times \mathrm{RHS}\right)^2+\left(0.2\mathrm{ln}(10)\pi \sigma _\mathrm{H}\times \mathrm{RHS}\right)^2,$$
(3)
with $`\sigma _\pi `$ the standard error in the parallax. This follows from the propagation-of-errors in Eq. (2). For the error (denoted $`\sigma _\mathrm{H}`$) in $`(V_0\delta \mathrm{log}P)`$ we follow FC’s “solution B” and adopt $`\sigma _\mathrm{H}=0.1`$ throughout this paper. Recently, L99 considered alternative weighting schemes, but concluded that the one used by FC and the present paper gives the most reliable zero point and the lowest dispersion.
The reddening is derived as follows. The intrinsic colours follow from the relation in Laney & Stobie (1994):
$$B_0V_0=0.416\mathrm{log}P+0.314,$$
(4)
which has a dispersion of 0.091 mag. The visual extinction ($`A_\mathrm{V}=R_\mathrm{V}\times E(BV)`$) is calculated using (Laney & Stobie 1993):
$$R_\mathrm{V}=3.07+0.28(BV)_0+0.04E(BV).$$
(5)
No dispersion is given for this relation, only an error of 0.03 in the zero point. We will assume that the dispersion in Eq. (5) is slightly larger than this, namely 0.05 mag.
For overtone pulsators, the fundamental period has to be estimated from the observed period. This was done, following FC, using:
$$P_1/P_0=0.7160.027\mathrm{log}P_1,$$
(6)
with a dispersion we estimate from the original data (Alcock et al. 1995) to be of order 0.002.
For pulsators in the second overtone, the fundamental period is calculated, following FC, using:
$$P_2/P_0=0.55.$$
(7)
This completes the description of the method used by FC.
Alternative methods which are described now, follow this description closely but are based on $`V`$ and $`I`$, respectively $`J`$ and $`K`$ photometry instead of $`B`$ and $`V`$. The rationale being that the extinction in $`I`$ and $`K`$ is less than in $`V`$, and that the scatter in the $`M_\mathrm{I}P`$\- and $`M_\mathrm{K}P`$-relations is less than in the $`M_\mathrm{V}P`$-relation (Tanvir 1999, Laney & Stobie 1994).
First consider the case based on $`V`$ and $`I`$. The intrinsic colour is derived from (Caldwell & Coulson 1986):
$$V_0I_0=0.292\mathrm{log}P+0.443$$
(8)
which has a dispersion of 0.064 mag. This relation was derived for magnitude-mean $`(VI)_0`$ but we will assume it to hold for intensity-mean magnitudes as well. In a recent paper, Feast (1999, Appendix D), presents a correction formula that implies that the difference intensity-mean minus magnitude-mean $`(VI)`$ colour is of order $`0.017`$ mag for a typical $`V`$-band amplitude of 0.7 mag.
There seems not to exist a relation similar to Eq. (5), where $`R(I)A_\mathrm{I}/E(VI)`$ is related to $`(VI)_0`$ and/or $`E(VI)`$. We have derived such a relation from the available $`BVI`$ data.
For each star with $`BVI`$ photometry, $`(BV)_0`$ and $`(VI)_0`$ can be calculated from Eqs. (4) and (8). Then, $`A_\mathrm{I}`$ is calculated using Gieren et al. (1998):
$$R_\mathrm{I}=1.82+0.205(BV)_0+0.022E(BV),$$
(9)
and $`A_\mathrm{I}=R_\mathrm{I}\times E(BV)`$. In Fig. 2 $`R(I)`$ is plotted versus $`(VI)_0`$. A least-square fit to 183 data points gives:
$$R(I)=1.422$$
(10)
with no significant dependence on $`(VI)_0`$ and with a standard error of 0.19. For $`\sigma _\mathrm{H}`$ (in this case the error in $`(I_0\delta _\mathrm{I}\mathrm{log}P`$)) we adopt a value of 0.15 mag (Tanvir 1999).
A variation on this method that treats the problem of reddening in a different way, is to use the reddening-free so-called “Wesenheit-index” (see for example Tanvir 1999), that uses the observed colours but is essentially reddening-free when defined as:
$$W=V2.42(VI).$$
(11)
For $`\sigma _\mathrm{H}`$ (in this case the error in $`(W\delta _\mathrm{W}\mathrm{log}P`$)) we adopt a value of 0.11 mag (Tanvir 1999).
Now consider the case based on $`J`$ and $`K`$ colours. The intrinsic color is derived from Laney & Stobie (1994):
$$J_0K_0=0.149\mathrm{log}P+0.310$$
(12)
which has a dispersion of 0.044 mag. Again, there seems not to exist a relation similar to Eq. (5), where $`R(K)A_\mathrm{K}/E(JK)`$ is related to $`(JK)_0`$ and/or $`E(JK)`$. We have derived such a relation from the available $`BVJK`$ data.
For each star with $`BVJK`$ photometry, $`(BV)_0`$ and $`(JK)_0`$ can be calculated from Eqs. (4) and (12). Then, $`A_\mathrm{K}`$ is calculated using Laney & Stobie (1993):
$$A_\mathrm{K}=0.279E(BV).$$
(13)
In Fig. 3 $`R(K)`$ is plotted versus $`(JK)_0`$. A least-square fit to 55 data points gives:
$$R(K)=1.0351.063(JK)_0$$
(14)
with a standard error of 0.091. For $`\sigma _\mathrm{H}`$ (in this case the error in $`(K_0\delta _\mathrm{K}\mathrm{log}P`$)) we adopt a value of 0.12 mag (Gieren et al. 1998).
## 4 Zero points of the $`PL`$-relations
### 4.1 Zero point of the $`M_\mathrm{V}\mathrm{log}P`$-relation
In Tables 1-4 we present the results from the “reduced parallax” method presented in the previous section for different samples of stars, and based on different colors ($`BV`$, $`VI`$ or $`JK`$). To test the implementation of the method we have calculated the zero point for some of the solutions considered in FC, adopting a slope $`\delta =2.81`$ and working with $`BV`$ colours, as they did. The data for the sample used by FC come from Feast & Whitelock (1997; their Table 1). These are our solutions 1-7 in Table 1, and the value for the zero point and total weight are in perfect agreement to within the listed number of decimal places in FC. The error in the zero point determination is slightly different, ours being larger by a few 1/100-th of a magnitude. Using the FC sample, we also calculated the solution for the overtones only, and the overtones excluding Polaris (solutions 8-9).
For the present sample and using a slope $`\delta =2.81`$ solutions 10-12 give the value for the zero point for the whole sample, and for fundamental mode and overtone pulsators separately. In this case, V473 Lyr was considered to be pulsating in the second overtone. This interesting object is thought to be the only Galactic Cepheid pulsating in the second overtone (Van Hoolst & Waelkens 1995 and Andrievsky et al. 1998). The values of $`\rho `$ assuming that V473 Lyr is a fundamental mode, first overtone or second overtone pulsator, respectively, are $`2.68`$, $`2.07`$ and $`1.61`$, all with an error of 0.70. With such a large error only fundamental mode pulsation can be excluded from the hipparcos parallax alone. The zero point assuming second overtone pulsation is consistent with the zero point obtained for the whole sample, and we will hence assume in our zero point determinations that V473 Lyr is indeed a second overtone pulsator.
The zero point for the entire sample of $`1.41\pm 0.10`$ compares to the value of FC of $`1.43\pm 0.10`$, and the value of L99 of $`1.44\pm 0.05`$ (they used a slightly different slope of $`2.77`$). These values are all very similar, and the differences are mainly due to the different samples of Cepheids, and to a lesser extent to the slight differences in the adopted photometry. We believe that the present sample is the most carefully selected sample of the three with respect to completeness, the flagging of overtone pulsators, and the removal of Cepheids that may influence the analysis for various reasons (double-mode Cepheids, unreliable hipparcos solutions, possible contaminated photometry due to binary companions), as explained in the introduction.
#### 4.1.1 $`P_0`$ versus $`P_1`$
A first comment is on the different solution from the fundamental mode and overtone pulsators. In FC the difference between the solution using the fundamental pulsators or the full sample was only 0.02 mag. In our case it is about 0.08. FC did not present the solution for the overtones only, but we have calculated it from their data (solution 8). The difference using only the overtones or the fundamental pulsators is 0.05 mag in their case. In our case it is 0.15 mag, which is a difference at the 1$`\sigma `$ level. This indicates how important it is to carefully flag the overtone pulsators.
#### 4.1.2 Selecting on visual magnitude
L99 derive a zero point of $`1.44`$ (after correcting for a bias of $`0.01`$ mag) with a very small error of 0.05, using a selection on $`V5.5`$. We confirm (solutions 14 and 15) that the derived zero point is not significantly different from that for the whole sample, but we do not confirm such a small error. In fact, Pont (1999) argues that the error of 0.10 in FC is even underestimated based on his numerical simulations. Our simulations confirm this (Sect. 7) and so it is not clear how L99 arrived at such a small error.
#### 4.1.3 Selecting on parallax
For solutions 16-18 only positive parallaxes have been selected to highlight the effect of LK-bias. The zero points are fainter, as expected. Interestingly enough, the zero point for the overtone pulsators is hardly changed. Only 5 overtones (15%) have negative observed parallaxes and those carry very little weight, while for the fundamental mode pulsators a selection on positive parallaxes reduces the number by 32%.
#### 4.1.4 Selecting on weight
FCs final choice for the zero point relied heavily on a subsample of 26 stars, selected on weight. The question arises if this introduces some bias. For solutions 19-24 we have made different selections on weight, in particular solutions 19-22 have been devised such that each bin selected on individual weight has about the same total weight, so that the errors on the zero point are comparable. The differences are at the 1$`\sigma `$ level. From numerical simulations performed in Sect. 7 (see Table 6) we confirm that there are no indications at the present level of accuracy that a selection on weight is an (indirect) selection on parallax, and so a sample selected on weight appears not to be subject to LK-bias. The error on the zero point determination is larger when selecting on (individual) weight, simply because of the smaller value of the total weight (see Table 6). This is different from the result in FC, who quote a smaller error on the zero point for the sample selected on weight compared to their full sample (cf. their solution 6 and 1).
#### 4.1.5 Selecting on period
Solutions 25-36 represent cases where the sample was split up in bins in $`\mathrm{log}P`$ carrying approximately equal weight, for the whole sample (solutions 25-28), and the fundamental pulsators only (solutions 29-32, and 33-36 for a slope of $`2.22`$). In both cases, the star with the highest individual weight was excluded in defining the bins. There is a significant dependence of the zero point on $`\mathrm{log}P`$. This is particularly clear in the case of the fundamental mode pulsators where the zero point for stars with $`\mathrm{log}P0.85`$ differs at the 3$`\sigma `$ level from the shortest period bin.
To expand on this matter further, we show in Fig. 4 how 10<sup>0.2ρ</sup> depends on $`\mathrm{log}P`$ for the full sample (top panel), and for the 47 stars with an individual weight $`>5`$. Shown are the weighted-mean values of 10<sup>0.2ρ</sup> (solid lines), and weighted least-square fits to the data of the form 10$`{}_{}{}^{0.2\rho }=\alpha \mathrm{log}P+\beta `$ (dashed line). The ($`\alpha ,\beta `$) found are (0.11 $`\pm `$ 0.12, 0.44 $`\pm `$ 0.10) for the full sample, and (0.094 $`\pm `$ 0.115, 0.43 $`\pm `$ 0.09) for the 47 stars. The analysis was repeated for the fundamental mode pulsators only, giving samples of 204 and 35 stars, respectively. The ($`\alpha ,\beta `$) found are (0.17 $`\pm `$ 0.15, 0.35 $`\pm `$ 0.14) for the full sample, and (0.17 $`\pm `$ 0.15, 0.34 $`\pm `$ 0.14) for the stars with weight $`>5`$. The slopes derived are significant at the 1$`\sigma `$ level only.
These results depend on the adopted slope in the $`PL`$-relation. For a slope $`\delta =3.1`$ we find that the effect of the dependence of the zero point on the binning in $`\mathrm{log}P`$ is enhanced, and that the slope in the 10<sup>0.2ρ</sup> versus $`\mathrm{log}P`$ relation becomes significant at the 2$`\sigma `$ level. We also calculated the results for a slope of $`\delta =2.22`$ (Bono et al. 1999). The results are listed as solutions 33-36 for the fundamental pulsators. The dependence on period is not significant in this case, although the shortest bin remains the brightest. Least-square fits give the following results for ($`\alpha ,\beta `$): all stars ($`0.039\pm `$ 0.090, 0.45 $`\pm `$ 0.08), 56 stars with weight $`>5`$ ($`0.050\pm `$ 0.099, 0.45 $`\pm `$ 0.08), all fundamental pulsators (0.00 $`\pm `$ 0.11, 0.39 $`\pm `$ 0.11), 42 fundamental pulsators with weight $`>5`$ (0.00 $`\pm `$ 0.13, 0.38 $`\pm `$ 0.12). The slope derived is no longer significant. From this analysis one may conclude that there is weak evidence for the fact that the slope in the $`M_\mathrm{V}\mathrm{log}P`$-relation is shallower than in the LMC, or, alternatively that there is a change of slope at the short period end. A similar effect is found for the $`M_\mathrm{I}\mathrm{log}P`$-relation (see next section). The implications are discussed further in Sect. 8.
### 4.2 Zero point of the $`M_\mathrm{I}\mathrm{log}P`$-relation
In Table 2 the zero points of the $`M_\mathrm{I}\mathrm{log}P`$-relation are listed based on $`V,I`$ photometry. Unless otherwise noted, a slope of $`3.05`$ is used. This is the slope adopted in L99, and is based on work of Tanvir (1999) and Gieren et al. (1998). In every case we have calculated both the zero points for the $`M_\mathrm{I}\mathrm{log}P`$\- and $`M_\mathrm{V}\mathrm{log}P`$-relation for the respective samples. For the whole sample we find $`\rho =1.89\pm 0.11`$. This is 0.08 mag brighter than in L99, but within their and our quoted errors. L99 used the magnitude-mean magnitudes in Caldwell & Coulson (1987), and then applied a correction of $`0.03`$ mag to convert these to intensity-mean magnitudes. We use for most stars the intensity-mean magnitudes as calculated in G99 from the original data.
As for the solutions based on $`BV`$ photometry we find that the zero point using only the fundamental modes is brighter than using only the overtone pulsators, but the difference is now less than 1$`\sigma `$. Polaris again is the star with the highest individual weight. Solutions 6-8 illustrate the effect of LK-bias when selecting stars with positive parallaxes.
#### 4.2.1 Selecting on period
As before, we have split up the sample according to period in bins of approximate equal total weight, for all stars (solutions 9-12, excluding Polaris), and for the fundamental pulsators (solutions 13-16, excluding $`\delta `$ Cep, the fundamental pulsator with the highest individual weight). Again we see that the zero point depends on the period bin chosen, but the effect is not as systematic as in the $`V`$-band. In Fig. 5 10<sup>0.2ρ</sup> is plotted against $`\mathrm{log}P`$. We have made least-square fits as before and find for the whole sample ($`\alpha ,\beta `$) = (0.11 $`\pm `$ 0.11, 0.331 $`\pm `$ 0.091), for the 48 stars with an individual weight $`>5`$ ($`\alpha ,\beta `$) = (0.101 $`\pm `$ 0.098, 0.332 $`\pm `$ 0.079), for all 163 fundamental mode pulsators ($`\alpha ,\beta `$) = (0.17 $`\pm `$ 0.14, 0.26 $`\pm `$ 0.13), and for the 35 fundamental mode pulsators with a weight $`>5`$ ($`\alpha ,\beta `$) = (0.17 $`\pm `$ 0.12, 0.25 $`\pm `$ 0.11). The slope is again significant at at most the 1$`\sigma `$ level.
Solutions 17-21 give the division in period bins for a slope of the $`PL`$-relation of $`2.35`$ for the fundamental mode pulsators (the slope in the $`PL`$-relation in $`V`$ is $`2.22`$). After taking into account the difference in zero points in $`V`$, the dependence of the zero point in $`I`$ on period bin is not significant. The least-square fits give the following results for ($`\alpha ,\beta `$): the whole sample ($`0.024\pm `$ 0.076, 0.343 $`\pm `$ 0.065), for the 62 stars with an individual weight $`>5`$ ($`0.013`$ $`\pm `$ 0.077, 0.317 $`\pm `$ 0.069), for all 163 fundamental mode pulsators (0.012 $`\pm `$ 0.093, 0.294 $`\pm `$ 0.089), and for the 50 fundamental mode pulsators with a weight $`>5`$ (0.032 $`\pm `$ 0.093, 0.264 $`\pm `$ 0.090). As we found previously, a shallower slope gives rise to a smaller or no dependence of the zero point on period.
### 4.3 Zero point of the $`M_\mathrm{K}\mathrm{log}P`$-relation
Table 3 lists the results for the zero point of the $`M_\mathrm{K}\mathrm{log}P`$-relation and the zero point for the $`PL`$-relations in $`V`$ and $`I`$ for the same samples. Unless noted otherwise we have used a slope of $`\delta =3.27`$ from Gieren et al. (1998). For the whole sample we find a zero point of $`2.61\pm 0.17`$, but note that the corresponding zero points in the $`V`$ and $`I`$ band of this sub-sample are too bright by about 0.11 magnitude compared to the full samples in $`V`$ and $`I`$, and so the true zero point is probably closer to $`2.50`$. These off-sets indicate biases due to the small number of available measurements in the NIR. This is especially clear when further sub-divisions into smaller samples are made, like a division in period (solutions 9-13). Especially in one bin (solution 10) the result depends very much on one star which was taken out (solution 11). After shifting the zero points in $`K`$ by an amount such as to make the zero points in $`V`$ all equal there is no evidence for a dependence of the zero point on period.
### 4.4 Zero point of the $`M_\mathrm{W}\mathrm{log}P`$-relation
Table 4 lists the results for the zero point of the $`M_\mathrm{W}\mathrm{log}P`$-relation. We have used a slope of $`\delta =3.411`$ (Tanvir 1999). As before the sample is divided into period bins. After shifting the zero points in $`W`$ by an amount such as to make the zero points in $`V`$ all equal (this value can be taken from the corresponding solutions in Table 2) there is no evidence for a dependence of the zero point on period.
## 5 Summary of the results
### 5.1 On the use of longer wavelengths
One argument to consider $`PL`$-relations in $`I`$ and $`K`$ besides the traditional relation in $`V`$, is because of the smaller extinction at longer wavelengths. We have calculated the zero points for the whole sample considering a systematic shift for all stars of 0.005 mag in $`V`$ (respectively $`I`$ and $`K`$), 0.007 in $`BV`$ (resp. $`VI`$, $`JK`$), a shift in the period-colour relations (Eqs. (4), (8), (12)) equal to 1/10-th of the quoted dispersions, and a shift in the selective reddening (Eqs. (5), (10), (14)) equal to the quoted dispersion. The results are listed in Table 1 (solutions 48-51), Table 2 (solutions 23-26), and Table 3 (solutions 15-18). Adding the differences between these zero points and that for the default case in quadrature, one arrives at estimated uncertainties due to errors in the photometry, reddening and period-colour relations of 0.038 mag in $`V`$, 0.032 mag in $`I`$, and 0.011 mag in $`K`$. This illustrates that the $`K`$-band $`PL`$-relation is indeed the least sensitive to these effects. For the Wesenheit-index the error due to the photometry and reddening coefficients amounts to 0.028 mag (solutions 13-15 in Table 4).
Unfortunately, the advantages of using the infrared, like the intrinsically tighter $`PL`$-relation and the insensitivity to reddening, are countered by the fact that so few stars have been measured in the NIR so far. Of the approximately 48 stars with in 1 kpc, $`BV`$ photometry is available for all of them, $`VI`$ for 41 of them, but $`JHK`$ data (of sufficient quality) for only 24. It is estimated that determining the intensity-mean NIR magnitudes of the remaining 24 stars alone would bring the scatter in the zero point determination in the $`K`$-band down from about 0.17 to less than 0.1 mag.
### 5.2 On the slopes of Galactic $`PL`$-relations
Another important issue concerns the slopes in the respective Galactic $`PL`$-relations. Common practice is to adopt the slope determined for Cepheids in the LMC, but the slope could be different for Galactic Cepheids. In Table 5 we have collected slopes of the $`PL`$-relations from the literature for Cepheids in the Galaxy, LMC and SMC, in $`V,I,K`$, both observationally determined and from two recent theoretical papers (Alibert et al. 1999, Bono et al. 1999 and private communication). Table 5 includes ths slopes in $`V`$ and $`I`$ by Madore & Freedman (1991) used by the hst $`H_0`$ Key Project (see for example Gibson et al. 1999). From the results of Gieren et al. (1998) on the Cepheids in the LMC, we have calculated the error in the slope using the data they kindly provided (their Tables 8 and 9). In addition we have calculated the slope and zero point for their data set but using cut-offs in period, for reasons that will be explained later.
There are several interesting features to be noted about the slopes. Observationally, the slopes in the $`V`$ and $`I`$ $`PL`$-relations in the LMC are very well established, respectively, about $`2.81`$ and $`3.05`$, and these are the default slopes adopted in the present study, with errors of about 0.08 and 0.06. For $`M_\mathrm{K}\mathrm{log}P`$-relation there are fewer observational data available but the slope in Gieren et al. (1998) is determined very accurately and is in reasonable agreement with the result of Madore & Freedman (1991).
One fact needs to be mentioned however, and that is that the period distribution of the calibrating Cepheids in the LMC is very different from that of the Galactic Cepheids in the hipparcos catalog. In Table 8 in Gieren et al. (1998) there are 53 LMC Cepheids listed with $`V,I`$ photometry that define their $`PL`$-relation. Nine have $`\mathrm{log}P0.555`$, then there is a gap, and 42 have $`\mathrm{log}P0.846`$. In $`JHK`$ (their Table 9), there are 59 Cepheids, including 5 that have $`\mathrm{log}P0.680`$, then there is a gap, and 54 have $`\mathrm{log}P0.834`$. The same dichotomy of Cepheids in period can be seen from Tanvir (1999). Note that in more distant Galaxies than the Magellanic Clouds, due to observational bias, most known Cepheids have periods longer than 10 days. We therefore have included in Table 5 the slopes in $`V,I,K`$ based on the data in Gieren et al. (1998), but have divided the sample according to period. The slopes for the long period sub-sample differ only slightly from that for the whole sample but are characterised by slightly larger errors. The slopes for the short-period sample have larger errors because of the small number of stars involved but nevertheless are systematically steeper at the 1$`\sigma `$ level in all three colours. For comparison, in our sample only 107 of the 236 Cepheids have $`\mathrm{log}P0.846`$, and the zero point of this sample differs at the 1.4$`\sigma `$ level from that of the whole sample (solution 52 in Table 1). This is yet another indication that for the default slope of $`2.81`$ the zero point may depend on the period range chosen. For a slope of $`\delta =2.22`$ (compare solutions 38, 53) this is not the case.
A second issue concerns the predictions of the theoretical models and the comparison with observations. In $`V,K`$ for the LMC, and in $`I`$ for the SMC the models of Bono et al. (1999) are in very good agreement with the observed slope, in $`I`$ for the LMC and $`V`$ for the SMC the agreement is poor. In $`I,K`$ for the LMC, the models of Alibert et al. (1999) are in good agreement with the observed slopes; in $`V`$ the agreement for the LMC and $`V,I`$ for the SMC is fair. However, the prediction both models make for the Galaxy are very different. Gieren et al. (1998) have derived $`PL`$-relations for Galactic Cepheids using the surface brightness technique. The slopes they find in all three bands are steeper than the corresponding slopes in the LMC at the 2-3 $`\sigma `$ level. In their paper they ascribe this to small number statistics, and in the end assume the LMC slopes to hold for the Galactic Cepheids as well, also, they add, because the slopes for the LMC Cepheids are better determined. The models of Alibert et al. (1999) predict slopes for the Galactic $`PL`$-relations that are steeper than the observed ones in the LMC in $`V,I,K`$ as well, although by a small amount only (the main conclusion of their paper is actually that the slope and zero point of the $`PL`$-relations do not depend on metallicity). By contrast, the Bono et al. (1999) paper predicts slopes that are significantly shallower in the Galaxy compared to the LMC especially in $`V,I`$ but also in $`K`$. In addition, the Bono et al. models actually predict a non-linear $`PL`$-relation in $`V`$, for all metallicities considered. Furthermore, Alibert et al. mention that a change of slope in $`PL`$-relations is expected at short periods due to the reduction of the blue loop during core He burning, and that this change of slope occurs near $`\mathrm{log}P=0.2`$ for $`Z`$ = 0.004, $`\mathrm{log}P=0.35`$ for $`Z`$ = 0.01, and $`\mathrm{log}P=0.5`$ for $`Z`$ = 0.02. Such a change of slope (in the sense that the slope of the $`PL`$-relation is steeper for periods below this limit) was recently observed for SMC Cepheids (Bauer et al. 1999) with the break occurring near $`\mathrm{log}P=0.3`$, near the predicted value.
### 5.3 Consistency between individual distances based on different colours
For a given slope and zero point of the $`PL`$-relation one can calculate the photometric distance, and since we have determined the zero point for three photometric band we can intercompare photometric distances to individual Cepheids. This is illustrated in Fig. 6, for the default slopes of the $`PL`$-relations. The zero point of the $`M_\mathrm{V}\mathrm{log}P`$ relation is fixed at $`1.411`$ (solution 10, Table 1). The zero point of the $`M_\mathrm{I}\mathrm{log}P`$ relation is determined to give a mean difference in $`(mM)_{0,\mathrm{I}}(mM)_{0,\mathrm{V}}`$ of zero, and is found to be $`1.918`$ (top panel Fig. 6). The rms dispersion is 0.14 mag. Similarly, to create the bottom panel, the zero point of the $`M_\mathrm{K}\mathrm{log}P`$ relation was determined to give a mean difference in $`(mM)_{0,\mathrm{K}}(mM)_{0,\mathrm{I}}`$ of zero, and is found to be $`2.600`$, with an rms dispersion is 0.10. The values for the zero points in $`I`$ and $`K`$ derived in this way differ only slightly from the solutions 1 in Tables 2 and 3.
Interestingly, least-square fitting shows that the slopes of the relations are not unity, but 0.980 $`\pm `$ 0.007 (top panel), and 0.965 $`\pm `$ 0.007 (bottom panel). Using the same procedure, but adopting steeper slopes of $`3.04`$, $`3.33`$, $`3.60`$ (see Sect. 8 for the reason of these choices) in, respectively, the $`M_\mathrm{V}\mathrm{log}P`$, $`M_\mathrm{I}\mathrm{log}P`$ and $`M_\mathrm{K}\mathrm{log}P`$ relations and fixing the zero point of the $`M_\mathrm{V}\mathrm{log}P`$ relation at $`1.234`$ we find in the same way the zero points of the $`M_\mathrm{I}\mathrm{log}P`$ and $`M_\mathrm{K}\mathrm{log}P`$-relations to be $`1.696`$ and $`2.328`$. The slopes are 0.983 $`\pm `$ 0.007 and 0.971 $`\pm `$ 0.007.
Using the same procedure, but adopting shallower slopes of $`2.22`$, $`2.35`$, $`3.05`$ in, respectively, the $`M_\mathrm{V}\mathrm{log}P`$, $`M_\mathrm{I}\mathrm{log}P`$ and $`M_\mathrm{K}\mathrm{log}P`$ relations and fixing the zero point of the $`M_\mathrm{V}\mathrm{log}P`$ relation at $`1.885`$ (solution 38 in Table 1) we find the zero point of the $`M_\mathrm{I}\mathrm{log}P`$ and $`M_\mathrm{K}\mathrm{log}P`$-relations to be $`2.491`$ and $`2.692`$. The slopes are 0.974 $`\pm `$ 0.007 and 1.010 $`\pm `$ 0.010.
The conclusion is the the photometric distances based on different $`PL`$-relations are consistent with each other at a level of 0.10-0.14 mag, similar to the uncertainties in the individually derived zero points in $`V,I,K`$ for the full sample. The data does not allow to discriminate between different choices of the slopes of the $`PL`$-relations.
## 6 A volume complete sample
The question remains how representative the hipparcos Cepheid sample is. As with all stars that made it in the hipparcos catalog, they were proposed by Principal Investigators in solicited proposals. If proposers preferred their ‘favorite’ objects to be observed, there may be a ‘Human’ bias, which is impossible to correct for. For example, at some point in the selection process of the hipparcos Input Catalog the number of Cepheids with magnitudes between 9.5 and 11 were ‘tuned-up’ (ESA 1989, page 80). Also, the distribution in pulsation period of the about 72 Cepheids that were proposed to be observed but did not end up in the hipparcos Input Catalog is somewhat different from the 270 that did make it (ESA 1989, page 96).
In view of this it may be instructive to try to construct a volume complete sample. The hipparcos Input Catalog (ESA 1989, page 94) mentions that all 55 Cepheids within 1 kpc have been included. This depends of course on the adopted slope and zero point to calculate the photometric parallax. The F95 database lists 51 stars within 1 kpc for their adopted relation of $`M_\mathrm{V}=2.902\mathrm{log}P1.203`$.
In Fig. 7 the cumulative (photometric) parallax distribution is plotted for the 236 stars in the sample, with distances from F95, or calculated in the same way for the stars not listed there. Also indicated is a line with slope $`2`$ as expected for a disk population (with $`N(r)drr`$ it follows that $`N(\pi )d\pi \pi ^3`$, and so the cumulative distribution is proportional to $`\pi ^2`$). From this figure it is confirmed that the hipparcos catalog is volume complete to approximately 1 kpc.
On the other hand, three new Cepheids were discovered with hipparcos. For the same slope and zero point as used in F95 the photometric parallaxes of CK Cam, V898 Cen, V411 Lac are 1.77, 0.73 and 1.17 mas, respectively.
As hipparcos discovered 2 Cepheids that are closer than 1 kpc, it implies that the present sample may now be indeed complete down to 1 mas, but is possibly only complete for stars with a parallax $`\stackrel{>}{}`$1.8 mas as no new Cepheids closer than this limit have been discovered.
In Tables 1-4 we include solutions for volume complete samples, or, in other words, complete in photometric parallax. Since this depends on the zero point itself, this is an iterative process. For the $`M_\mathrm{V}\mathrm{log}P`$-relation we give the results in Table 1 for a cut-off at 1 and 1.8 mas. In solutions 39, 41, 43 we list the zero point for an input zero point of $`1.411`$ (solution 10) for the whole sample, the fundamental pulsators, and the overtones, respectively. The next line (solutions 40, 42, 44) lists the final result after iteration on the zero point. Solutions 45-47 give the final results for a higher cut-off in photometric parallax. In most cases the zero points are slightly brighter than the corresponding solutions 10, 11, 12. This is a bit surprising because of the Malmquist bias one expects the volume complete sample to be dimmer. The expected Malmquist bias for a disk population with an intrinsic spread of $`\sigma _\mathrm{H}`$ = 0.10 around the mean is about 0.009 mag (Stobie et al. 1989). This is much smaller than the typical error estimate in the zero points, and so the fact that the volume complete sample is brighter probably reflects the slightly different nature of that sample. For example, the mean value of $`\mathrm{log}P_0`$ for the whole sample is 0.86, while that for the volume complete sample is 0.77. As we noticed earlier, since the shortest period bins give brighter zero points this could be an explanation.
For the $`M_\mathrm{I}\mathrm{log}P`$-relation the construction of a volume complete sample is slightly more complicated due to the fact that for some of the 47 stars with $`BV`$ photometry and a photometric parallax $`>1`$ mas no $`I`$-band photometry is available. This implies that one has to apply a more stringent criterion to obtain a volume complete sample of stars that also have $`I`$-band photometry. This turns out to be a photometric parallax limit of 1.8 mas, and solution 23 lists the result. The zero point is slightly brighter, and this again may be due to the fact that the average $`\mathrm{log}P_0`$ value is smaller (0.74) than for the whole sample (0.90).
In the $`K`$-band no meaningful volume complete sample can be constructed. As there is no NIR data available for Polaris, any volume complete sample could be constructed for fundamental pulsators only. In any case, even then, a volume complete sample of stars with NIR data would have a cut-off at 2.6 mas, and would only include 3 stars. For completeness, it has been included in Table 3 (sol. 14).
The conclusion is that the zero point of the $`PL`$-relation based on a volume complete sample are within $`1\sigma `$ of the results for the full sample. Surprisingly, the zero points are brighter, contrary to one would expect from Malmquist bias. However, the Malmquist correction is expected to be small (see the next section for an estimate based on numerical simulations), and certainly much smaller than the error in the zero point, so that this effect is due to the slightly different nature of the volume complete samples compared to the full samples.
## 7 Properties of the hipparcos Cepheids, and numerical simulations
In this section we describe a numerical model to first of all construct synthetic samples of stars that fit some observed properties of the hipparcos Cepheids. Second, this model is used to investigate the numerical bias involved in applying the “reduced parallax” method.
The observed distributions of the fundamental period, the $`V`$ magnitude, de-reddened $`(BV)_0`$, colour excess, and absolute distance to the Galactic plane are plotted in Fig. 8 for the 236 Cepheids in our sample (the dotted lines). For the overtones, the fundamental period was calculated following Eqs. (6-7). $`(BV)_0`$ follows from Eq. (4), and $`E(BV)`$ from the observed value of $`(BV)`$ minus $`(BV)_0`$. The absolute distance to the Galactic plane is calculated from the galactic latitude combined with the de-reddened $`V`$-magnitude and a photometric distance based on Eq. (1) with $`\delta =2.81`$ and $`\rho =1.43`$. The distributions are compared to simulations described now.
We have devised a numerical code to simulate the distributions described above. The input period distribution is (assumed ad-hoc to be) a Gaussian in $`\mathrm{log}P`$ with mean $`X_\mathrm{P}`$, and spread $`X_{\sigma \mathrm{P}}`$. Based on the observed properties of the sample only $`\mathrm{log}P`$ values between 0.43 and 1.66 are allowed. A random number is drawn and a $`\mathrm{log}P`$ selected following this distribution. The values of $`X_\mathrm{P}`$ and $`X_{\sigma \mathrm{P}}`$ directly influence the resulting distribution and so are easily determined. We find that values of $`X_\mathrm{P}`$ between 0.65 and 0.75 and of $`X_{\sigma \mathrm{P}}`$ between 0.2 and 0.25 give acceptable fits.
The galactic distribution of Cepheids is assumed to be a double exponential disk with a scale height $`H`$ in the $`z`$-direction (the coordinate perpendicular to the galactic plane), and a scale height $`R_{\mathrm{GC}}`$ in the galacto-centric direction. The coordinate system used is cylindrical coordinates centered on the Galactic Centre. Three random numbers are drawn to select the distance to the Galactic plane, the distance to the Galactic centre and a random angle $`\varphi `$ between 0 and 2$`\pi `$ in the Galactic plane centered on the Galactic centre. From this the distance $`d`$ to the Sun is calculated. Based on the observed photometric parallaxes of the sample only stars closer than 7800 pc to the Sun are allowed. We find no evidence for a gradient in the number of Cepheids with galacto-centric radius, in other words $`R_{\mathrm{GC}}=\mathrm{}`$ gives good fits, probably indicating that the volume sampled is too small to detect such a gradient, or that the distribution of Cepheids is at least equally determined by another factor, for example the location of the spiral arms. The value of $`H`$ is directly determined by the distribution of the number of stars as a function of $`z`$, and is found to be between 60 and 80 pc. This is consistent with the scale height of a relative massive population of stars, as the Cepheids are.
The value of $`M_\mathrm{V}`$ and $`(BV)_0`$ are correlated in the sense that brighter Cepheids are also bluer for a given period. Contrary to other studies we do not assume a Gaussian spread around the $`M_\mathrm{V}\mathrm{log}P`$ and $`(BV)_0\mathrm{log}P`$ relations, but instead a ‘box’ like distribution which is more physical because of the finite width of the instability strip. However, this does assume that the instability strip is uniformly filled (for all periods).
A single uniform random number, $`Rn`$, is drawn and then
$$M_\mathrm{V}=2.81\mathrm{log}P1.43+(0.42+0.84\times \mathrm{Rn})$$
(15)
and
$$(BV)_0=0.416\mathrm{log}P+0.314+(0.15+0.3\times \mathrm{Rn})$$
(16)
are calculated. The full width of the instability strip in $`M_\mathrm{V}`$ is taken to be 0.84 magnitude (derived from plots in Gieren et al. 1998, Tanvir 1999), and that of the $`(BV)_0`$ relation to be 0.3 mag. (Laney & Stobie 1994).
The reddening is calculated from:
$$A_\mathrm{V}=0.09\frac{1\mathrm{exp}\left(0.0111d\mathrm{sin}b\right)}{\mathrm{sin}b}$$
(17)
where $`b`$ is the absolute value of the galactic latitude. The colour excess is then calculated using Eq. (5). The factor in front was varied to fit the $`E(BV)`$ distribution. We also tried the reddening model of Arenou et al. (1992), but found that it could not fit the $`E(BV)`$ distribution.
The simulated ‘observed’ visual magnitude is calculated from $`V=M_\mathrm{V}+5\mathrm{log}d5+A_\mathrm{V}`$. Then a term is added simulating the uncertainty in the observed $`V`$, which is described by a Gaussian distribution with a dispersion of 0.005 mag.
hipparcos was complete only down to about $`V=7`$ and the following function was used to determine if a star was ‘observed’ or not. A random number ($`Rn`$) was drawn to calculate:
$$V_{\mathrm{lim}}=7.0\frac{\mathrm{log}(1.0\mathrm{Rn})}{C}$$
(18)
with $`C`$ empirically determined to be:
$$C=0.1070.030y0.00482y^2+0.00361y^3$$
(19)
with $`y=(V7)`$. A star is ‘observed’ if $`V<V_{\mathrm{lim}}`$.
The simulation is continued until 25 sets of 236 stars fulfill all criteria. Typically 100 000 stars needed to be drawn to arrive at this. The simulated distributions are depicted in Fig. 8 using the solid lines (normalised to the observed number of 236 stars). Typical parameters $`X_\mathrm{P}=0.70`$, $`X_{\sigma \mathrm{P}}=0.25`$, and $`H=70`$ pc have been used. The overall fit is good.
This numerical code for the quoted parameters provides us with a tool with which synthetic samples of Cepheids can be generated that obey the observed distributions.
The main difference with the simulations of L99 is in the conclusion about the space distribution of Cepheids. They assume a homogeneous 3D distribution, and consider a box centered on the Sun of 4200 pc on a side. They justify this because of “the relative small depth of the hipparcos survey \[with respect to\] the depth of the Galactic disk”. This is not true however. First of all, using any reasonable combination of $`\delta `$ and $`\rho `$ it is clear that hipparcos sampled Cepheids to much greater distances than 2.1 kpc, in fact out to 7-8 kpc. Furthermore, since the scale height of 3-10 M main-sequence stars (the progenitors of the Cepheids variables) is of order 100 pc it is clear that hipparcos sampled to distances much larger than the scale height of the Cepheid population. This is directly confirmed from our simulations from which we derive a scale height of about 70 pc. In other words, the space distribution of the Cepheids is a disk population, not a homogeneous population.
This might have consequences for the results of the L99 paper concerning biases which are difficult to judge by us. Another consequence is related to the fact that both Malmquist- and LK-bias depend on the underlying distribution of stars. For example, Oudmaijer et al. (1998) calculated the LK-bias assuming a homogeneous distribution of Cepheids. For a disk population the values of both the Malmquist- and LK-bias are smaller (Stobie et al. 1989, Koen 1992).
Now, we will discuss whether the zero points derived by the method outlined in Sect. 4 are subject to bias or not. Similar simulations were also performed by L99 and Pont (1999). L99 concluded that zero points derived for the whole sample in $`V`$ are too bright by 0.01 mag. Pont (1999) concluded that any bias is less than 0.03 mag.
What remains to be discussed in relation to our numerical model is how the ‘observed’ parallax and the error in the parallax are calculated.
We define several quantities. First of all a minimum error in the parallax, calculated from (in mas):
$$\begin{array}{cccc}\sigma _\pi ^{\mathrm{min}}\hfill & =\hfill & 0.45\hfill & V8\hfill \\ & =\hfill & 0.45\times (V/8)^{3.2}\hfill & \mathrm{else}\hfill \end{array}$$
(20)
Second of all, an error in the error on the parallax, calculated from (in mas):
$$\begin{array}{cccc}\sigma _\sigma \hfill & =\hfill & 0.18\hfill & V8\hfill \\ & =\hfill & 0.180.2167y+0.325y^20.06833y^3\hfill & \mathrm{else}\hfill \end{array}$$
(21)
with $`y=(V8)`$. Third, a ‘mean’ error on the parallax, calculated from (in mas):
$$\begin{array}{cccc}\sigma \hfill & =\hfill & 0.80\hfill & V5\hfill \\ & =\hfill & 0.800.0280y0.00229y^2+0.008651y^3\hfill & \mathrm{else}\hfill \end{array}$$
(22)
with $`y=(V5)`$. All fits have been made with idl using the routine polyfit from the full sample of 236 stars. It was verified that the synthetic samples have the same distribution of parallax and parallax error compared to the observed sample.
The error on the parallax, $`\sigma _\pi `$, was then determined from $`\sigma `$ plus a quantity randomly selected from a Gaussian distribution with dispersion $`\sigma _\sigma `$. The result is only accepted when $`\sigma _\pi >\sigma _\pi ^{\mathrm{min}}`$ however. Finally, the ‘observed’ parallax (in mas) is calculated from 1000/$`d`$, with $`d`$ the true distance in pc, plus a quantity randomly selected from a Gaussian distribution with dispersion $`\sigma _\pi `$.
The simulation is continued until 100 sets of 236, 191, or 63 stars fulfill all criteria, depending whether the simulation relates to $`V`$, $`I`$ or $`K`$. About 410 000 stars have to been drawn to arrive at 23 600 ‘observed’ stars in $`V`$. This was done with the parameter set that best described the observed distributions, as discussed above, i.e. with $`X_\mathrm{P}=0.70`$, $`X_{\sigma \mathrm{P}}=0.25`$, and $`H=70`$ pc.
In the case the simulation relates to $`I`$ and $`K`$ the procedure is as follows. First the stars are selected according to the procedure outlined above, that is, selection based on $`V`$-photometry. For the stars that fulfill the criteria, the ‘observed’ $`I`$, $`VI`$ (respectively $`K`$, $`JK`$) colours are determined, using the period-colour and reddening relations outlined in Sect. 4. The same random number that was used in Eqs. (15-16) to calculate the true $`M_\mathrm{V}`$ and $`(BV)`$ is used to calculate $`M_\mathrm{I}`$ and $`(VI)`$, respectively $`M_\mathrm{K}`$ and $`(JK)`$. This is to simulate the fact that if the synthetic star is ‘observed’ to be fainter and bluer than the mean in $`V`$, this is also the case at other wavelengths. This procedure ignores any phase shifts between the light curves at different colours. For the full-width of the instability strip (cf. Eqs. (15-16)) in $`M_\mathrm{I}`$ we take 0.70 mag, in $`(VI)`$ 0.20 mag, in $`M_\mathrm{K}`$ 0.60 and in $`(JK)`$ 0.16 mag (Gieren et al. 1998).
For every set of 236, 191, or 63 stars we apply the “reduced parallax” method and derive the zero point. From the distribution of the zero points of the 100 sets, we determine the mean, and the dispersion. The results are given in Table 6, where we list the zero point assumed in the numerical simulation, the zero point retrieved, the average number of stars in the simulation that fulfilled the selection criteria, and to which solution the simulation refers to. From the results we see that the zero point of the volume complete sample is dimmer than for the whole sample, as expected from Malmquist bias. In any case, the biases are very small, of order 0.01 mag, or less, similar to the results found by L99 and Pont (1999). Second, we confirm the result by Pont (1999) that the errors derived are somewhat larger compared to the outcome of the “reduced parallax” method. However the differences are not so large as compared to the error bars quoted in FC.
From the simulations we find that the smallest error is for the sample selected on parallax (as was found for the ‘real’ sample), but that it is subject to Lutz-Kelker bias.
## 8 Discussion
### 8.1 The finally adopted zero points
Based on the results obtained in Sects. 4, 5 and 6, we will present three sets of solutions, a ‘traditional’ one, and two alternatives. The traditional one follows FC and L99 closely. The zero point adopted is the one for the entire sample (which has the lowest error of the samples that are not selected on observed parallax), adopting the slope of the $`PL`$-relation observed for Cepheids in the LMC. This would then be solution 10 from Table 1 ($`\rho =1.41\pm 0.10`$), solution 1 from Table 2 ($`\rho =1.89\pm 0.11`$) and solution 1 from Table 3 corrected for the off-set in the corresponding solution in $`V`$ and $`I`$ as discussed in Sect. 4.3 ($`\rho =2.50\pm 0.17`$). These zero points have to be corrected for Malmquist bias. If the Malmquist bias would be evaluated in magnitude space, this bias would amount to 0.0092, 0.021 and 0.013 magnitude in $`V,I,K`$, respectively, for a disk population and the adopted values for $`\sigma _\mathrm{H}`$ of, respectively, 0.10, 0.15 and 0.12 mag (Stobie et al. 1989). However, if evaluated using the reduced parallax method the Malmquist bias is smaller (see Oudmaijer et al. 1999), and from the numerical simulations (Table 6) we find a Malmquist bias of 0.007 in $`V`$, 0.003 in $`I`$, and an unphysical negative value in $`K`$, probably due to the smaller number of stars involved. For the Malmquist bias we will assume a (round number) of 0.01 mag in all three bands.
Also, we have increased the errors in the zero points to reflect the sensitivity to uncertainties in the photometry, reddening and period-colour relations (see Sect. 5.1). For adopted slopes of $`2.81`$, $`3.05`$ and $`3.27`$, the finally adopted zero points of the $`PL`$-relations in $`V`$, $`I`$ and $`K`$ are, respectively, $`\rho =1.40\pm 0.11`$, $`1.88\pm 0.12`$ and $`2.49\pm 0.17`$. For the Wesenheit-index, after correcting for Malmquist bias and increasing the error bar due to the uncertainties described in Sect. 5.1, we adopt a zero point of $`2.55\pm 0.11`$ for a slope of $`3.411`$.
The ‘traditional’ solution above has certain advantages, like the fact that the errors are smallest compared to solutions that do not include all stars, or, since the slope adopted is the one in the LMC, a distance determination to the LMC is essentially a comparison of zero points only (apart from systematic effects). On the other hand, alternative solutions can be presented which also have merits. These solutions are based on a volume complete sample (at least in $`V`$ and $`I`$) to avoid Malmquist bias. Although small, its value does depend on the distribution of stars, and the intrinsic spread in the $`PL`$-relation and selecting a volume complete sample avoids Malmquist bias outright. In $`K`$, no volume complete sample could be constructed and we used the full sample instead. The alternative solutions take into account the theoretical prediction that the slope of the $`PL`$-relation may change at $`\mathrm{log}P0.5`$ for solar metallicities (Alibert et al. 1999). Where the two alternative solutions differ are in the adopted slopes for the Galactic $`PL`$-relations. One alternative solution takes into account the, admittedly at the 1$`\sigma `$ level of significance, evidence presented in Sects. 4.1-4.2 that the slope of the Galactic $`PL`$-relations are shallower than the ones in the LMC, and in accordance we have adopted the theoretical slopes predicted by Bono et al. (1999) for solar metallicities. The second alternative method adopts the slopes in Gieren et al. (1998), who derived distances from the infrared version of the surface brightness technique to Galactic Cepheids, which yields steeper slopes than for the LMC.
These alternative solutions are presented in Tables 1-3 (solutions 54-55, 28-28 and 19-20, respectively). The zero points in $`K`$ have to be corrected for Malmquist bias (+0.01 mag) and for the off-sets in the corresponding $`V`$ and $`I`$ solutions, and the errors in all three zero points are increased for reasons indicated above. For the adopted slopes of $`2.22`$, $`2.35`$ and $`3.05`$, the zero points of the $`PL`$-relations in $`V`$, $`I`$ and $`K`$ are, respectively, $`\rho =1.95\pm 0.12`$, $`2.47\pm 0.15`$ and $`2.68\pm 0.18`$. For the adopted slopes of $`3.04`$, $`3.33`$ and $`3.60`$, the zero points of the $`PL`$-relations in $`V`$, $`I`$ and $`K`$ are, respectively, $`\rho =1.30\pm 0.13`$, $`1.75\pm 0.14`$ and $`2.30\pm 0.18`$.
### 8.2 $`PL`$-relation of LMC Cepheids
Before proceeding we have to adopt $`PL`$-relations for the LMC Cepheids. In $`V`$, for a slope of $`2.81`$, Caldwell & Laney (1991) find a zero point of 17.23 $`\pm `$ 0.02. Tanvir (1999), for the same slope, gives an observed zero point of 17.451 $`\pm `$ 0.043. Adopting a mean reddening to the LMC of $`E(BV)`$ of 0.08 (Caldwell & Laney 1991) and a ratio of total-to-selective reddening of $`R_\mathrm{V}=A_\mathrm{V}/E(BV)=`$ 3.27 (Eq. (5) for typical colors) we derive a de-reddened zero point of 17.19 $`\pm `$ 0.04. For the Cepheids in the LMC we adopt the weighted mean of these two values, or a zero point of 17.22 $`\pm `$ 0.02 in $`V`$ for a slope of $`2.81`$.
In $`I`$, for a slope of $`3.041`$, Gieren et al. (1998) find a zero point of 16.74 $`\pm `$ 0.06 (the error is calculated by us, from their data). Tanvir (1999), for a slope of $`3.078`$, gives an observed zero point of 16.904 $`\pm `$ 0.031. Adopting $`A_\mathrm{I}=0.69A_\mathrm{V}=0.18`$ mag for the reddening in $`V`$ calculated as above, we derive a zero point of 16.72 $`\pm `$ 0.03. For the default slope of $`3.05`$ we adopt the weighted mean of these two values, or a de-reddened zero point in the $`I`$ band for the Cepheids in the LMC of 16.72 $`\pm `$ 0.02.
In $`K`$, for a slope of $`3.267`$, Gieren et al. (1998) find a zero point of 16.03 $`\pm `$ 0.05 (the error is calculated by us from their data) for the Cepheids in the LMC. This value is adopted by us.
For the Wesenheit-index, for a slope of $`3.411`$, Tanvir (1999) finds a zero point of 16.051 $`\pm `$ 0.017. This value is adopted by us.
### 8.3 Metallicity correction
We will now consider the effect of metallicity on the zero point. For comparison, FC applied a correction of +0.042 mag to the zero point in the $`V`$-band, based on Laney & Stobie (1994). The theoretical models of Bono et al. (1999), and Alibert et al. (1999) provide $`PL`$-relations and from those the difference $`\mathrm{\Delta }M=M`$(Gal) $`M`$(LMC) can be determined which will depend on the photometric band and period. We have determined this difference for two periods, namely for $`\mathrm{log}P_0=0.77`$ which we have determined to be the mean period of the volume complete sample of Galactic Cepheids in hipparcos and for $`\mathrm{log}P_0=0.47`$ which is the mean period of Cepheids in the LMC (Alcock et al. 1999). The results for $`\mathrm{\Delta }M`$ are listed in Table 7 for the three photometric bands. This illustrates the difference between the two theoretical models, for the Alibert et al. (1999) models predict essentially no dependence on metallicity, while the Bono et al. (1999) models predict a significant metallicity dependence, which mostly is in the sense that the metal-rich pulsators are fainter than the metal-poor ones. This is at variance with various empirical estimates that give the opposite result, and that, in the $`V`$-band, vary between $`0.24\pm 0.16`$ (Kennicutt et al. 1998) and about $`0.4`$ mag/dex (Kochanek 1997, Sasselov et al. 1997, Storm et al. 1999<sup>2</sup><sup>2</sup>2Recently, Storm (2000) suggested that this result may not be confirmed from his latest analysis and that the correction may have a positive sign instead.). A mean of these four determinations is $`0.38\pm 0.09`$ mag/dex, which for a difference in metallicity of 0.4 dex, implies $`\mathrm{\Delta }M=0.15\pm 0.04`$ in the $`V`$-band. In the $`I`$-band we assume the same value following Sasselov et al. (1997) and Kochanek (1997), but the reader should realise that this number is less well established than the correction in $`V`$, and in the $`K`$-band adopt $`\mathrm{\Delta }M=0.07`$. Note however, that a metallicity dependence as large as 0.4 mag in $`V`$ as suggested by Sekiguchi & Fukugita (1998) can be excluded at the 6$`\sigma `$ level (Laney 1999, 2000). In recent papers, Saio & Gautschy (1998) and Sandage et al. (1999) found no significant metallicity dependence on the bolometric $`PL`$-relation, and slopes of $`0.08`$ mag/dex in $`V`$ and $`0.1`$ mag/dex in $`I`$ (Sandage et al. 1999), which represents a shallower dependence than the values listed above, that are adopted in the present study, and which therefore may be an extreme view.
### 8.4 The distance to the LMC
In Table 8 are listed the true distance moduli (DM) to the LMC for $`V,I,K`$, the three slopes (‘traditional’ meaning the observed slopes for the Cepheids in the LMC, ‘shallower’ adopting the theoretical slopes from Bono et al. and ‘steeper’ adopting the observed slopes for Galactic Cepheids from Gieren et al. (1998) and the three metallicity corrections (‘0’ means no correction, ‘+’ means a longer distance scale as implied by the models from Bono et al., and ‘$``$’ means a shorter distance scale as implied from empirical evidence). The error quoted includes the error in the zero point of both the Galactic and LMC Cepheids, and where appropriate, the error due to the metallicity correction, and for the solutions with either steeper or shallower slope, the uncertainty due to the difference in DM at $`\mathrm{log}P=0.47`$ and 0.77.
Also included are the weighted mean DM, averaged over $`V,I,K`$ (with internal error), and, for reference, the (unweighted) mean DM per photometric band of all the solutions (with the one-sided range in the solutions). Also included is the solution based on the Wesenheit-index assuming the slope observed in the LMC, and no metallicity correction. The DM range from 18.45 $`\pm `$ 0.18 to 18.86 $`\pm `$ 0.12. Several important conclusions may be drawn:
(1) For every combination, the $`PL`$-relation in $`K`$ gives the shortest distance, and the difference between the distance based on $`V`$ and $`K`$ can be as large as 0.24 mag (solutions 4, 5 in Table 8).
This systematic effect is worrying and merits investigation. It could hint to errors in the reddening, or dereddening procedure. It is illustrative to note that the (minimum, maximum, mean) extinction for the whole sample is ($`0.21`$, 3.4, 1.3) in $`A_\mathrm{V}`$, ($`0.08`$, 2.1, 0.75) in $`A_\mathrm{I}`$, and ($`0.05`$, 0.26, 0.10) in $`A_\mathrm{K}`$. This implies that any uncertainty in reddening is less in $`K`$. In Sect. 5.1 we have estimated these uncertainties (about 0.04 in $`V`$, 0.03 in $`I`$ and 0.01 mag in $`K`$) and added them as a random errors. Possibly these are errors of a systematic nature instead. It is interesting to note that application of the procedures outlined in Sect. 3.2 results in negative reddening in some cases.
One can raise the question how much bluer the period-colour-relations need to be to give positive reddening for all stars. It turns out that Eq. (4) needs to be bluer by 0.065 mag, Eq. (8) by 0.055 mag, and Eq. (11) by 0.075 mag. For the default slopes and using all stars, the zero points in $`V,I,K`$ would change to, respectively, $`1.620`$ (from $`1.411`$), $`1.970`$ (from $`1.892`$), and $`2.660`$ (from $`2.607`$).
On the other hand, the DM based on the Wesenheit-index, which avoids the use of a $`PC`$-relation to estimate the individual reddenings, is in perfect agreement with the DM based both on the $`PL`$-relations in $`V`$ and $`I`$. This suggests that the reddening is not the main reason for the systematically shorter DM in the $`K`$-band.
As pointed out in Sect. 4.3 there may be a bias in the $`K`$-band zero point due to the smaller number of Cepheids with accurate intensity-mean magnitudes. The correction for this bias was estimated by comparing, for the same sample of stars with $`K`$-band data, the zero point of the $`PL`$-relations in $`V`$ and $`I`$ to those for the full samples in $`V`$ and $`I`$, and this correction is about 0.1 mag, in the sense that it makes the DM based on $`K`$ shorter than they would be without this correction. This uncertainty can only be eliminated if more intensity-mean NIR magnitudes come available.
(2) The uncertainty in the type of metallicity correction introduces a range in DM of up to 0.20 mag in $`V,I`$, and 0.12 mag in $`K`$.
(3) The uncertainty in the slope of the Galactic $`PL`$-relations introduces a range in DM of about 0.16 mag in $`V,I`$, and about 0.05 mag in $`K`$.
Taking the case with the observed slopes of LMC Cepheids with no metallicity correction as default, one may summarise the results as follows. Based on the $`PL`$-relation in $`V`$ and $`I`$, and the Wesenheit-index, the true DM to the LMC is 18.60 $`\pm `$ 0.11 ($`\pm `$ 0.08 slope) ($`{}_{0.15}{}^{}{}_{}{}^{+0.08}`$ metallicity). Based on the $`PL`$-relation in $`K`$ it is 18.52 $`\pm `$ 0.18 ($`\pm `$ 0.03 slope) ($`\pm 0.06`$ metallicity) ($`{}_{0}{}^{}{}_{}{}^{+0.10}`$ sample bias). The terms between parenthesis indicate the possible systematic uncertainties due to the slope of the Galactic $`PL`$-relations, the metallicity corrections, and in the $`K`$-band, due to the limited number of stars. Recent work by Sandage et al. (1999) indicate that the effect of metallicity towards shorter distances may be smaller in $`V`$ and $`I`$ than indicated here. A more accurate determination is not possible without more definite information on the slope of the Galactic $`PL`$-relations and the metallicity correction.
Our prefered distance modulus is the one based on the $`PL`$-relation in $`V`$, $`I`$ and the Wesenheit index, and puts the LMC 0.10 mag in DM closer than the value of 18.70 derived by FC. The difference is due to four effects that all work in the same direction, namely, (1) FC apply a metallicty correction of +0.042 mag, (2) the difference in the zero point in the Galactic $`PL`$-relation in $`V`$ between FC and this study is +0.03 mag (+0.01 mag is due to Malmquist bias which FC did not take into account, while +0.02 mag is due to the different sample and slightly different photometry in some cases), (3) the difference in the DM based on $`V`$ compared to the mean of the DM based on $`V`$, $`I`$ and the Wesenheit-index is +0.02 mag, and (4) the difference between FC and this study in the adopted zero point of the $`PL`$-relation in $`V`$ of the LMC Cepheids is +0.01 mag.
We finally note that the influence of the choice of slope and the metallicity correction are (predicted to be) smallest in the $`K`$-band as well as the uncertainty in the extinction correction. If the 20-30 closest Cepheids without published NIR photometry could have their NIR intensity-mean magnitudes determined, then the uncertainty due to the small number of stars could be eliminated and the zero point in $`K`$ could be determined with an error that is a factor of two smaller than is possible at present.
### Acknowledgements
We thank Pascal Fouqué, Michael Feast and Patricia Whitelock for providing tabular material in Gieren et al. (1998), respectively Feast & Whitelock (1997) in electronic format. We thank Giuseppe Bono and Santi Cassisi in calculating and communicating additional $`PL(C)`$-relations to us. Frederic Pont and Frederic Arenou are thanked for lively and interesting discussions. This research has made use of the SIMBAD database, operated at CDS, Strasbourg, France.
## Appendix
In this appendix we list the sample of 236 Cepheids in the hipparcos catalog considered in this study. Listed in Table A1 are the HIP number and variable star name, the parallax and error in the parallax from the hipparcos catalog except for RY Sco and Y Lac (Falin & Mignard 1999), the intensity-mean $`V`$ and $`BV`$ adopted, the $`\mathrm{log}`$ of the fundamental period, and the assumed pulsation mode. The $`VI`$ colours and $`JHK`$ photometry can be found in G99, as described in Sect. 2. |
warning/0002/math0002100.html | ar5iv | text | # On the combinatorics of Forrester-Baxter modelsResearch supported by the Australian Research Council (ARC)
## 0 Introduction
### 0.1 Motivation
The physical spectrum of exactly-solvable lattice models can be described in the language of highest-weight infinite dimensional representations of affine and Virasoro algebras . The characters of these representations are $`q`$-series that contain detailed information on the structure and symmetries of the corresponding models. In the following discussion, we wish to restrict attention to the characters of Virasoro highest-weight representations.
The earliest known expressions for these characters are due to Feigen and Fuchs and Rocha-Caridi . These expressions have alternating-signs. A number of years ago, the Stony Brook group discovered completely new expressions for the character formulae<sup>2</sup><sup>2</sup>2For references to the original Stony Brook papers, please refer to .. These expressions have constant-signs<sup>3</sup><sup>3</sup>3The characterisation of the different expressions of the characters as ‘alternating-sign’ and ‘constant-sign’ $`q`$-series is valid only for Virasoro but not for affine characters..
For physical reasons that are beyond the scope of this work, the original alternating-sign expressions are also known as ‘bosonic characters’. Correspondingly, the constant-sign expressions are also known as ‘fermionic characters’ <sup>4</sup><sup>4</sup>4For a complete discussion of the physical motivation of the terms ‘bosonic’ and ‘fermionic’, please refer to the original literature on the subject as cited in ..
The structure of these new character formulae hints at the presence of a completely new formulation of exactly-solvable models<sup>5</sup><sup>5</sup>5Analogous developments in the context of highest-weight representations of affine algebras also took place. They are outside the scope of this work.. This possibility has attracted attention for a number of reasons. One of these reasons is the fact that certain physical problems, such as the long-distance asymptotics of the correlation functions, are too difficult to handle in the current formulation. Further, there are reasons to believe that the new formulation could be the right starting point to tackle them (see and references therein). At a more technical level, the availability of two distinct formulations is mathematically enriching, as we can use one to learn about the other.
However, although the bosonic characters are technically simple to write down, and are completely known for all Virasoro representations, the structure of the fermionic characters is strictly-speaking known explicitly only in special cases, and generally only conceptually. In particular, the characters of the ‘non-unitary’ Virasoro representations have turned out to be rather resistant to a complete formulation in fermionic form<sup>6</sup><sup>6</sup>6The reason for that may of course eventually turn out to be the fact that we are not using the most efficient tools to tackle this problem..
This work is part of a series of papers that aim at a complete and explicit derivation of the fermionic characters of a certain class of models first discussed by Forrester and Baxter . The characters of the Forrester-Baxter models correspond to the complete set of Virasoro characters of the discrete, though not necessarily unitary, Virasoro algebras with central charge $`c<1`$, first discussed in . As such, they form the largest class of Virasoro characters with no $`W`$-symmetries.
As in previous works, our approach is purely combinatorial. Further, the exposition is self-contained, in the sense that we have included all concepts required in the derivations. Our main result is a combinatorial derivation of two related finitised fermionic forms for the characters of a certain class of Forrester-Baxter models. The first of these requires the use of the classical form of Gaussian polynomials and can be interpreted combinatorially using the concept of particles. The second has already appeared in the works of Berkovich, McCoy and Schilling , requires the use of a modified form of Gaussian polynomials, and has a combinatorial interpretation in terms of particles and particle annihilation.
In a forthcoming paper, we further extend and refine the techniques of this work to obtain a complete and explicit derivation of the fermionic characters of the complete set of Forrester-Baxter models .
### 0.2 Overview of content of paper
The aim of this paper is to obtain fermionic expressions for $`\chi _{a,b,c}^{p,p^{}}(L)`$, the generating function for the set $`𝒫_{a,b,c}^{p,p^{}}(L)`$ of restricted length-$`L`$ paths that have startpoint $`a`$ and endpoint $`b`$.
These functions<sup>7</sup><sup>7</sup>7To be precise, a certain renormalisation thereof. first arose in the calculation of one-point functions of the Forrester-Baxter models . The weighting originally assigned in to the paths is significantly different from that used here. The weighting described in the current paper arose by obtaining a ‘weight-preserving’ bijection between partitions with prescribed hook-differences that were considered in , and the paths of . This bijection is described in .
The paths in $`𝒫_{a,b,c}^{p,p^{}}(L)`$ may be depicted on a $`(p^{}2)\times L`$ grid that we refer to as the $`(p,p^{})`$-model, as described in Section 1.1. Of particular importance is the shading of the $`(p,p^{})`$-model, which determines the weights $`wt(h)`$ that we assign to the paths $`h`$.
A bosonic expression for $`\chi _{a,b,c}^{p,p^{}}(L)`$ is given in Section 1.3. This expression is readily proved using $`L`$-recurrence relations , or by using the generating function for partitions with prescribed hook-differences given in , and the bijection of . The polynomial $`\chi _{a,b,c}^{p,p^{}}(L)`$ is seen to be a finitisation of a Virasoro character.
In this paper, we tackle the particular cases where $`a`$ and $`b`$ are each one of the Takahashi lengths $`𝒯`$, or one of $`𝒯^{}=\{p^{}s:s𝒯\}`$. These values depend on $`p`$ and $`p^{}`$, and are defined in Section 5.1. Our methods and results are a common generalisation of those of .
On equating the bosonic expression for $`\chi _{a,b,c}^{p,p^{}}(L)`$ with either of the fermionic expressions, we obtain boson-fermion polynomial identities. Taking the $`L\mathrm{}`$ limit (using, for example, the variable change employed in ), these become $`q`$-series identities. Amongst them, in particular, are the Rogers-Ramanujan identities, and their generalisations by Andrews and Gordon . In fact, the techniques employed by Agarwal and Bressoud in their combinatorial proof of the Andrews-Gordon identities provided the genesis of the techniques employed here.
Before we develop a generalisation of Agarwal and Bressoud’s ‘Volcanic activity’, we define in Section 2, a slightly different set $`𝒫_{a,b,e,f}^{p,p^{}}(L)`$ of paths, which have assigned pre-segments and post-segments that are determined by $`e,f\{0,1\}`$. Their generating function $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L)`$ is defined in terms of a path weighting that differs slightly from that defined earlier.
The $``$-transform, which is described in Section 3, enables $`\stackrel{~}{\chi }_{a^{},b^{},e,f}^{p,p^{}+p}(L^{})`$, for certain $`a^{},b^{}`$ to be expressed in terms of $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}+p}(L)`$. We derive this transform combinatorially in three steps. The first step is known as the $`_1`$-transform and enlarges the features of a path, so that the resultant path resides in a larger model. The second step, referred to as a $`_2(k)`$-transform, lengthens a path by appending $`k`$ pairs of segments to the path. Each of these pairs is known as a particle. The third step, the $`_3(\lambda )`$-transform deforms the path in a particular way. This process may be viewed as the particles moving through the path. The resulting transformation of generating functions is given in Corollary 3.14.
In Section 4, we see that $`\stackrel{~}{\chi }_{a,b,1e,1f}^{p^{}p,p^{}}(L)`$ may be obtained from $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L)`$ in a combinatorially trivial way. This process is referred to as a $`𝒟`$-transform. In fact, it is more convenient to use the $`𝒟`$-transform combined with the $``$-transform. The resulting transformation of generating functions is given in Corollary 4.6.
To obtain a particular generating function $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L)`$, where $`p`$ and $`p^{}`$ are co-prime, we begin with one of the trivial generating functions $`\stackrel{~}{\chi }_{a^{},b^{},e^{},f^{}}^{1,3}(L)`$ given in Lemma 2.5, and perform a sequence of $``$\- and $`𝒟`$-transforms. This sequence is determined by the continued fraction of $`p^{}/p`$ which is described in Section 5.1.
In fact, a basic application of the transforms does not generate all elements of $`𝒫_{a,b,e,f}^{p,p^{}}(L)`$ in some cases. In these instances, the set generated is deficient in the full set of paths that do not rise above (or below) a certain height. Various results obtained in Section 6 enable us to keep track of this height. Lemma 6.4 shows that this height bounds a portion of the $`(p,p^{})`$-model which is identical to a smaller $`(\widehat{p},\widehat{p}^{})`$-model. This property enables (in one case), the final generating function to be expressed using the generating function for paths in the $`(\widehat{p},\widehat{p}^{})`$-model.
Section 7 provides one further ingredient for the final construction. There, it is shown how appending or removing the first segment of the path affects the generating function.
Everything is now in place to carry out the proof of the main results. These results are stated in Section 8.1. We provide two similar expressions for $`\chi _{a,b,c}^{p,p^{}}(L)`$. These are Theorems 8.1 and 8.2. The first of these makes use of the classical definition of the Gaussian polynomial:
$$\left[\genfrac{}{}{0pt}{}{A}{B}\right]_q=\{\begin{array}{cc}\frac{\left(q\right)_A}{\left(q\right)_{AB}\left(q\right)_B}& \text{if }0BA;\hfill \\ 0& \text{otherwise},\hfill \end{array}$$
(1)
where $`(q)_0=1`$ and $`(q)_n=_{i=1}^n(1q^i)`$ for $`n>0`$. In some cases, the expression also includes a term $`\chi _{a,b,c}^{\widehat{p},\widehat{p}^{}}(L)`$ for $`\widehat{p}^{}<\widehat{p}`$. Thus this expression may be viewed as a recursive fermionic expression for $`\chi _{a,b,c}^{p,p^{}}(L)`$. In the cases where this additional term is not present (for $`a`$ and $`b`$ further restricted in a certain way), the expressions were first stated in .
The expression of Theorem 8.2 makes use of a modified definition of the Gaussian polynomial ():
$$\left[\genfrac{}{}{0pt}{}{A}{B}\right]_q^{}=\{\begin{array}{cc}\frac{\left(q^{AB+1}\right)_B}{\left(q\right)_B}& \text{if }0B;\hfill \\ 0& \text{otherwise},\hfill \end{array}$$
(2)
where $`(z)_0=1`$ and $`(z)_n=_{i=0}^{n1}(1zq^i)`$ for $`n>0`$. These expressions were first presented and proved in . In fact, invoking the definition (2) is somewhat overkill, since the only value of $`\left[\genfrac{}{}{0pt}{}{A}{B}\right]^{}`$ that we require that differs from $`\left[\genfrac{}{}{0pt}{}{A}{B}\right]`$ is $`\left[\genfrac{}{}{0pt}{}{1}{0}\text{ }\right]^{}=1`$.
In , expressions for $`\chi _{a,b,c}^{p,p^{}}(L)`$ are presented, where $`b`$ is now any value with $`1bp^{}1`$. However, only $`a𝒯𝒯^{}`$ is still permitted. In , we show that it is Theorem 8.1, and not Theorem 8.2, that generalises to provide fermionic expressions for the most general $`\chi _{a,b,c}^{p,p^{}}(L)`$.
The remainder of Section 8 is concerned with the detailed derivation of the expression for first $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L)`$, and then converting it to $`\chi _{a,b,c}^{p,p^{}}(L)`$. Section 8.3 describes the $`𝒎`$$`𝒏`$-system which aids the actual evaluation of the fermionic expressions obtained. Section 8.4 describes how the proof for Theorem 8.1 modifies to provide a proof for Theorem 8.2. Here we see that the appearance of $`\left[\genfrac{}{}{0pt}{}{1}{0}\text{ }\right]^{}`$ may be viewed in terms of ‘particle annihilation’.
## 1 Paths
### 1.1 Paths and the $`(p,p^{})`$-model
Let $`p`$ and $`p^{}`$ be positive co-prime integers for which $`0<p<p^{}`$. Then, given $`a,b,c,L\mathrm{}_0`$ such that $`1a,b,cp^{}1`$, $`b=c\pm 1`$, $`L+ab0`$ ($`\text{mod}\mathrm{\hspace{0.17em}2}`$), a path $`h𝒫_{a,b,c}^{p,p^{}}(L)`$ is a sequence $`h_0,h_1,h_2,\mathrm{},h_L,`$ of integers such that:
1. $`1h_ip^{}1`$ for $`0iL`$,
2. $`h_{i+1}=h_i\pm 1`$ for $`0i<L`$,
3. $`h_0=a,h_L=b.`$
Note that the values of $`p`$ and $`c`$ do not feature in the above restrictions. As described below, they specify how the elements of $`𝒫_{a,b,c}^{p,p^{}}(L)`$ are weighted.
The integers $`h_0,h_1,h_2,\mathrm{},h_L,`$ are readily depicted as a sequence of heights on a two-dimensional $`L\times (p^{}2)`$ grid. Adjacent heights are connected by line segments passing from $`(i,h_i)`$ to $`(i+1,h_{i+1})`$ for $`0i<L`$.
Scanning the path from left to right, each of these line segments points either in the NE direction or in the SE direction. Fig. 1 shows a typical path in the set $`𝒫_{2,4,3}^{3,8}(14)`$. The shadings in Fig. 1 are explained below.
In the grid introduced above, the horizontal strip between adjacent heights is referred to as a band. There are $`p^{}2`$ bands. The $`h`$th band lies between heights $`h`$ and $`h+1`$.
We now assign a parity to each band: the $`h`$th band is said to be an even band if $`hp/p^{}=(h+1)p/p^{}`$; and an odd band if $`hp/p^{}(h+1)p/p^{}`$. The array of odd and even bands so obtained will be referred to as the $`(p,p^{})`$-model. It may immediately be deduced that the $`(p,p^{})`$-model has $`p^{}p1`$ even bands and $`p1`$ odd bands. In addition, it is easily shown that for $`1r<p`$, the band lying between heights $`rp^{}/p`$ and $`rp^{}/p+1`$ is odd: it will be referred to as the $`r`$th odd band.
When drawing the $`(p,p^{})`$-model, we distinguish the bands by shading the odd bands. This was done in Fig. 1 for the $`(3,8)`$-model.
We note that the band structure of the $`(p,p^{})`$-model is up-down symmetrical, and that if $`p^{}>2p`$ then the 1st band and the $`(p^{}2)`$th band are both even, and there are no two adjacent odd bands.
For $`2ap^{}2`$, we say that $`a`$ is interfacial if $`(a+1)p/p^{}=(a1)p/p^{}+1`$. Thus $`a`$ is interfacial if and only if $`a`$ lies between an odd and even band in the $`(p,p^{})`$-model. Thus for the case of the $`(3,8)`$-model depicted in Fig. 1, $`a`$ is interfacial for $`a=2,3,5,6`$. Note that if $`a`$ is interfacial, the odd band that it borders is the $`(a+1)p/p^{}`$th.
As is easily seen, the $`(p^{}p,p^{})`$-model differs from the $`(p,p^{})`$-model in that each band has changed parity. It follows that if $`a`$ is interfacial in the $`(p,p^{})`$-model then $`a`$ is also interfacial in the $`(p^{}p,p^{})`$-model.
### 1.2 Weighting the paths
Given a path $`h`$ of length $`L`$, for $`1i<L`$, the values of $`h_{i1}`$, $`h_i`$ and $`h_{i+1}`$ determine the shape of the vertex at the point $`i`$. The four possible shapes are given in Fig. 2.
The four types of vertices shown in Fig. 2 are referred to as a straight-up vertex, a straight-down vertex, a peak-up vertex and a peak-down vertex respectively. Each vertex is also assigned a parity: this is the parity of the band in which the segment between $`(i,h_i)`$ and $`(i+1,h_{i+1})`$ lies. Thus, there are eight types of paritied vertex.
For paths $`h𝒫_{a,b,c}^{p,p^{}}(L)`$, we define $`h_{L+1}=c`$, whereupon the shape and parity of the vertex at $`i=L`$ is well-defined.
The weight function for the paths is best specified in terms of a $`(x,y)`$-coordinate system which is inclined at $`45^o`$ to the original $`(i,h)`$-coordinate system and whose origin is at the path’s initial point at $`(i=0,h=a)`$. Specifically,
$$x=\frac{i(ha)}{2},y=\frac{i+(ha)}{2}.$$
Note that at each step in the path, either $`x`$ or $`y`$ is incremented and the other is constant. In this system, the path depicted in Fig. 1 has its first few coordinates at $`(0,0)`$, $`(0,1)`$, $`(0,2)`$, $`(0,3)`$, $`(1,3)`$, $`(1,4)`$, $`(1,5)`$, $`(1,6)`$, $`(2,6)`$, $`\mathrm{}`$
Now, for $`1iL`$, we define the weight $`c_i=c(h_{i1},h_i,h_{i+1})`$ of the $`i`$th vertex according to its shape, its parity and its $`(x,y)`$-coordinate, as specified in Table 1.
In Table 1, the lightly shaded bands can be either even or odd bands (or when $`h_i=p^{}1`$ or $`h_i=1`$ in the lowermost four cases, not a band in the model at all). Note that for each vertex shape, only one parity case has non-zero weight in general. We shall refer to those four vertices, with assigned parity, for which in general, the weight is non-zero, as scoring vertices. The other four vertices will be termed non-scoring.
We now define:
$$wt(h)=\underset{i=1}{\overset{L}{}}c_i.$$
(3)
To illustrate this procedure, consider again the path $`h`$ depicted in Fig. 1. The above table indicates that there are scoring vertices at $`i=3`$, $`4`$, $`5`$, $`7`$, $`8`$, $`13`$ and $`14`$. This leads to
$$wt(h)=0+3+1+1+6+7+6=24.$$
The generating function $`\chi _{a,b,c}^{p,p^{}}(L)`$ for the set of paths $`𝒫_{a,b,c}^{p,p^{}}(L)`$ is defined to be:
$$\chi _{a,b,c}^{p,p^{}}(L;q)=\underset{h𝒫_{a,b,c}^{p,p^{}}(L)}{}q^{wt(h)}.$$
(4)
Often, we drop the base $`q`$ from the notation so that $`\chi _{a,b,c}^{p,p^{}}(L)=\chi _{a,b,c}^{p,p^{}}(L;q)`$. The same will be done for other functions without comment.
### 1.3 Bosonic generating function
By setting up recurrence relations for $`\chi _{a,b,c}^{p,p^{}}(L)`$, it may be readily verified that:
$`\chi _{a,b,c}^{p,p^{}}(L)`$ $`=`$ $`{\displaystyle \underset{\lambda =\mathrm{}}{\overset{\mathrm{}}{}}}q^{\lambda ^2pp^{}+\lambda (p^{}rpa)}\left[{\displaystyle \genfrac{}{}{0pt}{}{L}{\frac{L+ab}{2}p^{}\lambda }}\right]_q`$
$`{\displaystyle \underset{\lambda =\mathrm{}}{\overset{\mathrm{}}{}}}q^{(\lambda p+r)(\lambda p^{}+a)}\left[{\displaystyle \genfrac{}{}{0pt}{}{L}{\frac{L+ab}{2}p^{}\lambda a}}\right]_q,`$
where
$$r=pc/p^{}+(bc+1)/2.$$
(6)
In the limit $`L\mathrm{}`$, we obtain
$$\underset{L\mathrm{}}{lim}\chi _{a,b,c}^{p,p^{}}(L)=\chi _{r,a}^{p,p^{}},$$
(7)
where $`r`$ is defined in (6) and
$$\chi _{r,s}^{p,p^{}}=\frac{1}{(q)_{\mathrm{}}}\underset{\lambda =\mathrm{}}{\overset{\mathrm{}}{}}(q^{\lambda ^2pp^{}+\lambda (p^{}rps)}q^{(\lambda p+r)(\lambda p^{}+s)})$$
(8)
is, up to a normalisation, the Rocha-Caridi expression for the Virasoro character of central charge $`c=16(p^{}p)^2/pp^{}`$ and conformal dimension $`\mathrm{\Delta }_{r,s}^{p,p^{}}=\left((p^{}rps)^2(p^{}p)^2\right)/4pp^{}`$. Therefore, $`\chi _{a,b,c}^{p,p^{}}(L)`$ provides a finite analogue of the character $`\chi _{r,a}^{p,p^{}}`$.
## 2 Winged generating functions
For $`h𝒫_{a,b,c}^{p,p^{}}(L)`$, the values of $`b`$ and $`c`$ serve to specify a path post-segment that extends between $`(L,b)`$ and $`(L+1,c)`$. We now define another set of paths which specifies both the direction of a post-segment and a pre-segment.
Let $`p`$ and $`p^{}`$ be positive co-prime integers for which $`0<p<p^{}`$. Then, given $`a,b,L\mathrm{}_0`$ such that $`1a,bp^{}1`$, $`L+ab0`$ ($`\text{mod}\mathrm{\hspace{0.17em}2}`$), and $`e,f\{0,1\}`$, a path $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ is a sequence $`h_0,h_1,h_2,\mathrm{},h_L,`$ of integers such that:
1. $`1h_ip^{}1`$ for $`0iL`$,
2. $`h_{i+1}=h_i\pm 1`$ for $`0i<L`$,
3. $`h_0=a,h_L=b.`$
If $`f=0`$ (resp. $`f=1`$) then the post-segment of each $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ is defined to be in the NE (resp. SE) direction. If $`e=0`$ (resp. $`e=1`$) then the pre-segment of each $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ is defined to be in the SE (resp. NE) direction. This enables a shape and a parity to be assigned to both the zeroth and the $`L`$th vertices of $`h`$. For $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$, we define $`e(h)=e`$ and $`f(h)=f`$.
We now define a weight $`\stackrel{~}{wt}(h)`$, for $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$. For $`1i<L`$, set $`\stackrel{~}{c}_i=c(h_{i1},h_i,h_{i+1})`$ as above. Then, set
$$\stackrel{~}{c}_L=\{\begin{array}{cc}x\hfill & \text{if }h_Lh_{L1}=1\text{ and }f(h)=1;\hfill \\ y\hfill & \text{if }h_Lh_{L1}=1\text{ and }f(h)=0;\hfill \\ 0\hfill & \text{otherwise,}\hfill \end{array}$$
where $`(x,y)`$ is the coordinate of the $`L`$th vertex of $`h`$. We then designate this vertex as scoring if it is a peak vertex ($`h_L=h_{L1}(1)^{f(h)}`$), and as non-scoring otherwise.
We define:
$$\stackrel{~}{wt}(h)=\underset{i=1}{\overset{L}{}}\stackrel{~}{c}_i.$$
(9)
Consider the corresponding path $`h^{}𝒫_{a,b,c}^{p,p^{}}(L)`$ with $`c=b+(1)^f`$, defined by $`h_i^{}=h_i`$ for $`0iL`$. From Table 1, we see that $`\stackrel{~}{wt}(h)=wt(h^{})`$ if the post-segment of $`h`$ lies in an even band.
In what follows, we work entirely in terms of $`\stackrel{~}{wt}(h)`$, and the generating functions that we derive from it. Only at the end of our work, do we revert back to $`wt(h)`$ to obtain fermionic expressions for $`\chi _{a,b,c}^{p,p^{}}(L)`$.
Define the generating function
$$\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L;q)=\underset{h𝒫_{a,b,e,f}^{p,p^{}}(L)}{}q^{\stackrel{~}{wt}(h)},$$
(10)
where $`\stackrel{~}{wt}(h)`$ is given by (9). Of course, $`\stackrel{~}{\chi }_{a,b,0,f}^{p,p^{}}(L)=\stackrel{~}{\chi }_{a,b,1,f}^{p,p^{}}(L)`$.
### 2.1 Striking sequence of a path
For each path $`h`$, define $`\pi (h)\{0,1\}`$ to be the parity of the band between heights $`h_0`$ and $`h_1`$ (if $`L(h)=0`$, we set $`h_1=h_0+(1)^{f(h)}`$). Thus, for the path $`h`$ shown in Fig. 1, we have $`\pi (h)=1`$. In addition, define $`d(h)=0`$ when $`h_1h_0=1`$ and $`d(h)=1`$ when $`h_1h_0=1`$. We then see that if $`e(h)+d(h)+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2})`$ then the $`0`$th vertex is a scoring vertex, and if $`e(h)+d(h)+\pi (h)1(\text{mod}\mathrm{\hspace{0.17em}2})`$ then it is a non-scoring vertex.
Now consider each path $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ as a sequence of straight lines, alternating in direction between NE and SE. Then, reading from the left, let the lines be of lengths $`w_1`$, $`w_2`$, $`w_3,\mathrm{},w_l,`$ for some $`l`$, with $`w_i>0`$ for $`1il`$. Thence $`w_1+w_2+\mathrm{}+w_l=L(h)`$, where $`L(h)=L`$ is the length of $`h`$.
For each of these lines, the last vertex will be considered to be part of the line but the first will not. Then, the $`i`$th of these lines contains $`w_i`$ vertices, the first $`w_i1`$ of which are straight vertices. Then write $`w_i=a_i+b_i`$ so that $`b_i`$ is the number of scoring vertices in the $`i`$th line. The striking sequence of $`h`$ is then the array:
$$\left(\begin{array}{ccccc}a_1& a_2& a_3& \mathrm{}& a_l\\ b_1& b_2& b_3& \mathrm{}& b_l\end{array}\right)^{(e(h),f(h),d(h))}.$$
With $`\pi =\pi (h)`$, $`e=e(h)`$ and $`d=d(h)`$, we define
$$m(h)=\{\begin{array}{cc}(e+d+\pi )\text{mod}\mathrm{\hspace{0.17em}2}+_{i=1}^la_i\hfill & \text{if }L>0;\hfill \\ |fe|\hfill & \text{if }L=0,\hfill \end{array}$$
whence $`m(h)`$ is the number of non-scoring vertices possessed by $`h`$ (altogether, $`h`$ has $`L(h)+1`$ vertices). We also define $`\alpha (h)=(1)^d((w_1+w_3+\mathrm{})(w_2+w_4+\mathrm{}))`$ and for $`L>0`$,
$$\beta (h)=\{\begin{array}{c}(1)^d((b_1+b_3+\mathrm{})(b_2+b_4+\mathrm{}))\hfill \\ \text{if }e+d+\pi 0(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ (1)^d((b_1+b_3+\mathrm{})(b_2+b_4+\mathrm{}))+(1)^e\hfill \\ \text{otherwise.}\hfill \end{array}$$
For $`L=0`$, we set $`\beta (h)=fe`$.
For example, for the path shown in Fig. 1 for which $`d(h)=0`$ and $`\pi (h)=1`$, the striking sequence is:
$$\left(\genfrac{}{}{0pt}{}{2011120}{1121011}\right)^{(e,1,0)}.$$
In this case, $`m(h)=8e`$, $`\alpha (h)=2`$, and $`\beta (h)=2e`$.
We note that given the startpoint $`h_0=a`$ of the path, the path can be reconstructed from its striking sequence<sup>8</sup><sup>8</sup>8We only need $`w_1,w_2,\mathrm{},w_l`$ together with $`d`$.. In particular, $`h_L=b=a+\alpha (h)`$. In addition, the nature of the final vertex may be deduced from $`a_l`$ and $`b_l`$<sup>9</sup><sup>9</sup>9Thus the value of $`f`$ in the striking sequence is redundant — we retain it for convenience.
###### Lemma 2.1
Let the path $`h`$ have the striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)},`$ with $`w_i=a_i+b_i`$ for $`1il`$. Then
$$\stackrel{~}{wt}(h)=\underset{i=1}{\overset{l}{}}b_i(w_{i1}+w_{i3}+\mathrm{}+w_{1+imod2}).$$
Proof: For $`L=0`$, both sides are clearly $`0`$. So assume $`L>0`$. First consider $`d=0`$. For $`i`$ odd, the $`i`$th line is in the NE direction and its $`x`$-coordinate is $`w_2+w_4+\mathrm{}+w_{i1}`$. By the prescription of the previous section, and the definition of $`b_i`$, this line thus contributes $`b_i(w_2+w_4+\mathrm{}+w_{i1})`$ to the weight $`\stackrel{~}{wt}(h)`$ of $`h`$. Similarly, for $`i`$ even, the $`i`$th line is in the SE direction and contributes $`b_i(w_1+w_3+\mathrm{}+w_{i1})`$ to $`\stackrel{~}{wt}(h)`$. The lemma then follows for $`d=0`$. The case $`d=1`$ is similar. $`\mathrm{}`$
### 2.2 Path parameters
We make the following definitions:
$$\begin{array}{cc}\alpha _{a,b}^{p,p^{}}\hfill & =ba;\hfill \\ \beta _{a,b,e,f}^{p,p^{}}\hfill & =\frac{bp}{p^{}}\frac{ap}{p^{}}+fe;\hfill \\ \delta _{a,e}^{p,p^{}}\hfill & =\{\begin{array}{cc}0\hfill & \text{if }\frac{(a+(1)^e)p}{p^{}}=\frac{ap}{p^{}};\hfill \\ 1\hfill & \text{if }\frac{(a+(1)^e)p}{p^{}}\frac{ap}{p^{}}.\hfill \end{array}\hfill \end{array}$$
(The superscripts of $`\alpha _{a,b}^{p,p^{}}`$ are superfluous, of course.) It may be seen that the value of $`\delta _{a,e}^{p,p^{}}`$ gives the parity of the band in which the path pre-segment resides.
###### Lemma 2.2
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$. Then $`\alpha (h)=\alpha _{a,b}^{p,p^{}}`$ and $`\beta (h)=\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: That $`\alpha (h)=\alpha _{a,b}^{p,p^{}}`$ follows immediately from the definitions.
The second result is proved by induction on $`L`$. If $`h𝒫_{a,b,e,f}^{p,p^{}}(0)`$ then $`a=b`$, whence $`\beta _{a,b,e,f}^{p,p^{}}=fe=\beta (h)`$, immediately from the definitions.
For $`L>0`$, let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ and assume that the result holds for all $`h^{}𝒫_{a,b^{},e,f^{}}^{p,p^{}}(L1)`$. We consider a particular $`h^{}`$ by setting $`h_i^{}=h_i`$ for $`0i<L`$, $`b^{}=h_{L1}`$ and choosing $`f^{}\{0,1\}`$ so that $`f^{}=0`$ if either $`bb^{}=1`$ and the $`L`$th segment of $`h`$ lies in an even band, or $`bb^{}=1`$ and the $`L`$th segment of $`h`$ lies in an odd band; and $`f^{}=1`$ otherwise. It may easily be checked that the $`(L1)`$th vertex of $`h^{}`$ is scoring if and only if the $`(L1)`$th vertex of $`h`$ is scoring. Then, from the definition of $`\beta (h)`$, we see that:
$$\beta (h)=\{\begin{array}{cc}\beta (h^{})+1\hfill & \text{if }bb^{}=1\text{ and }f=1;\hfill \\ \beta (h^{})1\hfill & \text{if }bb^{}=1\text{ and }f=0;\hfill \\ \beta (h^{})\hfill & \text{otherwise.}\hfill \end{array}$$
The induction hypothesis gives $`\beta (h^{})=b^{}p/p^{}ap/p^{}+f^{}e`$. Then when the $`L`$th segment of $`h`$ lies in an even band so that $`bp/p^{}=b^{}p/p^{}`$, consideration of the four cases of $`bb^{}=\pm 1`$ and $`f\{0,1\}`$ shows that $`\beta (h)=bp/p^{}ap/p^{}+fe`$. When the $`L`$th segment of $`h`$ lies in an odd band so that $`bp/p^{}=b^{}p/p^{}+bb^{}`$, consideration of the four cases of $`bb^{}=\pm 1`$ and $`f\{0,1\}`$ again shows that $`\beta (h)=bp/p^{}ap/p^{}+fe`$. The result follows by induction. $`\mathrm{}`$
### 2.3 Scoring generating functions
We now define a generating function for paths that have a particular number of non-scoring vertices. First define $`𝒫_{a,b,e,f}^{p,p^{}}(L,m)`$ to be the subset of $`𝒫_{a,b,e,f}^{p,p^{}}(L)`$ comprising those paths $`h`$ for which $`m(h)=m`$. Then define:
$$\chi _{a,b,e,f}^{p,p^{}}(L,m;q)=\underset{h𝒫_{a,b,e,f}^{p,p^{}}(L,m)}{}q^{\stackrel{~}{wt}(h)}.$$
(11)
###### Lemma 2.3
Let $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$. Then
$$\chi _{a,b,e,f}^{p,p^{}}(L)=\underset{\genfrac{}{}{0pt}{}{mL+\beta }{\left(\text{mod}\mathrm{\hspace{0.17em}2}\right)}}{}\chi _{a,b,e,f}^{p,p^{}}(L,m).$$
Proof: Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$. We claim that $`m(h)+L(h)+\beta (h)0(\text{mod}\mathrm{\hspace{0.17em}2})`$. This will follow from showing that $`L(h)m(h)+(1)^{d(h)}\beta (h)`$ is even. If $`h`$ has striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)},`$ then $`L(h)m(h)=(b_1+b_2+\mathrm{}+b_l)(e+d+\pi )\text{mod}\mathrm{\hspace{0.17em}2}`$, where $`\pi =\pi (h)`$. For $`e+d+\pi 0(\text{mod}\mathrm{\hspace{0.17em}2})`$, we immediately obtain $`L(h)m(h)+(1)^d\beta (h)=2(b_1+b_3+\mathrm{})`$. For $`e+d+\pi 0(\text{mod}\mathrm{\hspace{0.17em}2})`$, we obtain $`L(h)m(h)+(1)^d\beta (h)=2(b_1+b_3+\mathrm{})1+(1)^{d+e}`$, whence the claim is proved in all cases. The lemma then follows, once it is noted, via Lemma 2.2, that $`\beta (h)=\beta _{a,b,e,f}^{p,p^{}}`$. $`\mathrm{}`$
###### Note 2.4
Since each element of $`𝒫_{a,b,e,f}^{p,p^{}}(L,m)`$ has $`L+1`$ vertices, it follows that $`\chi _{a,b,e,f}^{p,p^{}}(L,m)`$ is non-zero only if $`0mL+1`$. Therefore the sum in Lemma 2.3 may be further restricted to $`0mL+1`$.
### 2.4 A seed
The following result provides a seed on which the results of later sections will act.
###### Lemma 2.5
If $`L0`$ is even then:
$$\chi _{1,1,0,0}^{1,3}(L,m)=\chi _{2,2,1,1}^{1,3}(L,m)=\delta _{m,0}q^{\frac{1}{4}L^2}.$$
If $`L>0`$ is odd then:
$$\chi _{1,2,0,1}^{1,3}(L,m)=\chi _{2,1,1,0}^{1,3}(L,m)=\delta _{m,0}q^{\frac{1}{4}(L^21)}.$$
Proof: The $`(1,3)`$-model comprises one even band. Thus when $`L`$ is even, there is precisely one $`h𝒫_{1,1,0,0}^{1,3}(L)`$. It has $`h_i=1`$ for $`i`$ even, and $`h_i=2`$ for $`i`$ odd. We see that $`h`$ has striking sequence $`\left(\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{0}{1}\right)^{(0,0,0)}`$ and $`m(h)=0`$. Lemma 2.1 then yields $`\stackrel{~}{wt}(h)=0+1+1+2+2+3+\mathrm{}+(\frac{1}{2}L1)+\frac{1}{2}L=(L/2)^2`$, as required.
The other expressions follow in a similar way. $`\mathrm{}`$
### 2.5 Partitions
A partition $`\lambda =(\lambda _1,\lambda _2,\mathrm{},\lambda _k)`$ is a sequence of $`k`$ integer parts $`\lambda _1,\lambda _2,\mathrm{},\lambda _k,`$ satisfying $`\lambda _1\lambda _2\mathrm{}\lambda _k>0`$. It is to be understood that $`\lambda _i=0`$ for $`i>k`$. The weight $`\mathrm{wt}(\lambda )`$ of $`\lambda `$ is given by $`\mathrm{wt}(\lambda )=_{i=1}^k\lambda _i`$.
We define $`𝒴(k,m)`$ to be the set of all partitions $`\lambda `$ with at most $`k`$ parts, and for which $`\lambda _1m`$. A proof of the following well known result may be found in .
###### Lemma 2.6
The generating function,
$$\underset{\lambda 𝒴(k,m)}{}q^{\mathrm{wt}(\lambda )}=\left[\genfrac{}{}{0pt}{}{m+k}{m}\right]_q.$$
## 3 The $``$-transform
In this section, we introduce the $``$-transform which maps paths $`𝒫_{a,b,e,f}^{p,p^{}}(L)`$ into $`𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{})`$ for certain $`a^{}`$, $`b^{}`$ and various $`L^{}`$.
The band structure of the $`(p,p^{}+p)`$-model is easily obtained from that of the $`(p,p^{})`$-model. Indeed, according to Section 1.1, for $`1r<p`$, the $`r`$th odd band of the $`(p,p^{}+p)`$-model lies between heights $`r(p^{}+p)/p=rp^{}/p+r`$ and $`r(p^{}+p)/p+1=rp^{}/p+r+1`$. Thus the height of the $`r`$th odd band in the $`(p,p^{}+p)`$-model is $`r`$ greater than that in the $`(p,p^{})`$-model. Therefore, the $`(p,p^{}+p)`$-model may be obtained from the $`(p,p^{})`$-model by increasing the distance between neighbouring odd bands by one unit and appending an extra even band to both the top and the bottom of the grid. For example, compare the $`(3,8)`$-model of Fig. 1 with the $`(3,11)`$-model of Fig. 3.
The $``$-transform has three components, which we refer to as path-dilation, particle-insertion, and particle-motion. These three components will also be known as the $`_1`$-, $`_2`$\- and $`_3`$-transforms respectively. In fact, particle-insertion is dependent on a parameter $`k\mathrm{}_0`$, and particle-motion is dependent on a partition $`\lambda `$ that has certain restrictions. Consequently, we sometimes refer to particle-insertion and particle-motion as $`_2(k)`$\- and $`_3(\lambda )`$-transforms respectively. Then, combining the $`_1`$-, $`_2(k)`$\- and $`_3(\lambda )`$-transforms produces the $`(k,\lambda )`$-transform.
### 3.1 Path-dilation
The $`_1`$-transform acts on a path $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ to yield a path $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(0)})`$, for certain $`a^{}`$, $`b^{}`$ and $`L^{(0)}`$. First, the starting point $`a^{}`$ of the new path $`h^{(0)}`$ is specified to be:
$$a^{}=a+\frac{ap}{p^{}}+e.$$
If $`r=ap/p^{}`$ then $`r`$ is the number of odd bands below $`h=a`$ in the $`(p,p^{})`$-model. Since the height of the $`r`$th odd band in the $`(p,p^{}+p)`$-model is $`r`$ greater than that in the $`(p,p^{})`$-model, we thus see that under path-dilation, the height of the startpoint above the next lowermost odd band (or if there isn’t one, the bottom of the grid) has either increased by one or remained constant.
We define $`d(h^{(0)})=d(h)`$. The above definition specifies that $`e(h^{(0)})=e(h)`$ and $`f(h^{(0)})=f(h)`$.
In the case that $`L=0`$ and $`e=f`$, we specify $`h^{(0)}`$ by setting $`L^{(0)}=L(h^{(0)})=0`$. When $`L=0`$ and $`ef`$, we leave the action of the $`_1`$-transform on $`h`$ undefined (it will not be used in this case). Thus in Lemmas 3.3, 3.6, 3.7, 3.10, 3.13, 4.4, 4.5 and and Corollary 3.4, we implicitly exclude consideration of the case $`L=0`$ and $`ef`$. However, it must be considered in the proofs of Corollaries 3.14 and 4.6.
In the case $`L>0`$ consider, as in Section 2.1, $`h`$ to comprise $`l`$ straight lines that alternate in direction, the $`i`$th of which is of length $`w_i`$ and possesses $`b_i`$ scoring vertices. $`h^{(0)}`$ is then defined to comprise $`l`$ straight lines that alternate in direction (since $`d(h^{(0)})=d(h)`$, the direction of the first line in $`h^{(0)}`$ is the same as that in $`h`$), the $`i`$th of which has length
$$w_i^{}=\{\begin{array}{c}w_i+b_i\text{if }i2\text{ or }e(h)+d(h)+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ w_1+b_1+2\pi (h)1\hfill \\ \text{if }i=1\text{ and }e(h)+d(h)+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2}).\hfill \end{array}$$
In particular, this determines $`L^{(0)}=L(h^{(0)})`$ and $`b^{}=h_{L^{(0)}}^{(0)}`$.
As an example, consider the path $`h`$ shown in Fig. 1 as an element of $`𝒫_{2,4,e,1}^{3,8}(14)`$. Here $`d(h)=0`$, $`\pi (h)=1`$ and $`ap/p^{}=0`$.
Thus when $`e=0`$, the action of path-dilation on $`h`$ produces the path given in Fig. 3.
This path is an element of $`𝒫_{2,6,e,1}^{3,11}(22)`$.
When $`e=1`$, the action of path-dilation on $`h`$ produces the element of $`𝒫_{3,6,e,1}^{3,11}(21)`$ given in Fig. 4.
The situation at the start point may be considered as falling into one of eight cases, corresponding to $`e(h),d(h),\pi (h)\{0,1\}`$.<sup>10</sup><sup>10</sup>10Theses cases may be seen to correspond to the eight cases of vertex type as listed in Table 1. In Table 2, we illustrate the four cases that arise when $`d(h)=0`$ (the four cases for $`d(h)=1`$ may be obtained from these by an up-down reflection and changing the value of $`e(h)`$).<sup>11</sup><sup>11</sup>11The examples here are such that $`w_13`$.
###### Lemma 3.1
Let $`1p<p^{}`$, $`1a<p^{}`$, $`e\{0,1\}`$ and $`a^{}=a+ap/p^{}+e`$. Then $`a^{}p/(p^{}+p)=ap/p^{}`$ and $`\delta _{a^{},e}^{p,p^{}+p}=0`$.
Proof: Let $`r=ap/p^{}`$ whence $`p^{}rpa<p^{}(r+1)`$. Then, for $`x\{0,1\}`$, we have $`(p^{}+p)rp(a+r+x)<(p^{}+p)r+p^{}+xp`$, so that $`(a+r+x)p/(p^{}+p)=r`$. In particular, $`a^{}p/(p^{}+p)=r`$, and $`(a+r+e+(1)^e)p/(p^{}+p)=r`$. Thus $`r=a^{}p/(p^{}+p)=(a^{}+(1)^e)p/(p^{}+p)`$ which gives the required results. $`\mathrm{}`$
This result asserts, amongst other things, that the pre-segment of $`h^{(0)}`$ always lies in an even band. This is also evident from Table 2.
###### Note 3.2
The action of path-dilation on $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ yields a path $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(0)})`$ that has, including the vertex at $`i=0`$, no adjacent scoring vertices, except in the case where $`\pi (h)=1`$ and $`e(h)=d(h)`$, when a single pair of scoring vertices occurs in $`h^{(0)}`$ at $`i=0`$ and $`i=1`$.
Also note that $`\pi (h^{(0)})=\pi (h)`$ unless $`\pi (h)=1`$ and $`e(h)=d(h)`$, in which case $`\pi (h^{(0)})=0`$.
Now compare the $`i`$th line of $`h^{(0)}`$ (which has length $`w_i^{}`$) with the $`i`$th line of $`h`$ (which has length $`w_i`$). Now for the sake of the following argument, assume that there are odd bands immediately below (i.e. between heights $`0`$ and $`1`$), and immediately above (i.e. between heights $`p^{}1`$ and $`p^{}`$) the $`(p,p^{})`$-model and do likewise for the $`(p,p^{}+p)`$-model.
If the lines in question are in the NE direction, we claim that the height of the final vertex of that in $`h^{(0)}`$ above the next lower odd band is one greater than that in $`h`$. If the lines in question are in the SE direction, we claim that the height of the final vertex of that in $`h^{(0)}`$ below the next higher odd band is one greater than that in $`h`$. In particular, if either the first or last segment of the $`i`$th line is in an odd band, then the corresponding segment of $`h^{(0)}`$ lies in the same odd band.
We also claim that if that of $`h`$ has a straight vertex that passes into the $`k`$th odd band in the $`(p,p^{})`$-model then that of $`h^{(0)}`$ has a straight vertex that passes into the $`k`$th odd band in the $`(p,p^{}+p)`$-model.
These claims follow because in passing from the $`(p,p^{})`$-model to the $`(p,p^{}+p)`$-model, the distance between neighbouring odd bands has increased by one, and because the length of each line has increased by one for every scoring vertex and possibly a small adjustment made to the length of the first line. In effect, a new straight vertex has been inserted immediately prior to each scoring vertex and, if $`e(h)+d(h)+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2})`$, adjusting the length of the resulting first line by $`2\pi (h)1`$.
###### Lemma 3.3
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ have striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)},`$ and let $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(0)})`$ be obtained from the action of the $`_1`$-transform on $`h`$. If $`e(h)+d(h)+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2})`$ then $`h^{(0)}`$ has striking sequence:
$$\left(\begin{array}{ccccc}a_1+b_1& a_2+b_2& a_3+b_3& \mathrm{}& a_l+b_l\\ b_1& b_2& b_3& \mathrm{}& b_l\end{array}\right)^{(e,f,d)},$$
and if $`e(h)+d(h)+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2})`$ then $`h^{(0)}`$ has striking sequence:
$$\left(\begin{array}{ccccc}a_1+b_1+\pi 1& a_2+b_2& a_3+b_3& \mathrm{}& a_l+b_l\\ b_1+\pi & b_2& b_3& \mathrm{}& b_l\end{array}\right)^{(e,f,d)}.$$
Moreover, if $`m=m(h)`$:
* $`m(h^{(0)})=L`$;
* $`L^{(0)}=\{\begin{array}{cc}2Lm+2\hfill & \text{if }\pi =1\text{ and }e=d,\hfill \\ 2Lm\hfill & \text{otherwise};\hfill \end{array}`$
* $`\alpha (h^{(0)})=\alpha (h)+\beta (h)`$;
* $`\beta (h^{(0)})=\beta (h)`$.
Proof: The form of the striking sequence for $`h^{(0)}`$ follows because, for $`i>1`$, every scoring vertex in the $`i`$th line of $`h`$ accounts for an extra non-scoring vertex in that line. The same is true when $`i=1`$, except in the case $`(e(h)+d(h)+\pi (h))1`$ (throughout this paper, in proofs, we take all equivalences, modulo 2.) when the length of the new $`1`$st line becomes $`a_1+2b_1+2\pi 1`$. That there are $`b_1+\pi `$ scoring vertices in this case, follows from examining Table 2.
Let $`e=e(h)`$, $`d=d(h)`$, $`\pi =\pi (h)`$ and $`\pi ^{}=\pi (h^{(0)})`$. Then $`e(h^{(0)})=e`$ and $`d(h^{(0)})=d`$.
If $`(e+d+\pi )0`$ then $`(e+d+\pi ^{})0`$ by Note 3.2. Thereupon $`m^{(0)}=_{i=1}^l(a_i+b_i)=L`$. Additionally, $`L^{(0)}=_{i=1}^l(a_i+2b_i)=2L_{i=1}^la_i=2Lm`$. That $`\beta (h^{(0)})=\beta (h)`$ and $`\alpha (h^{(0)})=\alpha (h)+\beta (h)`$ both follow immediately in this case.
On the other hand, if $`(e+d+\pi )0`$ then $`\pi =0ed`$ and $`\pi =1e=d`$. In each instance, Note 3.2 implies that $`\pi ^{}=0`$. Thereupon, $`m^{(0)}=(e+d+\pi ^{})\text{mod}\mathrm{\hspace{0.17em}2}+\pi 1+_{i=1}^l(a_i+b_i)=_{i=1}^l(a_i+b_i)=L`$. Additionally, $`L^{(0)}=2\pi 1+_{i=1}^l(a_i+2b_i)=2L(1+_{i=1}^la_i)+2\pi =2Lm+2\pi `$. This is the required value. Now in this case, $`\beta (h)=(1)^d((b_1+b_3+\mathrm{})(b_2+b_4+\mathrm{}))+(1)^e`$. When $`\pi =0`$ so that $`(e+d+\pi ^{})1`$ then $`\beta (h^{(0)})=\beta (h)`$ follows immediately. When $`\pi =1`$, we have $`\beta (h^{(0)})=(1)^d((b_1+1+b_3+\mathrm{})(b_2+b_4+\mathrm{}))`$. $`\beta (h^{(0)})=\beta (h)`$ now follows in this case because $`(e+d+\pi )0`$ implies that $`e=d`$. Finally, $`\alpha (h^{(0)})=\alpha (h)+(1)^d((b_1+b_3+\mathrm{})(b_2+b_4+\mathrm{}))+(1)^d(2\pi 1)`$. Since $`(1)^d(2\pi 1)=(1)^d(1)^\pi =(1)^e`$, the lemma then follows. $`\mathrm{}`$
###### Corollary 3.4
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ and $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(0)})`$ be the path obtained by the action of the $`_1`$-transform on $`h`$. Then $`a^{}=a+ap/p^{}+e`$ and $`b^{}=b+bp/p^{}+f`$.
Proof: $`a^{}=a+ap/p^{}+e`$ is by definition. Lemma 3.3 gives $`\alpha (h^{(0)})=\alpha (h)+\beta (h)`$, whence Lemma 2.2 implies that $`\alpha _{a^{},b^{}}^{p,p^{}+p}=\alpha _{a,b}^{p,p^{}}+\beta _{a,b,e,f}^{p,p^{}}`$. Expanding this gives $`b^{}a^{}=ba+bp/p^{}ap/p^{}+fe`$, whence $`b^{}=b+bp/p^{}+f`$. $`\mathrm{}`$
The above result implies that the $`_1`$-transform maps $`𝒫_{a,b,e,f}^{p,p^{}}(L)`$ into a set of paths that have the same startpoint as one another and the same endpoint as one another. However, the lengths of these paths are not necessarily equal. We also see that the transformation of the endpoint is analogous to that which occurs at the startpoint. In particular, Lemma 3.1 implies that $`\delta _{b^{},f}^{p,p^{}+p}=0`$ so that the path post-segment of $`h^{(0)}`$ always resides in an even band. For the four cases where $`h_L=h_{L1}1`$, the $`_1`$-transform affects the endpoint as in Table 3 (the value $`\pi ^{}(h)`$ is the parity of the band in which the $`L`$th segment of $`h`$ lies).
###### Lemma 3.5
Let $`1p<p^{}`$, $`1a,b<p^{}`$, $`e,f\{0,1\}`$, $`a^{}=a+ap/p^{}+e`$, and $`b^{}=b+bp/p^{}+f`$. Then $`\alpha _{a^{},b^{}}^{p,p^{}+p}=\alpha _{a,b}^{p,p^{}}+\beta _{a,b,e,f}^{p,p^{}}`$ and $`\beta _{a^{},b^{},e,f}^{p,p^{}+p}=\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: Lemma 3.1 implies that $`a^{}p/(p^{}+p)=ap/p^{}`$, $`b^{}p/(p^{}+p)=bp/p^{}`$. The results then follow immediately from the definitions. $`\mathrm{}`$
###### Lemma 3.6
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ and $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(0)})`$ be the path obtained by the action of the $`_1`$-transform on $`h`$. Then
$$\stackrel{~}{wt}(h^{(0)})=\stackrel{~}{wt}(h)+\frac{1}{4}\left((L^{(0)}m^{(0)})^2\beta ^2\right),$$
where $`m^{(0)}=m(h^{(0)})`$ and $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: Let $`h`$ have striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)}`$, and let $`\pi =\pi (h)`$. If $`(e+d+\pi )0(\text{mod}\mathrm{\hspace{0.17em}2})`$, then Lemmas 3.3 and 2.1 show that
$$\stackrel{~}{wt}(h^{(0)})\stackrel{~}{wt}(h)=(b_1+b_3+b_5+\mathrm{})(b_2+b_4+b_6+\mathrm{}).$$
Via Lemma 3.3, we obtain $`L^{(0)}m^{(0)}=Lm(h)=b_1+b_2+\mathrm{}+b_l`$. Then since $`\beta (h)=\pm ((b_1+b_3+b_5+\mathrm{})(b_2+b_4+b_6+\mathrm{}))`$, it follows that
$$\stackrel{~}{wt}(h^{(0)})\stackrel{~}{wt}(h)=\frac{1}{4}((L^{(0)}m^{(0)})^2\beta (h)^2).$$
If $`(e+d+\pi )0(\text{mod}\mathrm{\hspace{0.17em}2})`$, then Lemmas 3.3 and 2.1 show that
$`\stackrel{~}{wt}(h^{(0)})\stackrel{~}{wt}(h)`$ $`=`$ $`(2\pi 1+b_1+b_3+b_5+\mathrm{})(b_2+b_4+b_6+\mathrm{})`$
$`=`$ $`{\displaystyle \frac{1}{4}}((L^{(0)}m^{(0)})^2\beta (h)^2),`$
the second equality resulting because $`L^{(0)}m^{(0)}=Lm(h)+2\pi =b_1+b_2+\mathrm{}+b_l+2\pi 1`$ and
$`\beta (h)`$ $`=`$ $`(1)^d((b_1+b_3+b_5+\mathrm{})(b_2+b_4+b_6+\mathrm{}))+(1)^e`$
$`=`$ $`\pm ((2\pi 1+b_1+b_3+b_5+\mathrm{})(b_2+b_4+b_6+\mathrm{})),`$
on using $`(1)^{e+d}=(1)^\pi =2\pi 1`$.
Finally, Lemma 2.2 gives $`\beta (h)=\beta _{a,b,e,f}^{p,p^{}}=\beta `$. $`\mathrm{}`$
### 3.2 Particle insertion
Let $`p^{}>2p`$ so that the $`(p,p^{})`$-model has no two neighbouring odd bands, and let $`\delta _{a^{},e}^{p,p^{}}=0`$. Then if $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}}(L^{(0)})`$, the pre-segment of $`h^{(0)}`$ lies in an even band. By inserting a particle into $`h^{(0)}`$, we mean displacing $`h^{(0)}`$ two positions to the right and inserting two segments: the leftmost of these is in the NE (resp. SE) direction if $`e=0`$ (resp. $`e=1`$), and the rightmost is in the opposite direction, which is thus the direction of the pre-segment of $`h^{(0)}`$. In this way, we obtain a path $`h^{(1)}`$ of length $`L^{(0)}+2`$. We assign $`e(h^{(1)})=e`$ and $`f(h^{(1)})=f`$. Note also that $`d(h^{(1)})=e`$ and $`\pi (h^{(1)})=0`$.
Thereupon, we may repeat this process of particle insertion. After inserting $`k`$ particles into $`h^{(0)}`$, we obtain a path $`h^{(k)}𝒫_{a^{},b^{},e,f}^{p,p^{}}(L^{(0)}+2k)`$. We say that $`h^{(k)}`$ has been obtained by the action of a $`_2(k)`$-transform on $`h^{(0)}`$.
In the case of the element of $`𝒫_{3,6,1,1}^{3,11}(21)`$ shown in Fig. 4, the insertion of two particles produces the element of $`𝒫_{3,6,1,1}^{3,11}(25)`$ shown in Fig. 5.
###### Lemma 3.7
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$. Apply a $`_1`$-transform to $`h`$ to obtain the path $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(0)})`$. Then obtain $`h^{(k)}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{(k)})`$ by applying a $`_2(k)`$-transform to $`h^{(0)}`$. If $`m^{(k)}=m(h^{(k)})`$, then $`L^{(k)}=L^{(0)}+2k`$, $`m^{(k)}=m^{(0)}`$ and
$$\stackrel{~}{wt}(h^{(k)})=\stackrel{~}{wt}(h)+\frac{1}{4}((L^{(k)}m^{(k)})^2\beta ^2),$$
(12)
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: That $`L^{(k)}=L^{(0)}+2k`$ follows immediately from the definition of a $`_2`$-transform. Lemma 3.6 yields:
$$\stackrel{~}{wt}(h^{(0)})=\stackrel{~}{wt}(h)+\frac{1}{4}\left((L^{(0)}m(h^{(0)}))^2\beta ^2\right).$$
Let the striking sequence of $`h^{(0)}`$ be $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)},`$ and let $`\pi =\pi (h^{(0)})`$.
If $`e=d`$, we are restricted to the case $`\pi =0`$, since $`\delta _{a^{},e}^{p,p^{}+p}=0`$ by Lemma 3.1. The striking sequence of $`h^{(1)}`$ is then $`\left(\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,e)}`$. Thereupon $`m(h^{(1)})=_{i=1}^la_i=m(h^{(0)})`$. In this case, Lemma 2.1 shows that $`\stackrel{~}{wt}(h^{(1)})\stackrel{~}{wt}(h^{(0)})=1+b_1+b_2+\mathrm{}+b_l=L^{(0)}m^{(0)}+1`$.
If $`ed`$, the striking sequence of $`h^{(1)}`$ is $`\left(\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{a_1+1\pi }{b_1+\pi }\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,e)}`$. Then $`m(h^{(1)})=1\pi +_{i=1}^la_i`$ which equals $`m(h^{(0)})=(e+d+\pi )\text{mod}\mathrm{\hspace{0.17em}2}+_{i=1}^la_i`$ for both $`\pi =0`$ and $`\pi =1`$. Here, Lemma 2.1 shows that $`\stackrel{~}{wt}(h^{(1)})\stackrel{~}{wt}(h^{(0)})=\pi +b_1+b_2+\mathrm{}+b_l`$. Since $`L^{(0)}m^{(0)}=(e+d+\pi )\text{mod}\mathrm{\hspace{0.17em}2}+b_1+b_2+\mathrm{}+b_l`$, we once more have $`\stackrel{~}{wt}(h^{(1)})\stackrel{~}{wt}(h^{(0)})=L^{(0)}m^{(0)}+1`$.
Repeated application of these results, yields $`m(h^{(k)})=m(h^{(0)})`$ and
$$\stackrel{~}{wt}(h^{(k)})=\stackrel{~}{wt}(h^{(0)})+k\left(L^{(0)}m(h^{(0)})\right)+k^2.$$
Then, on using (12) and $`L^{(k)}=L^{(0)}+2k`$, the lemma follows. $`\mathrm{}`$
### 3.3 Particle moves
In this section, we once more restrict to the case $`p^{}>2p`$ so that the $`(p,p^{})`$-model has no two neighbouring odd bands, and consider only paths $`h𝒫_{a^{},b^{},e,f}^{p,p^{}}(L^{})`$, where $`\delta _{a^{},e}^{p,p^{}}=\delta _{b^{},f}^{p,p^{}}=0`$.
We specify six types of local deformations of a path. These deformations will be known as particle moves. In each of the six cases, a particular sequence of four segments of a path is changed to a different sequence, the remainder of the path being unchanged. The moves are as follows — the path portion to the left of the arrow is changed to that on the right:
Move 1.
Move 2.
Move 3.
Move 4.
Move 5.
Move 6.
Since $`p^{}>2p`$, each odd band is straddled by a pair of even bands. Thus, there is no impediment to enacting moves 2 and 5 for paths in $`𝒫_{a,b,e,f}^{p,p^{}}(L)`$.
Note that moves 4–6 are inversions of moves 1–3. Also note that moves 2 and 3 (likewise moves 5 and 6) may be considered to be the same move since in the two cases, the same sequence of three edges is changed.
In addition to the six moves described above, we permit certain deformations of a path close to its left and right extremities in certain circumstances. Each of these moves will be referred to as an edge-move. They, together with their validity, are as follows:
If $`e=1`$:
Edge-move 1.
If $`e=0`$:
Edge-move 2.
If $`f=0`$:
Edge-move 3.
If $`f=1`$:
Edge-move 4.
In fact, the above four edge-moves may be considered as instances of moves 1 and 4 described beforehand, if for edge-moves 1 and 2, we append the appropriate pre-segment to the path, and for edge-moves 3 and 4, we append the appropriate post-segment to the path.
###### Lemma 3.8
Let the path $`\widehat{h}`$ differ from the path $`h`$ in that four consecutive segments have changed according to one of the six moves described above, or in that three consecutive segments have changed according to one of the four edge-moves described above (subject to their restrictions). Then
$$\stackrel{~}{wt}(\widehat{h})=\stackrel{~}{wt}(h)+1.$$
Additionally, $`L(\widehat{h})=L(h)`$ and $`m(\widehat{h})=m(h)`$.
Proof: For each of the six moves and four edge-moves, take the $`(x,y)`$-coordinate of the leftmost point of the depicted portion of $`h`$ to be $`(x_0,y_0)`$. Now consider the contribution to the weight of the three vertices in question before and after the move (although the vertex at $`(x_0,y_0)`$ may change, its contribution doesn’t). In each of the ten cases, the contribution is $`x_0+y_0+1`$ before the move and $`x_0+y_0+2`$ afterwards. Thus $`\stackrel{~}{wt}(\widehat{h})=\stackrel{~}{wt}(h)+1`$. The other statements are immediate on inspecting all ten moves. $`\mathrm{}`$
Now observe that for each of the ten moves specified above, the sequence of path segments before the move consists of an adjacent pair of scoring vertices followed by a non-scoring vertex. The specified move replaces this combination with a non-scoring vertex followed by two scoring vertices. As anticipated above, the pair of adjacent scoring vertices is viewed as a particle. Thus each of the above ten moves describes a particle moving to the right by one step.
When $`p^{}>2p`$, so that there are no two adjacent odd bands in the $`(p,p^{})`$-model, and noting that $`\delta _{b^{},f}^{p,p^{}}=0`$, we see that each sequence comprising two scoring vertices followed by a non-scoring vertex is present amongst the ten configurations prior to a move, except for the case depicted in Fig. 6 and its up-down reflection.
Only in these cases, where the 0th and 1st segments are scoring and the first two segments are in the same direction, do we not refer to the adjacent pair of scoring vertices as a particle.
Also note that when $`p^{}>2p`$ and $`\delta _{a^{},e}^{p,p^{}}=\delta _{b^{},f}^{p,p^{}}=0`$, each sequence of a non-scoring vertex followed by two scoring vertices appears amongst the ten configurations that result from a move. In such cases, the move may thus be reversed.
### 3.4 The $`_3`$-transform
Since in each of the moves described in Section 3.3, a pair of scoring vertices shifts to the right by one step, we see that a succession of such moves is possible until the pair is followed by another scoring vertex. If this itself is followed by yet another scoring vertex, we forbid further movement. However, if it is followed by a non-scoring vertex, further movement is allowed after considering the latter two of the three consecutive scoring vertices to be the particle (instead of the first two).
As in Section 3.2, let $`h^{(k)}`$ be a path resulting from a $`_2(k)`$-transform acting on a path that itself is the image of a $`_1`$ transform. We now consider moving the $`k`$ particles that have been inserted.
###### Lemma 3.9
Let $`\delta _{b^{},f}^{p,p^{}}=0`$. There is a bijection between the set of paths obtained by moving the particles in $`h^{(k)}`$ and $`𝒴(k,m)`$, where $`m=m(h^{(k)})`$. This bijection is such that if $`\lambda 𝒴(k,m)`$ is the bijective image of a particular $`h`$ then
$$\stackrel{~}{wt}(h)=\stackrel{~}{wt}(h^{(k)})+\mathrm{wt}(\lambda ).$$
Additionally, $`L(h)=L(h^{(k)})`$ and $`m(h)=m(h^{(k)})`$.
Proof: Since each particle moves by traversing a non-scoring vertex, and there are $`m`$ of these to the right of the rightmost particle in $`h^{(k)}`$, and there are no consecutive scoring vertices to its right, this particle can make $`\lambda _1`$ moves to the right, with $`0\lambda _1m`$. Similarly, the next rightmost particle can make $`\lambda _2`$ moves to the right with $`0\lambda _2\lambda _1`$. Here, the upper restriction arises because the two scoring vertices would then be adjacent to those of the first particle. Continuing in this way, we obtain that all possible final positions of the particles are indexed by $`\lambda =(\lambda _1,\lambda _2,\mathrm{},\lambda _k)`$ with $`m\lambda _1\lambda _2\mathrm{}\lambda _k0`$, that is, by partitions of at most $`k`$ parts with no part exceeding $`m`$. Moreover, since by Lemma 3.8 the weight increases by one for each move, the weight increase after the sequence of moves specified by a particular $`\lambda `$ is equal to $`\mathrm{wt}(\lambda )`$. The final statement also follows from Lemma 3.8. $`\mathrm{}`$
We say that a path obtained by moving the particles in $`h^{(k)}`$ according to the partition $`\lambda `$ has been obtained by the action of a $`_3(\lambda )`$-transform.
Having defined $`_1`$, $`_2(k)`$ for $`k0`$ and $`_3(\lambda )`$ for $`\lambda `$ a partition with at most $`k`$ parts, we now define a $`(k,\lambda )`$-transform as the composition $`(k,\lambda )=_3(\lambda )_2(k)_1`$.
###### Lemma 3.10
Let $`h^{}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{})`$ be obtained from $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ by the action of the $`(k,\lambda )`$-transform. If $`\pi =\pi (h)`$ and $`m=m(h)`$ then:
$$\begin{array}{c}L^{}=\{\begin{array}{cc}2Lm+2k+2\hfill & \text{if }\pi =1\text{ and }e=d,\hfill \\ 2Lm+2k\hfill & \text{otherwise};\hfill \end{array}\hfill \\ m(h^{})=L;\hfill \\ \stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+\frac{1}{4}\left((L^{}L)^2\beta ^2\right)+\mathrm{wt}(\lambda ),\hfill \end{array}$$
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: These results follow immediately from Lemmas 3.3, 3.7 and 3.9. $`\mathrm{}`$
###### Note 3.11
Since particle insertion and the particle moves don’t change the startpoint, endpoint or value $`e(h)`$ or $`f(h)`$ of a path $`h`$, then in view of Lemma 3.1 and Corollary 3.4, we see that the action of a $``$-transform on $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ yields a path $`h^{}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{})`$, where $`a^{}=a+ap/p^{}+e`$, $`b^{}=b+bp/p^{}+f`$, and $`\delta _{a^{},e}^{p,p^{}+p}=\delta _{b^{},f}^{p,p^{}+p}=0`$.
### 3.5 Particle content of a path
Again restrict to the case $`p^{}>2p`$ so that the $`(p,p^{})`$-model has no two neighbouring odd bands, and let $`h^{}𝒫_{a^{},b^{},e,f}^{p,p^{}}(L^{})`$. In the following lemma, we once more restrict to the cases for which $`\delta _{a,e}^{p,p^{}}=\delta _{b,f}^{p,p^{}}=0`$, and thus only consider the cases for which the pre-segment and the post-segment of $`h^{}`$ lie in even bands.
###### Lemma 3.12
For $`1p<p^{}`$ with $`p^{}>2p`$, let $`1a^{},b^{}<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a^{},e}^{p,p^{}}=\delta _{b^{},f}^{p,p^{}}=0`$. If $`h^{}𝒫_{a^{},b^{},e,f}^{p,p^{}}(L^{})`$, then there is a unique triple $`(h,k,\lambda )`$ where $`h𝒫_{a,b,e,f}^{p,p^{}p}(L)`$ for some $`a,b,L`$, such that the action of a $`(k,\lambda )`$-transform on $`h`$ results in $`h^{}`$.
Proof: This is proved by reversing the constructions described in the previous sections. Locate the leftmost pair of consecutive scoring vertices in $`h^{}`$, and move them leftward by reversing the particle moves, until they occupy the $`0`$th and $`1`$st positions. This is possible in all cases when $`\delta _{a^{},e}^{p,p^{}}=\delta _{b^{},f}^{p,p^{}}=0`$. Now ignoring these two vertices, do the same with the next leftmost pair of consecutive scoring vertices, moving them leftward until they occupy the third and fourth positions. Continue in this way until all consecutive scoring vertices occupy the leftmost positions of the path. Denote this path by $`h^{()}`$. At the leftmost end of $`h^{()}`$, there will be a number of even segments (possibly zero) alternating in direction. Let this number be $`2k`$ or $`2k+1`$ according to whether is it even or odd. Clearly $`h^{}`$ results from $`h^{()}`$ by a $`_3(\lambda )`$-transform for a particular $`\lambda `$ with at most $`k`$ parts.
Removing the first $`2k`$ segments of $`h^{()}`$ yields a path $`h^{(0)}𝒫_{a^{},b^{},e,f}^{p,p^{}}`$. This path thus has no two consecutive scoring vertices, except possibly at the $`0`$th and $`1`$st positions, and then only if the first vertex is a straight vertex (as in Fig. 6). Moreover, $`h^{(k)}`$ arises by the action of a $`_2(k)`$-transform on $`h^{(0)}`$.
Ignoring for the moment the case where there are scoring vertices at the $`0`$th and $`1`$st positions, $`h^{(0)}`$ has by construction no pair of consecutive scoring vertices. Therefore, beyond the $`0`$th vertex, we may remove a non-scoring vertex before every scoring vertex to obtain a path $`h𝒫_{a,b,e,f}^{p,p^{}p}(L)`$ for some $`a,b,L`$, from which $`h^{(0)}`$ arises by the action of a $`_1`$-transform.
On examining the third case depicted in Table 2, we see that the case where $`h^{(0)}`$ has a pair of scoring vertices at the $`0`$th and $`1`$st positions, arises similarly from a particular $`h𝒫_{a,b,e,f}^{p,p^{}p}(L)`$ for some $`a,b,L`$. The lemma is then proved. $`\mathrm{}`$
The value of $`k`$ obtained above will be referred to as the particle content of $`h^{}`$.
###### Lemma 3.13
For $`1p<p^{}`$, let $`1a,b<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a,e}^{p,p^{}}=0`$. Set $`a^{}=a+e+ap/p^{}`$ and $`b^{}=b+f+bp/p^{}`$. Fix $`m_0,m_10`$. Then the map $`(h,k,\lambda )h^{}`$ effected by the action of a $`(k,\lambda )`$-transform on $`h`$, is a bijection between $`_k𝒫_{a,b,e,f}^{p,p^{}}(m_1,2k+2m_1m_0)\times 𝒴(k,m_1)`$ and $`𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,m_1)`$. Moreover,
$$\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+\frac{1}{4}\left((m_0m_1)^2\beta ^2\right)+\mathrm{wt}(\lambda ),$$
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: Given $`h𝒫_{a,b,e,f}^{p,p^{}}(m_1,m)`$, let $`h^{}`$ be the result of a $`(k,\lambda )`$-transform on $`h`$.
Since $`\delta _{a,e}^{p,p^{}}=0`$ so that $`(a+(1)^e)p/p^{})=ap/p^{}`$, it follows that if $`\pi (h)=1`$ then $`e(h)d(h)`$. Then, with $`m=2m_1+2km_0`$, we obtain $`h^{}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,m_1)`$ via Lemma 3.10.
Lemma 3.1 shows that $`\delta _{a^{},e}^{p,p^{}+p}=\delta _{b^{},f}^{p,p^{}+p}=0`$. Thereupon, Lemma 3.12 shows that each $`h^{}𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,m_1)`$ arises from a unique triple $`(h,k,\lambda )`$, with $`h𝒫_{a,b,e,f}^{p,p^{}}(m_1,m)`$ for some $`m`$. The bijection then follows.
The expression for $`\stackrel{~}{wt}(h^{})`$ also results from Lemma 3.10. $`\mathrm{}`$
Note that the above lemma excludes consideration of the case for which $`\delta _{a,e}^{p,p^{}}=1`$. In fact, similar results fail in that case. Nonetheless, it is necessary to tackle the $`\delta _{a,e}^{p,p^{}}=1`$ case for a restricted set of paths in the more general analysis of .
###### Corollary 3.14
For $`1p<p^{}`$, let $`1a,b<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a,e}^{p,p^{}}=0`$. Set $`a^{}=a+e+ap/p^{}`$ and $`b^{}=b+f+bp/p^{}`$. Fix $`m_0,m_10`$. Then
$$\stackrel{~}{\chi }_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,m_1)=q^{\frac{1}{4}\left((m_0m_1)^2\beta ^2\right)}\underset{\genfrac{}{}{0pt}{}{mm_0}{\left(\text{mod}\mathrm{\hspace{0.17em}2}\right)}}{}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_0+m)}{m_1}\right]_q\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(m_1,m),$$
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: Apart from the case where $`m_1=0`$ and $`ef`$, this follows immediately from Lemma 3.13 on setting $`m=2m_1+2km_0`$, once it is noted, via Lemma 2.6, that $`\left[\genfrac{}{}{0pt}{}{k+m_1}{m_1}\right]_q`$ is the generating function for $`𝒴(k,m_1)`$.
For the case $`m_1=0`$ and $`ef`$, both sides are zero unless $`a=b`$ and $`m_0`$ is odd. In this case, $`𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,0)`$ has precisely one element $`h`$ for which (via the same calculation as in the proof of 2.5) $`\stackrel{~}{wt}(h)=\frac{1}{4}(m_0^21)`$. Thus the two sides are also equal in this case. $`\mathrm{}`$
## 4 The $`𝒟`$-transform
The $`𝒟`$-transform is defined to act on each $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ to yield a path $`\widehat{h}𝒫_{a,b,1e,1f}^{p^{}p,p^{}}(L)`$ with exactly the same sequence of integer heights, i.e., $`\widehat{h}_i=h_i`$ for $`0iL`$. Note that, by definition, $`e(\widehat{h})=1e(h)`$ and $`f(\widehat{h})=1f(h)`$.
Since the band structure of the $`(p^{}p,p^{})`$-model is obtained from that of the $`(p,p^{})`$-model simply by replacing odd bands by even bands and vice-versa, then, ignoring the vertex at $`i=0`$, each scoring vertex maps to a non-scoring vertex and vice-versa. That $`e(h)`$ and $`e(\widehat{h})`$ differ implies that the vertex at $`i=0`$ is both scoring or both non-scoring in $`h`$ and $`\widehat{h}`$.
###### Lemma 4.1
Let $`\widehat{h}𝒫_{a,b,1e,1f}^{p^{}p,p^{}}(L)`$ be obtained from $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ by the action of the $`𝒟`$-transform. Then $`\pi (\widehat{h})=1\pi (h)`$. Moreover, if $`m=m(h)`$ then:
$$\begin{array}{c}L(\widehat{h})=L;\hfill \\ m(\widehat{h})=\{\begin{array}{cc}Lm\hfill & \text{if }e+d+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2}),\hfill \\ Lm+2\hfill & \text{if }e+d+\pi (h)0(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \end{array}\hfill \\ \stackrel{~}{wt}(\widehat{h})=\frac{1}{4}\left(L^2\alpha (h)^2\right)\stackrel{~}{wt}(h).\hfill \end{array}$$
Proof: Let $`h`$ have striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)}`$. Since, beyond the zeroth vertex, the $`𝒟`$-transform exchanges scoring vertices for non-scoring vertices and vice-versa, it follows that the striking sequence for $`\widehat{h}`$ is $`\left(\genfrac{}{}{0pt}{}{b_1}{a_1}\genfrac{}{}{0pt}{}{b_2}{a_2}\genfrac{}{}{0pt}{}{b_3}{a_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{b_l}{a_l}\right)^{(1e,f,d)}`$. It is immediate that $`L(\widehat{h})=L`$, $`\pi (\widehat{h})=1\pi (h)`$, $`e(\widehat{h})=1e(h)`$ and $`d(\widehat{h})=d(h)`$. Then $`m(\widehat{h})=(e(\widehat{h})+d(\widehat{h})+\pi (\widehat{h}))\text{mod}\mathrm{\hspace{0.17em}2}+_{i=1}^lb_i=(e+d+\pi (h))\text{mod}\mathrm{\hspace{0.17em}2}+L_{i=1}^la_i=2((e+d+\pi (h))\text{mod}\mathrm{\hspace{0.17em}2})+Lm(h)`$.
Now let $`w_i=a_i+b_i`$ for $`1il`$. Then, using Lemma 2.1, we obtain
$`\stackrel{~}{wt}(h)+\stackrel{~}{wt}(\widehat{h})`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{l}{}}}b_i(w_{i1}+w_{i3}+\mathrm{}+w_{1+imod2})`$
$`+{\displaystyle \underset{i=1}{\overset{l}{}}}a_i(w_{i1}+w_{i3}+\mathrm{}+w_{1+imod2})`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{l}{}}}w_i(w_{i1}+w_{i3}+\mathrm{}+w_{1+imod2})`$
$`=`$ $`(w_1+w_3+w_5+\mathrm{})(w_2+w_4+w_6+\mathrm{}).`$
The lemma then follows because $`(w_1+w_3+w_5+\mathrm{})+(w_2+w_4+w_6+\mathrm{})=L`$ and $`(w_1+w_3+w_5+\mathrm{})(w_2+w_4+w_6+\mathrm{})=\pm \alpha (h)`$. $`\mathrm{}`$
###### Lemma 4.2
Let $`1p<p^{}`$ with $`p`$ co-prime to $`p^{}`$ and $`1a<p^{}`$. Then $`a(p^{}p)/p^{}=a1ap/p^{}`$.
If, in addition, $`a`$ is interfacial in the $`(p,p^{})`$-model and $`\delta _{a,e}^{p,p^{}}=0`$ then $`a`$ is interfacial in the $`(p^{}p,p^{})`$-model and $`\delta _{a,1e}^{p^{}p,p^{}}=0`$.
Proof: Since $`p`$ and $`p^{}`$ are co-prime, $`ap/p^{}<ap/p^{}`$. Hence $`ap/p^{}+a(p^{}p)/p^{}=a1`$.
Since the $`(p,p^{})`$-model differs from the $`(p^{}p,p^{})`$-model only in that corresponding bands are of the opposite parity, $`a`$ being interfacial in one model implies that it also is in the other. The final part then follows immediately. $`\mathrm{}`$
###### Corollary 4.3
If $`1p<p^{}`$ with $`p`$ co-prime to $`p^{}`$, $`1a,b<p^{}`$ and $`e,f\{0,1\}`$ then $`\alpha _{a,b}^{p^{}p,p^{}}=\alpha _{a,b}^{p,p^{}}`$ and $`\beta _{a,b,1e,1f}^{p^{}p,p^{}}+\beta _{a,b,e,f}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}`$.
Proof: Lemma 4.2 gives $`ap/p^{}+a(p^{}p)/p^{}=a1`$ and likewise, $`bp/p^{}+b(p^{}p)/p^{}=b1`$. The required results then follow immediately. $`\mathrm{}`$
### 4.1 The $`𝒟`$-pair
It will often be convenient to consider the combined action of a $`𝒟`$-transform followed immediately by a $``$-transform. Such a pair will naturally be referred to as a $`𝒟`$-transform and maps a path $`h𝒫_{a,b,e,f}^{p^{}p,p^{}}(L)`$ to a path $`h^{}𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(L^{})`$, where $`a^{},b^{},L^{}`$ are determined by our previous results.
In what follows, the $`𝒟`$-transform will always follow a $``$-transform. Thus we restrict consideration to where $`2(p^{}p)<p^{}`$.
###### Lemma 4.4
With $`p^{}<2p`$, let $`h𝒫_{a,b,e,f}^{p^{}p,p^{}}(L)`$. Let $`h^{}𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(L^{})`$ result from the action of a $`𝒟`$-transform on $`h`$, followed by a $`(k,\lambda )`$-transform. Then:
$$\begin{array}{c}L^{}=\{\begin{array}{cc}L+m(h)+2k2\hfill & \text{if }\pi (h)=1\text{ and }e=d(h),\hfill \\ L+m(h)+2k\hfill & \text{otherwise};\hfill \end{array}\hfill \\ m(h^{})=L;\hfill \\ \stackrel{~}{wt}(h^{})=\frac{1}{4}\left(L^2+(L^{}L)^2\alpha ^2\beta ^2\right)+\mathrm{wt}(\lambda )\stackrel{~}{wt}(h),\hfill \end{array}$$
where $`\alpha =\alpha _{a,b}^{p,p^{}}`$ and $`\beta =\beta _{a,b,1e,1f}^{p,p^{}}`$.
Proof: Let $`\widehat{h}`$ result from the action of the $`𝒟`$-transform on $`h`$, and let $`d=d(h)`$, $`\pi =\pi (h)`$, $`\widehat{e}=e(\widehat{h})`$ $`\widehat{d}=d(\widehat{h})`$, $`\widehat{\pi }=\pi (\widehat{h})`$. Then we immediately have $`\widehat{d}=d`$, $`\widehat{e}=1e`$, and $`\widehat{\pi }=1\pi `$.
In the case where $`\pi =0`$ and $`ed`$, we then have, using Lemmas 3.10 and 4.1, $`L^{}=2L(\widehat{h})m(\widehat{h})+2k+2=2L(Lm(h)+2)+2k+2=L+m(h)+2k`$.
In the case where $`\pi =1`$ and $`e=d`$, we then have, using Lemmas 3.10 and 4.1, $`L^{}=2L(\widehat{h})m(\widehat{h})+2k=2L(Lm(h)+2)+2k=L+m(h)+2k2`$.
In the other cases, $`e+d+\pi 0(\text{mod}\mathrm{\hspace{0.17em}2})`$ and so $`\widehat{e}+\widehat{d}+\widehat{\pi }0(\text{mod}\mathrm{\hspace{0.17em}2})`$. Lemmas 3.10 and 4.1 yield $`L^{}=2L(\widehat{h})m(\widehat{h})+2k=2L(Lm(h))+2k=L+m(h)+2k`$.
The expressions for $`m(h^{})`$ and $`\stackrel{~}{wt}(h^{})`$ also follow immediately from Lemmas 3.10 and 4.1. $`\mathrm{}`$
We now obtain analogues of Lemma 3.13 and Corollary 3.14 which combine the $`𝒟`$-transform with the $``$-transform. As above, we restrict to where $`p^{}<2p`$.
###### Lemma 4.5
For $`1p<p^{}<2p`$, let $`1a,b<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a,e}^{p^{}p,p^{}}=0`$. Set $`a^{}=a+1e+ap/p^{}`$ and $`b^{}=b+1f+bp/p^{}`$. Fix $`m_0,m_10`$. Then the map $`(h,k,\lambda )h^{}`$ effected by the action of a $`𝒟`$-transform on $`h`$ followed by a $`(k,\lambda )`$-transform, is a bijection between $`_k𝒫_{a,b,e,f}^{p^{}p,p^{}}(m_1,m_0m_12k)\times 𝒴(k,m_1)`$ and $`𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(m_0,m_1)`$. Moreover,
$$\stackrel{~}{wt}(h^{})=\frac{1}{4}\left(m_1^2+(m_0m_1)^2\alpha ^2\beta ^2\right)+\mathrm{wt}(\lambda )\stackrel{~}{wt}(h),$$
where $`\alpha =\alpha _{a,b}^{p,p^{}}`$ and $`\beta =\beta _{a,b,1e,1f}^{p,p^{}}`$.
Proof: Given $`h𝒫_{a,b,e,f}^{p^{}p,p^{}}(m_1,m)`$, let $`\widehat{h}`$ result from the action of a $`𝒟`$-transform on $`h`$, and let $`h^{}`$ be the result of a $`(k,\lambda )`$-transform on $`\widehat{h}`$.
Since $`\delta _{a,e}^{p^{}p,p^{}}=0`$ so that $`(a+(1)^e)(p^{}p)/p^{})=a(p^{}p)/p^{}`$, it follows that if $`\pi (h)=1`$ then $`e(h)d(h)`$. Then, for $`m=m_0m_12k`$, we obtain $`h^{}𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(m_0,m_1)`$ via Lemma 4.4.
Lemma 3.1 gives $`\delta _{a^{},1e}^{p,p^{}+p}=\delta _{b^{},1f}^{p,p^{}+p}=0`$. Lemma 3.12 then shows that for arbitrary $`h^{}𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(m_0,m_1)`$, there is a unique triple $`(\widehat{h},k,\lambda )`$, with $`\widehat{h}𝒫_{a,b,1e,1f}^{p,p^{}}(m_1,m^{})`$ for some $`m^{}`$, such that the action of the $`(k,\lambda )`$-transform on $`\widehat{h}`$ yields $`h^{}`$. Then, via the $`𝒟`$-transform, we obtain a unique $`h𝒫_{a,b,e,f}^{p^{}p,p^{}}(m_1,m^{\prime \prime })`$, for some $`m^{\prime \prime }`$. The bijection then follows.
The expression for $`\stackrel{~}{wt}(h)`$ also results from Lemma 4.4. $`\mathrm{}`$
Note that the above lemma excludes the case for which $`\delta _{a,e}^{p^{}p,p^{}}=1`$. Once more, similar results fail in that case.
###### Corollary 4.6
For $`1p<p^{}<2p`$, let $`1a,b<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a,e}^{p^{}p,p^{}}=0`$. Set $`a^{}=a+1e+ap/p^{}`$ and $`b^{}=b+1f+bp/p^{}`$. Fix $`m_0,m_10`$. Then
$$\begin{array}{c}\stackrel{~}{\chi }_{a^{},b^{},1e,1f}^{p,p^{}+p}(m_0,m_1;q)=\hfill \\ q^{\frac{1}{4}\left(m_1^2+(m_0m_1)^2\alpha ^2\beta ^2\right)}\underset{\genfrac{}{}{0pt}{}{mm_0m_1}{\left(\text{mod}\mathrm{\hspace{0.17em}2}\right)}}{}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_0+m_1m)}{m_1}\right]_q\stackrel{~}{\chi }_{a,b,e,f}^{p^{}p,p^{}}(m_1,m;q^1),\hfill \end{array}$$
where $`\alpha =\alpha _{a,b}^{p,p^{}}`$ and $`\beta =\beta _{a,b,1e,1f}^{p,p^{}}`$.
Proof: Apart from the case where $`m_1=0`$ and $`ef`$, this follows immediately from Lemma 4.5 on setting $`m=m_0m_12k`$, once it is noted, via Lemma 2.6, that $`\left[\genfrac{}{}{0pt}{}{k+m_1}{m_1}\right]_q`$ is the generating function for $`𝒴(k,m_1)`$.
The case $`m_1=0`$ and $`ef`$ is dealt with exactly as in the proof of Corollary 3.14. $`\mathrm{}`$
###### Lemma 4.7
Let $`1p<p^{}<2p`$ with $`p`$ co-prime to $`p^{}`$, $`1a,b<p^{}`$ and $`e,f\{0,1\}`$ and set $`a^{}=a+1e+ap/p^{}`$ and $`b^{}=b+1f+bp/p^{}`$. Then $`a^{}p/(p^{}+p)=a1a(p^{}p)/p^{}`$ and $`b^{}p/(p^{}+p)=b1b(p^{}p)/p^{}`$. In addition, $`\alpha _{a^{},b^{}}^{p,p^{}+p}=2\alpha _{a,b}^{p^{}p,p^{}}\beta _{a,b,e,f}^{p^{}p,p^{}}`$ and $`\beta _{a^{},b^{},1e,1f}^{p,p^{}+p}=\alpha _{a,b}^{p^{}p,p^{}}\beta _{a,b,e,f}^{p^{}p,p^{}}`$.
Proof: By Lemma 4.2 and Corollary 4.3, $`ap/p^{}=a1a(p^{}p)/p^{}`$, $`bp/p^{}=b1b(p^{}p)/p^{}`$, $`\alpha _{a,b}^{p,p^{}}=\alpha _{a,b}^{p^{}p,p^{}}`$ and $`\beta _{a,b,1e,1f}^{p,p^{}}=\alpha _{a,b}^{p^{}p,p^{}}\beta _{a,b,e,f}^{p^{}p,p^{}}`$. The current lemma then follows immediately from Lemma 3.5. $`\mathrm{}`$
## 5 The structure of the $`(p,p^{})`$-model
### 5.1 Continued fractions
If $`p^{}`$ and $`p`$ are positive co-prime integers and
$$\frac{p^{}}{p}=c_0+\frac{1}{c_1+{\displaystyle \frac{1}{c_2+{\displaystyle \frac{1}{\frac{\mathrm{}}{c_{n1}+{\displaystyle \frac{1}{c_n}}}}}}}}$$
with $`c_00`$, $`c_i1`$ for $`0<i<n`$, and $`c_n2`$, then $`(c_0,c_1,c_2,\mathrm{},c_n)`$ is said to be the continued fraction for $`p^{}/p`$.
We refer to $`n`$ as the height of $`p^{}/p`$. We set $`t=c_0+c_1+\mathrm{}+c_n2`$ and refer to it as the rank of $`p^{}/p`$. The height and rank of $`𝒫_{a,b,c}^{p,p^{}}(L)`$ are then defined to be equal to those of $`p^{}/p`$.
For $`0kn+1`$, we also define
$$t_k=1+\underset{i=0}{\overset{k1}{}}c_i.$$
(13)
Then $`t_{n+1}=t+1`$ and $`t_nt1`$. We say that the index $`j`$ with $`0jt_{n+1}`$ is in zone $`k`$ if $`t_k<jt_{k+1}`$. We then write $`k=\zeta (j)`$. Note that there are $`n+1`$ zones and that for $`0kn`$, zone $`k`$ contains $`c_k`$ indices.
### 5.2 The Takahashi and string lengths
Given positive co-prime integers $`p`$ and $`p^{}`$ with $`p^{}/p`$ having rank $`t`$, define the set $`\{\kappa _i\}_{i=0}^t`$ of Takahashi lengths, the set $`\{\stackrel{~}{\kappa }_i\}_{i=0}^t`$ of truncated Takahashi lengths, and the set $`\{l_i\}_{i=0}^t`$ of string lengths as follows. First define $`y_k`$ and $`z_k`$ for $`1kn+1`$ by:
$$\begin{array}{cccccc}y_1\hfill & =& 0;\hfill & z_1\hfill & =& 1;\hfill \\ y_0\hfill & =& 1;\hfill & z_0\hfill & =& 0;\hfill \\ y_k\hfill & =& c_{k1}y_{k1}+y_{k2};\hfill & z_k\hfill & =& c_{k1}z_{k1}+z_{k2},(1kn+1).\hfill \end{array}$$
Now for $`t_k<jt_{k+1}`$ and $`0kn`$, set
$`\kappa _j`$ $`=`$ $`y_{k1}+(jt_k)y_k;`$
$`\stackrel{~}{\kappa }_j`$ $`=`$ $`z_{k1}+(jt_k)z_k;`$
$`l_j`$ $`=`$ $`y_{k1}+(jt_k1)y_k.`$
Note that $`\kappa _j=l_{j+1}`$ unless $`j=t_k`$ for some $`k`$, in which case $`\kappa _{t_k}=y_k`$ and $`l_{t_k+1}=y_{k1}`$. We define $`𝒯=\{\kappa _i\}_{i=0}^{t1}`$ and $`𝒯^{}=\{p^{}\kappa _i\}_{i=0}^{t1}`$. (We don’t include $`\kappa _t`$ in the former since it is present in the latter.) Then, for $`n>0`$, $`𝒯𝒯^{}=\mathrm{}`$.<sup>12</sup><sup>12</sup>12In fact, when $`n=0`$, $`𝒯𝒯^{}=\{2,3,\mathrm{},p^{}2\}`$. Then, if $`2ap^{}2`$, different fermionic expressions for $`𝒫_{a,b,c}^{p,p^{}}(L)`$ arise by considering either $`a𝒯`$ or $`a𝒯^{}`$. The same holds for $`2bp^{}2`$. This $`n=0`$ case was fully examined in .
For example, in the case $`p^{}=38`$, $`p=11`$, for which the continued fraction is $`(3,2,5)`$, so that $`n=2`$, $`(t_1,t_2,t_3)=(2,4,9)`$ and $`t=8`$. We then obtain:
$$\begin{array}{c}(y_1,y_0,y_1,y_2,y_3)=(0,1,3,7,38),\hfill \\ (z_1,z_0,z_1,z_2,z_3)=(1,0,1,2,11),\hfill \\ (\kappa _0,\kappa _1,\kappa _2,\kappa _3,\kappa _4,\kappa _5,\kappa _6,\kappa _7)=(1,2,3,4,7,10,17,24),\hfill \\ (l_1,l_2,l_3,l_4,l_5,l_6,l_7,l_8)=(1,2,1,4,3,10,17,24),\hfill \\ (\stackrel{~}{\kappa }_0,\stackrel{~}{\kappa }_1,\stackrel{~}{\kappa }_2,\stackrel{~}{\kappa }_3,\stackrel{~}{\kappa }_4,\stackrel{~}{\kappa }_5,\stackrel{~}{\kappa }_6,\stackrel{~}{\kappa }_7)=(1,1,1,1,2,3,5,7).\hfill \end{array}$$
An induction argument readily establishes that if $`1kn+1`$, then $`y_kz_{k1}y_{k1}z_k=(1)^k`$, that $`y_k`$ is co-prime to $`z_k`$, and that $`y_k/z_k`$ has continued fraction $`(c_0,c_1,\mathrm{},c_{k1})`$. Thus, in particular, $`y_{n+1}=p^{}`$ and $`z_{n+1}=p`$.
## 6 Segmenting the model
### 6.1 Model comparisons
Here, we relate the parameters associated with the $`(p,p^{})`$-model for which the continued fraction is $`(c_0,c_1,\mathrm{},c_n)`$ to those associated with certain ‘simpler’ models. In particular, if $`c_0>1`$, we compare them with those associated with the $`(p,p^{}p)`$-model and, if $`c_0=1`$, we compare them with those associated with the $`(p^{}p,p^{})`$-model.
In the following two lemmas, the parameters associated with those simpler models will be primed to distinguish them from those associated with the $`(p,p^{})`$-model. In particular if $`c_0>1`$, $`(p^{}p)/p`$ has continued fraction $`(c_01,c_1,\mathrm{},c_n)`$, so that in this case, $`t^{}=t1`$, $`n^{}=n`$ and $`t_k^{}=t_k1`$ for $`1kn`$. If $`c_0=1`$, $`p^{}/(p^{}p)`$ has continued fraction $`(c_1+1,c_2,\mathrm{},c_n)`$, so that in this case, $`t^{}=t`$, $`n^{}=n1`$ and $`t_k^{}=t_{k+1}`$ for $`1kn^{}`$.
###### Lemma 6.1
Let $`c_0>1`$. For $`1kn`$ and $`0jt`$, let $`y_k`$, $`z_k`$, $`\kappa _j`$ and $`\stackrel{~}{\kappa }_j`$ be the parameters associated with the $`(p,p^{})`$-model as defined in Section 5.2. For $`1kn`$ and $`0jt^{}`$, let $`y_k^{}`$, $`z_k^{}`$, $`\kappa _j^{}`$ and $`\stackrel{~}{\kappa }_j^{}`$ be the corresponding parameters for the $`(p,p^{}p)`$-model. Then:
* $`y_k=y_k^{}+z_k^{}(0kn)`$;
* $`z_k=z_k^{}(0kn)`$;
* $`\kappa _j=\kappa _{j1}^{}+\stackrel{~}{\kappa }_{j1}^{}(1jt)`$;
* $`\stackrel{~}{\kappa }_j=\stackrel{~}{\kappa }_{j1}^{}(1jt)`$.
Proof: This result is a straightforward consequence of the definitions. $`\mathrm{}`$
###### Lemma 6.2
Let $`c_0=1`$. For $`1kn`$ and $`0jt`$, let $`y_k`$, $`z_k`$, $`\kappa _j`$ and $`\stackrel{~}{\kappa }_j`$ be the parameters associated with the $`(p,p^{})`$-model as defined in Section 5.2. For $`1kn^{}`$ and $`0jt`$, let $`y_k^{}`$, $`z_k^{}`$, $`\kappa _j^{}`$ and $`\stackrel{~}{\kappa }_j^{}`$ be the corresponding parameters for the $`(p^{}p,p^{})`$-model. Then:
* $`y_k=y_{k1}^{}(1kn)`$;
* $`z_k=y_{k1}^{}z_{k1}^{}(1kn)`$;
* $`\kappa _j=\kappa _j^{}(1jt)`$;
* $`\stackrel{~}{\kappa }_j=\kappa _j^{}\stackrel{~}{\kappa }_j^{}(1jt)`$.
Proof: Again, this result is a straightforward consequence of the definitions. $`\mathrm{}`$
###### Lemma 6.3
If $`t_1jt`$ then<sup>13</sup><sup>13</sup>13We use the notation $`\delta _{i,j}^{(2)}=1`$ if $`ij(\text{mod}\mathrm{\hspace{0.17em}2})`$ and $`\delta _{i,j}^{(2)}=0`$ if $`ij(\text{mod}\mathrm{\hspace{0.17em}2})`$.
$$\frac{\stackrel{~}{\kappa }_jp^{}}{p}=\kappa _j\delta _{\zeta (j),1}^{(2)},$$
and if $`0jt`$ then
$$\frac{\kappa _jp}{p^{}}=\stackrel{~}{\kappa }_j\delta _{\zeta (j),0}^{(2)}.$$
Proof: We prove the first of these two results by induction on the sum of the height and rank of $`p^{}/p`$. Since $`\kappa _{t_1}=c_0`$ and $`\stackrel{~}{\kappa }_{t_1}=1`$ and $`\zeta (t_1)=0`$, the required result always holds for the case $`j=t_1`$. In particular, it certainly holds in the case where the sum of the height and rank of $`p^{}/p`$ is at most 2.
Now assume that the first part holds in the case that sum of height and rank is $`n+t1`$, and consider the case where $`p^{}/p`$ has height $`n`$ and rank $`t`$. First assume that $`p^{}>2p`$. For $`jt_1`$, the induction hypothesis implies that $`\kappa _{j1}^{}\delta _{\zeta ^{}(j1),1}^{(2)}<\stackrel{~}{\kappa }_{j1}^{}(p^{}p)/p<\kappa _{j1}^{}\delta _{\zeta ^{}(j1),1}^{(2)}+1`$, where the primed quantities pertain to the continued fraction of $`(p^{}p)/p`$. Using Lemma 6.1 and noting that $`\zeta ^{}(j1)=\zeta (j)`$, readily yields $`\kappa _j\delta _{\zeta (j),1}^{(2)}<\stackrel{~}{\kappa }_jp^{}/p<\kappa _j\delta _{\zeta (j),1}^{(2)}+1`$. This immediately gives the required result.
In the case $`p^{}<2p`$, first let $`jt_2`$. The induction hypothesis implies that $`\kappa _j^{}\delta _{\zeta ^{}(j),1}^{(2)}<\stackrel{~}{\kappa }_j^{}p^{}/(p^{}p)<\kappa _j^{}\delta _{\zeta ^{}(j),1}^{(2)}+1`$, where the primed quantities pertain to the continued fraction of $`p^{}/(p^{}p)`$. Using Lemma 6.2 and noting that $`\zeta ^{}(j)=\zeta (j)1`$, readily yields $`\kappa _j\delta _{\zeta (j),1}^{(2)}(p^{}p)/p<\stackrel{~}{\kappa }_jp^{}/p<\kappa _j+(1\delta _{\zeta (j),1}^{(2)})(p^{}p)/p`$. Since $`(p^{}p)/p<1`$, this implies the required result.
When $`p^{}<2p`$, we have $`c_0=1`$ so that $`t_1=0`$ and $`t_2=c_1`$. Then $`\stackrel{~}{\kappa }_j=j`$ for $`t_1<jt_2`$, whereupon in view of the continued fraction expression for $`p^{}/p`$, we immediately obtain $`\stackrel{~}{\kappa }_jp^{}/p=j=\kappa _j1`$, as required.
The first part of the lemma then follows by induction. For $`t_1jt`$, the second part readily follows from the first. For $`0jt_1t`$, both sides are clearly equal to 0. $`\mathrm{}`$
If $`t_1jt`$, it follows from this result that, with $`k`$ such that $`t_k<jt_{k+1}`$, the $`\stackrel{~}{\kappa }_j`$th odd band in the $`(p,p^{})`$-model lies between heights $`\kappa _j1`$ and $`\kappa _j`$ when $`k`$ is odd, and between heights $`\kappa _j`$ and $`\kappa _j+1`$ when $`k`$ is even. Since there are no adjacent odd bands when $`p^{}>2p`$, it follows that $`\kappa _j`$ is interfacial when $`jt_1`$. On switching the parity of each band, we then obtain in the case $`p^{}<2p`$ that $`\kappa _j`$ is interfacial when $`jt_2`$.
###### Lemma 6.4
If $`1p<p^{}`$ and $`p`$ is co-prime to $`p^{}`$, then for $`1sy_n2`$, the $`s`$th band of the $`(p,p^{})`$-model is of the same parity as the $`s`$th band of the $`(z_n,y_n)`$-model.
Proof: We must establish that $`sz_n/y_n=sz_{n+1}/y_{n+1}`$ for $`1sy_n1`$.
With $`s`$ such that $`1s<y_n`$, let $`r=sz_{n+1}/y_{n+1}`$. Using $`y_nz_{n+1}=y_{n+1}z_n+(1)^n`$ then yields:
$$ry_n(1)^n\frac{s}{y_{n+1}}sz_n<(r+1)y_n(1)^n\frac{s}{y_{n+1}}.$$
Since $`1s<y_n<y_{n+1}`$, the first inequality here implies that $`sz_n/y_nr`$. For the same reasons, and noting that $`sz_n/y_n`$ is not integral, the second inequality here implies that $`sz_n/y_n<r+1`$. The lemma then follows. $`\mathrm{}`$
This lemma shows that the $`(z_n,y_n)`$-model resides within the $`(p,p^{})`$-model, between heights $`1`$ and $`y_n1`$. The up-down symmetry of the $`(p,p^{})`$-model then also implies that the $`(z_n,y_n)`$-model also resides within the $`(p,p^{})`$-model, between heights $`p^{}y_n+1`$ and $`p^{}1`$.
### 6.2 Interfacial retention
We now show that if $`h`$ attains an interfacial height, then the path resulting from the action of a $``$-transform on $`h`$ attains the corresponding interfacial height.
###### Lemma 6.5
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$, and let $`h𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(L^{})`$ result from the action of a $`(k,\lambda )`$-transform on $`h`$. Let $`s`$ be interfacial in the $`(p,p^{})`$-model with $`asb`$, and set $`r=(s+1)p/p^{}`$. Then $`s+r`$ is interfacial in the $`(p,p^{}+p)`$-model.
If $`h_i=s`$ for $`0iL`$ then $`h_j^{}=s+r`$ for some $`j`$ with $`0jL^{}`$. On the other hand, if $`h_j^{}=s+r`$ for $`0jL^{}`$ then $`h_i=s`$ for some $`i`$ with $`0iL`$.
Proof: First note that $`s`$ borders the $`r`$th odd band in the $`(p,p^{})`$-model. If $`s`$ is at the lower (resp. upper) edge of the $`r`$th odd band in the $`(p,p^{})`$-model then $`s+r`$ is at the lower (resp. upper) edge of the $`r`$th odd band in the $`(p,p^{}+p)`$-model. In particular, this implies that $`s+r`$ is interfacial in the $`(p,p^{}+p)`$-model. Then note that in the $`(p,p^{})`$-model, there is at least one even band between the two odd bands on either side of $`s`$ (assume that there is an odd band immediately above and immediately below the $`(p,p^{})`$-model grid if necessary). Thus there are at least two even bands between the two odd bands on either side of $`s+r`$ in the $`(p,p^{}+p)`$-model.
Let $`h^{(0)}`$ result from the action of the $`_1`$-transform on $`h`$. The definition of this transform implies that if $`h_i=s`$ for some $`i`$ then $`h_j^{(0)}=s+r`$ for some $`j`$ and vice-versa (when $`\delta _{a,e}^{p,p^{}}=1`$ or $`\delta _{b,f}^{p,p^{}}=1`$, this statement relies on $`asb`$).
If $`h^{(k)}`$ results from the action of the $`_2(k)`$-transform on $`h^{(0)}`$, then if $`h_j^{(0)}=s+r`$ for some $`j`$ then $`h_j^{}^{(k)}=s+r`$ for some $`j^{}`$ and vice-versa (this statement relies on the two odd bands either side of $`s+r`$ having at least two even bands between them).
If $`h^{}`$ results from the action of the $`_3(\lambda )`$-transform on $`h^{(k)}`$, then if $`h_j^{(0)}=s+r`$ for some $`j`$, examination of the ten particle moves and edge-moves described in Section 3.3, shows that $`h_j^{}^{(k)}=s+r`$ for some $`j^{}`$ and vice-versa (this statement also relies on the two odd bands either side of $`s+r`$ having at least two even bands between them). Combining these results proves the lemma. $`\mathrm{}`$
We also need the analogue of this result for the $`𝒟`$-transform.
###### Lemma 6.6
Let $`h𝒫_{a,b,e,f}^{p^{}p,p^{}}(L)`$ and let $`h^{}𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(L^{})`$ result from the action of a $`𝒟`$-transform on $`h`$ followed by a $`(k,\lambda )`$-transform. Let $`s`$ be interfacial in the $`(p,p^{})`$-model with $`asb`$, and set $`r=(s+1)p/p^{}`$. Then $`s+r`$ is interfacial in the $`(p,p^{}+p)`$-model.
If $`h_i=s`$ for $`0iL`$ then $`h_j^{}=s+r`$ for some $`j`$ with $`0jL^{}`$. On the other hand, if $`h_j^{}=s+r`$ for $`0jL^{}`$ then $`h_i=s`$ for some $`i`$ with $`0iL`$.
Proof: This follows immediately from the above result after noting that if $`s`$ is interfacial in the $`(p^{}p,p^{})`$-model then it is also in the $`(p,p^{})`$-model. $`\mathrm{}`$
A set $`𝒮`$ is said to be interfacial in the $`(p,p^{})`$-model if each $`s𝒮`$ is interfacial in the $`(p,p^{})`$-model. We now define $`𝒫_{a,b,e,f}^{p,p^{}}(L,m)\{𝒮\}`$ to be the subset of $`𝒫_{a,b,e,f}^{p,p^{}}(L,m)`$ comprising those paths $`h`$ for which for each $`s𝒮`$, there exists $`i`$ with $`0iL`$ such that $`h_i=s`$. The generating function for this set is
$$\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L;q)\{𝒮\}=\underset{h𝒫_{a,b,e,f}^{p,p^{}}(L)\{𝒮\}}{}q^{\stackrel{~}{wt}(h)}.$$
Of course, $`𝒫_{a,b,e,f}^{p,p^{}}(L,m)\{\mathrm{}\}=𝒫_{a,b,e,f}^{p,p^{}}(L,m)`$.
Given $`𝒮`$ as above, we now define $`𝒮^{}=\{s+(s+1)p/p^{}:s𝒮\}`$.
###### Corollary 6.7
For $`1p<p^{}`$, let $`1a,b<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a,e}^{p,p^{}}=0`$. Let $`𝒮`$ be interfacial in the $`(p,p^{})`$-model with $`asb`$ for all $`s𝒮`$. Set $`a^{}=a+e+ap/p^{}`$ and $`b^{}=b+f+bp/p^{}`$. Fix $`m_0,m_10`$. Then
$$\begin{array}{c}\stackrel{~}{\chi }_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,m_1)\{𝒮^{}\}\hfill \\ =q^{\frac{1}{4}\left((m_0m_1)^2\beta ^2\right)}\underset{\genfrac{}{}{0pt}{}{mm_0}{\left(\text{mod}\mathrm{\hspace{0.17em}2}\right)}}{}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_0+m)}{m_1}\right]_q\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(m_1,m)\{𝒮\},\hfill \end{array}$$
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
Proof: Combining Lemmas 3.13 and 6.5 implies that the map $`(h,k,\lambda )h^{}`$ effected by the action of a $`(k,\lambda )`$-transform on $`h`$, is a bijection between $`_k𝒫_{a,b,e,f}^{p,p^{}}(m_1,2k+2m_1m_0)\{𝒮\}\times 𝒴(k,m_1)`$ and $`𝒫_{a^{},b^{},e,f}^{p,p^{}+p}(m_0,m_1)\{𝒮^{}\}`$. The result then follows as in the proof of Corollary 3.14. $`\mathrm{}`$
###### Corollary 6.8
For $`1p<p^{}<2p`$, let $`1a,b<p^{}`$ and $`e,f\{0,1\}`$, with $`\delta _{a,e}^{p^{}p,p^{}}=0`$. Let $`𝒮`$ be interfacial in the $`(p,p^{})`$-model with $`asb`$ for all $`s𝒮`$. Set $`a^{}=a+1e+ap/p^{}`$ and $`b^{}=b+1f+bp/p^{}`$. Fix $`m_0,m_10`$. Then
$$\begin{array}{c}\stackrel{~}{\chi }_{a^{},b^{},1e,1f}^{p,p^{}+p}(m_0,m_1;q)\{𝒮^{}\}\hfill \\ =q^{\frac{1}{4}\left(m_1^2+(m_0m_1)^2\alpha ^2\beta ^2\right)}\hfill \\ \times \underset{\genfrac{}{}{0pt}{}{mm_0m_1}{\left(\text{mod}\mathrm{\hspace{0.17em}2}\right)}}{}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_0+m_1m)}{m_1}\right]_q\stackrel{~}{\chi }_{a,b,e,f}^{p^{}p,p^{}}(m_1,m;q^1)\{𝒮\},\hfill \end{array}$$
where $`\alpha =\alpha _{a,b}^{p,p^{}}`$ and $`\beta =\beta _{a,b,1e,1f}^{p,p^{}}`$.
Proof: Combining Lemmas 4.5 and 6.6 implies that the map $`(h,k,\lambda )h^{}`$ effected by the action of a $`𝒟`$-transform on $`h`$ immediately followed by a $`(k,\lambda )`$-transform, is a bijection between $`_k𝒫_{a,b,e,f}^{p^{}p,p^{}}(m_1,m_0m_12k)\{𝒮\}\times 𝒴(k,m_1)`$ and $`𝒫_{a^{},b^{},1e,1f}^{p,p^{}+p}(m_0,m_1)\{𝒮^{}\}`$. The result then follows as in the proof of Corollary 4.6. $`\mathrm{}`$
## 7 Extending and truncating paths
### 7.1 Extending paths
In this section, we specify a process by which a path $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ may be extended by a single unit to its left, or by a single unit to its right. One extension may follow the other to yield a path of length $`L+2`$.
Path extension on the left is restricted to where $`\delta _{a,e}^{p,p^{}}=0`$ so that the pre-segment of $`h`$ lies in the even band.
We obtain $`h^{}`$ by defining $`h_0^{}=a^{}=a+(1)^e`$ and $`h_i^{}=h_{i1}`$ for $`1iL+1`$. In particular, $`\pi (h^{})=0`$. We also define $`e(h^{})=e^{}=1e`$, so that then $`h^{}𝒫_{a^{},b,e^{},f}^{p,p^{}}(L+1)`$.
This extending process is depicted in Fig. 7.
###### Lemma 7.1
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$, where $`\delta _{a,e}^{p,p^{}}=0`$. Let $`h^{}𝒫_{a^{},b,e^{},f}^{p,p^{}}(L^{})`$ be obtained from $`h`$ by the above process of path extension. If $`\mathrm{\Delta }=a^{}a`$ then $`\mathrm{\Delta }=(1)^e=(1)^e^{}`$, and
$$\begin{array}{c}L^{}=L+1;\hfill \\ m(h^{})=m(h);\hfill \\ \stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+\frac{1}{2}(Lm(h)+\mathrm{\Delta }\beta (h)).\hfill \end{array}$$
Furthermore, $`\alpha _{a^{},b}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}\mathrm{\Delta }`$ and $`\beta _{a^{},b,e^{},f}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}\mathrm{\Delta }`$.
Proof: That $`\mathrm{\Delta }=(1)^e=(1)^e^{}`$ is immediate from the definition. Let $`h`$ have striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)}`$.
If $`e=d`$, we are restricted to the case $`\pi (h)=0`$, since $`\delta _{a,e}^{p,p^{}}=0`$. The striking sequence of $`h^{}`$ is then $`\left(\genfrac{}{}{0pt}{}{0}{1}\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e^{},f,e^{})}`$. Thereupon, since $`\pi (h^{})=0`$, we obtain $`m(h^{})=m(h)`$. In this case we immediately obtain, via Lemma 2.1, that $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+(b_1+b_3+\mathrm{})`$. Thereupon, since $`\mathrm{\Delta }=(1)^e=(1)^d`$, $`\beta (h)=(1)^d((b_1+b_3+\mathrm{})(b_2+b_4+\mathrm{}))`$ and $`m(h)=(a_1+a_2+a_3+\mathrm{})`$, we obtain $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+(Lm(h)+\mathrm{\Delta }\beta (h))/2`$.
If $`ed`$, the striking sequence of $`h^{}`$ is $`\left(\genfrac{}{}{0pt}{}{a_1+1\pi }{b_1+\pi }\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e^{},f,e^{})}`$. Then $`m(h^{})=1\pi +_{i=1}^la_i`$ which equals $`m(h)=(e+d+\pi )\text{mod}\mathrm{\hspace{0.17em}2}+_{i=1}^la_i`$ for both $`\pi =0`$ and $`\pi =1`$. Here Lemma 2.1 implies that $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+(b_2+b_4+\mathrm{})`$. Thereupon, since $`\mathrm{\Delta }=(1)^e=(1)^d`$, $`\beta (h)=(1)^d((b_1+b_3+\mathrm{})(1\pi +b_2+b_4+\mathrm{}))`$ and $`m(h)=(1\pi +a_1+a_2+a_3+\mathrm{})`$, we also obtain $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+(Lm(h)+\mathrm{\Delta }\beta (h))/2`$.
That $`\alpha _{a^{},b}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}\mathrm{\Delta }`$ is immediate. Since $`\pi (h^{})=0`$ then $`a^{}p/p^{}=ap/p^{}`$. That $`\beta _{a^{},b,e^{},f}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}\mathrm{\Delta }`$ now follows. $`\mathrm{}`$
In the following lemma, we consider the special case when $`2p<p^{}<3p`$ so that the first and second bands of the $`(p,p^{})`$-model are even and odd respectively. We then only consider path extension into the first or the $`(p^{}2)`$th band of the $`(p,p^{})`$-model.
###### Lemma 7.2
Let $`2<2p<p^{}<3p`$ and either $`a=2`$ and $`e=1`$, or $`a=p^{}2`$ and $`e=0`$. Then $`a`$ is interfacial in the $`(p,p^{})`$-model. Let $`𝒮`$ be interfacial in the $`(p,p^{})`$-model, and set $`\mathrm{\Delta }=(1)^e`$, $`a^{}=a+\mathrm{\Delta }`$ and $`e^{}=1e`$. Then:
$$\stackrel{~}{\chi }_{a^{},b,e^{},f}^{p,p^{}}(L,m)\{𝒮\{a\}\}=q^{\frac{1}{2}(L1m+\mathrm{\Delta }\beta )}\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L1,m)\{𝒮\},$$
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
In addition, $`\alpha _{a^{},b}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}\mathrm{\Delta }`$, $`\beta _{a^{},b,e^{},f}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}\mathrm{\Delta }`$.
Proof: Since $`2p<p^{}<3p`$, it follows that $`0=2p/p^{}3p/p^{}`$ wherepon $`2`$ and $`p^{}2`$ are both interfacial in the $`(p,p^{})`$-model, and $`\delta _{a,e}^{p,p^{}}=0`$.
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L1,m)\{𝒮\}`$. Extend $`h`$ on the left to obtain $`h^{}`$ with $`h_0^{}=a^{}=a+\mathrm{\Delta }`$. Clearly, $`h^{}`$ attains $`a`$. Then, Lemma 7.1 implies that $`h^{}𝒫_{a^{},b,e^{},f}^{p,p^{}}(L,m)\{𝒮\{a\}\}`$.
Conversely, any such $`h^{}`$ arises from some $`h𝒫_{a,b,e,f}^{p,p^{}}(L1,m)\{𝒮\}`$ in this way since either $`h_0^{}=1`$ and $`e^{}=0`$, or $`h_0^{}=p^{}1`$ and $`e^{}=1`$. The result then follows from the expression for $`\stackrel{~}{wt}(h^{})`$ given in Lemma 7.1, and $`\beta (h)=\beta _{a,b,e,f}^{p,p^{}}`$ from Lemma 2.2.
The final statement also follows from Lemma 7.1. $`\mathrm{}`$
For $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$, we now define path extension to the right in a similar way. Here we restrict path extension to the cases where $`\delta _{b,f}^{p,p^{}}=0`$ so that the post-segment of $`h`$ lies in the even band.
We obtain $`h^{}`$ by defining $`h_i^{}=h_i`$ for $`0iL`$ and $`h_{L+1}^{}=b^{}=b+(1)^f`$ and We also define $`f(h^{})=f^{}=1f`$, so that then $`h^{}𝒫_{a,b^{},e,f^{}}^{p,p^{}}(L+1)`$.
This extending process is depicted in Fig. 8.
###### Lemma 7.3
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$, where $`\delta _{b,f}^{p,p^{}}=0`$. Let $`h^{}𝒫_{a,b^{},e,f^{}}^{p,p^{}}(L^{})`$ be obtained from $`h`$ by the above process of path extension. If $`\mathrm{\Delta }=b^{}b`$ then $`\mathrm{\Delta }=(1)^f=(1)^f^{}`$, and
$$\begin{array}{c}L^{}=L+1;\hfill \\ m(h^{})=m(h);\hfill \\ \stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+\frac{1}{2}(L\mathrm{\Delta }\alpha (h)).\hfill \end{array}$$
Furthermore, $`\alpha _{a,b^{}}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}+\mathrm{\Delta }`$ and $`\beta _{a,b^{},e,f^{}}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}+\mathrm{\Delta }`$.
Proof: That $`\mathrm{\Delta }=(1)^f=(1)^f^{}`$ is immediate from the definition. Let $`h`$ have striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\right)^{(e,f,d)}`$. It is easily checked that the $`L`$th vertex of $`h^{}`$ is scoring if and only if the $`L`$th vertex of $`h`$ is scoring.
Then, if the extending segment is in the same direction as the $`L`$th segment, $`h^{}`$ has striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{a_3}{b_3}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l+1}\right)^{(e,f^{},d)}`$ and $`\mathrm{\Delta }=(1)^{d+l}`$. That $`m(h^{})=m(h)`$ is immediate.
When the extending segment is in the direction opposite to that of the $`L`$th segment, $`h^{}`$ has striking sequence $`\left(\genfrac{}{}{0pt}{}{a_1}{b_1}\genfrac{}{}{0pt}{}{a_2}{b_2}\genfrac{}{}{0pt}{}{\mathrm{}}{\mathrm{}}\genfrac{}{}{0pt}{}{a_l}{b_l}\genfrac{}{}{0pt}{}{0}{1}\right)^{(e,f^{},d)}`$ and $`\mathrm{\Delta }=(1)^{d+l}`$. We immediately obtain $`m(h^{})=m(h)`$ in this case.
For $`1il`$, let $`w_i=a_i+b_i`$. We find $`\alpha (h)=(1)^{d+l}((w_l+w_{l2}\mathrm{})(w_{l1}+w_{l3}+\mathrm{}))`$. In the first case above, Lemma 2.1 gives $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+(w_{l1}+w_{l3}+w_{l5}+\mathrm{})`$, whereupon we obtain $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+\frac{1}{2}(L(h)\mathrm{\Delta }\alpha (h))`$. In the second case above, Lemma 2.1 gives $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+(w_l+w_{l2}+w_{l4}+\mathrm{})`$, and we again obtain $`\stackrel{~}{wt}(h^{})=\stackrel{~}{wt}(h)+\frac{1}{2}(L(h)\mathrm{\Delta }\alpha (h))`$.
That $`\alpha _{a,b^{}}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}+\mathrm{\Delta }`$ is immediate. That $`\beta _{a,b^{},e,f^{}}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}+\mathrm{\Delta }`$ now follows because $`bp/p^{}=b^{}p/p^{}`$. $`\mathrm{}`$
###### Lemma 7.4
Let $`2<2p<p^{}<3p`$ and either $`b=2`$ and $`f=1`$, or $`b=p^{}2`$ and $`f=0`$. Then $`b`$ is interfacial in the $`(p,p^{})`$-model. Let $`𝒮`$ be interfacial in the $`(p,p^{})`$-model, and set $`\mathrm{\Delta }=(1)^f`$, $`b^{}=b+\mathrm{\Delta }`$ and $`f^{}=1f`$. Then:
$$\stackrel{~}{\chi }_{a,b^{},e,f^{}}^{p,p^{}}(L,m)\{𝒮\{b\}\}=q^{\frac{1}{2}(L1\mathrm{\Delta }\alpha )}\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L1,m)\{𝒮\},$$
where $`\alpha =\alpha _{a,b}^{p,p^{}}`$.
In addition, $`\alpha _{a,b^{}}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}+\mathrm{\Delta }`$ and $`\beta _{a,b^{},e,f^{}}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}+\mathrm{\Delta }`$.
Proof: Proof: Since $`2p<p^{}<3p`$, it follows that $`0=2p/p^{}3p/p^{}`$ wherepon $`2`$ and $`p^{}2`$ are both interfacial in the $`(p,p^{})`$-model, and $`\delta _{b,f}^{p,p^{}}=0`$.
Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L1,m)\{𝒮\}`$. Extend this path on the right to obtain $`h^{}`$ with $`h_L^{}=b^{}=b+\mathrm{\Delta }`$. Clearly, $`h^{}`$ attains height $`b`$. Then, via Lemma 7.3, $`h^{}𝒫_{a,b^{},e,f^{}}^{p,p^{}}(L,m)\{𝒮\{b\}\}`$. Conversely, any such $`h^{}`$ arises in this way from some $`h𝒫_{a,b,e,f}^{p,p^{}}(L1,m)\{𝒮\}`$, since either $`h_L^{}=1`$ and $`f^{}=1`$, or $`h_L^{}=p^{}1`$ and $`f^{}=0`$. The required result then follows from the expression for $`\stackrel{~}{wt}(h^{})`$ given in Lemma 7.3, and $`\alpha (h)=\alpha _{a,b}^{p,p^{}}`$ from Lemma 2.2.
The final statement follows from Lemma 7.3. $`\mathrm{}`$
### 7.2 Truncating paths
In this section, we specify a process by which a path $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$, for $`L>0`$ may be shortened by removing just the leftmost (first) segment, or by removing just the rightmost ($`L`$th) segment. Consequently, the new path $`h^{}`$ is of length $`L^{}=L1`$. One shortening may follow the other to yield a path of length $`L2`$.
In fact, we will only use these shortening processes when $`p^{}>2p`$, so that in particular, the 1st and the $`(p^{}2)`$th bands of the $`(p,p^{})`$-model are even.
Shortening on the left side will occur only when $`a=1`$ or $`a=p^{}1`$ so that the removed segment is in an even band, and will occur when the 0th vertex is scoring.
###### Lemma 7.5
Let $`p^{}>2p`$ and either $`a=1`$ and $`e=0`$, or $`a=p^{}1`$ and $`e=1`$. Let $`𝒮`$ be interfacial in the $`(p,p^{})`$-model, with $`a𝒮`$. Define $`\mathrm{\Delta }=(1)^e`$, $`e^{}=1e`$ and $`a^{}=a\mathrm{\Delta }`$. Then
$$\begin{array}{c}\stackrel{~}{\chi }_{a^{},b,e^{},f}^{p,p^{}}(L,m)\{𝒮\}=q^{\frac{1}{2}(L+1m+\mathrm{\Delta }\beta )}\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L+1,m)\{𝒮\},\hfill \end{array}$$
where $`\beta =\beta _{a,b,e,f}^{p,p^{}}`$.
In addition, $`\alpha _{a^{},b}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}+\mathrm{\Delta }`$, and $`\beta _{a^{},b,e^{},f}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}+\mathrm{\Delta }`$.
Proof: Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L+1,m)\{𝒮\}`$, and note that necessarily $`h_1=a^{}`$. Let $`h^{}𝒫_{a^{},b,e^{},f}^{p,p^{}}(L,m)\{𝒮\}`$ be defined by $`h_i^{}=h_{i+1}`$ for $`0iL`$. The lemma then follows on noting that $`\delta _{a^{},e^{}}^{p,p^{}}=0`$ and using Lemma 7.1 after switching the roles of $`h`$ and $`h^{}`$ there. $`\mathrm{}`$
Shortening on the right side will occur only when $`b=1`$ or $`b=p^{}1`$ so that the removed segment is in an even band, and will occur when the $`L`$th vertex is scoring.
###### Lemma 7.6
Let $`p^{}>2p`$ and either $`b=1`$ and $`f=0`$, or $`b=p^{}1`$ and $`f=1`$. Let $`𝒮`$ be interfacial in the $`(p,p^{})`$-model, with $`b𝒮`$. Define $`\mathrm{\Delta }=(1)^f`$, $`f^{}=1f`$ and $`b^{}=b\mathrm{\Delta }`$. Then
$$\stackrel{~}{\chi }_{a,b^{},e,f^{}}^{p,p^{}}(L,m)\{𝒮\}=q^{\frac{1}{2}(L+1\mathrm{\Delta }\alpha )}\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L+1,m)\{𝒮\},$$
where $`\alpha =\alpha _{a,b}^{p,p^{}}`$.
In addition, $`\alpha _{a,b^{}}^{p,p^{}}=\alpha _{a,b}^{p,p^{}}\mathrm{\Delta }`$, and $`\beta _{a,b^{},e,f^{}}^{p,p^{}}=\beta _{a,b,e,f}^{p,p^{}}\mathrm{\Delta }`$.
Proof: Let $`h𝒫_{a,b,e,f}^{p,p^{}}(L+1,m)\{𝒮\}`$, and note that necessarily $`h_L=b^{}`$. Let $`h^{}𝒫_{a,b^{},e,f^{}}^{p,p^{}}(L,m)\{𝒮\}`$ be defined by $`h_i^{}=h_i`$ for $`0iL`$. The lemma then follows on noting that $`\delta _{b^{},f^{}}^{p,p^{}}=0`$ and using Lemma 7.3 after switching the roles of $`h`$ and $`h^{}`$ there. $`\mathrm{}`$
## 8 Fermionic expressions
### 8.1 Results
In this section, we fix co-prime $`p`$ and $`p^{}`$, and fix $`a,b𝒯𝒯^{}`$, with $`1a,b<p^{}`$. We make use of the definitions of 5.1 and 5.2. For certain $`c`$, we present two fermionic expressions for $`𝒫_{a,b,c}^{p,p^{}}(L)`$. The value of $`c`$ depends on $`b`$ and, for $`p^{}>2p`$, is given by:
$$c=\{\begin{array}{cc}2\hfill & \text{if }b=1;\hfill \\ b1\hfill & \text{if }1<bt_1;\hfill \\ p^{}2\hfill & \text{if }b=p^{}1;\hfill \\ b+1\hfill & \text{if }p^{}t_1b<p^{}1;\hfill \\ b\pm 1\hfill & \text{otherwise.}\hfill \end{array}$$
(14)
For $`p^{}<2p`$, change $`t_1`$ to $`t_2`$ in this definition.
The statement of these fermionic expressions requires the following notation. For convenience, set $`a^L=a`$ and $`a^R=b`$. Now, for $`A\{L,R\}`$, define $`\sigma ^A`$ such that:
$$\kappa _{\sigma ^A}=\{\begin{array}{cc}a^A\hfill & \text{if }a^A𝒯;\hfill \\ p^{}a^A\hfill & \text{if }a^A𝒯^{}.\hfill \end{array}$$
(15)
For $`0jt`$, define<sup>14</sup><sup>14</sup>14In this paper, all vectors $`𝑸`$, $`𝒎`$, $`\widehat{𝒎}`$, $`𝒏`$, $`𝒖`$, $`𝚫`$ and $`𝒆`$ should be considered as column vectors. However, for typographical convenience, we shall express their components in row vector form. $`𝒆_j=(e_1,e_2,\mathrm{},e_t)`$ with $`e_i=\delta _{ij}`$. Then define
$$𝒖^A=𝒆_{\sigma ^A}\underset{k:\sigma ^At_k<t}{}𝒆_{t_k}+\{\begin{array}{cc}0\hfill & \text{if }a^A𝒯;\hfill \\ 𝒆_t\hfill & \text{if }a^A𝒯^{},\hfill \end{array}$$
(16)
and
$$𝚫^A=\{\begin{array}{cc}𝒆_{\sigma ^A}+\underset{k:\sigma ^At_k<t}{}𝒆_{t_k}\hfill & \text{if }a^A𝒯;\hfill \\ 𝒆_t+𝒆_{\sigma ^A}\underset{k:\sigma ^At_k<t}{}𝒆_{t_k}\hfill & \text{if }a^A𝒯^{}.\hfill \end{array}$$
(17)
We define the matrix $`𝑪`$ to be the $`t\times t`$ tri-diagonal matrix with entries $`𝑪_{ij}`$ for $`0i,jt1`$ where, when the indices are in this range,
$$\begin{array}{cccc}𝑪_{j,j1}=1,& 𝑪_{j,j}=1,& 𝑪_{j,j+1}=1,& \text{if }j=t_k,k=1,2,\mathrm{},n\text{;}\hfill \\ 𝑪_{j,j1}=1,& 𝑪_{j,j}=2,& 𝑪_{j,j+1}=1,& 0j<t\text{ otherwise.}\hfill \end{array}$$
(18)
It is also useful to define $`\widehat{𝑪}`$ to be the $`t\times t`$ upper-triangular matrix with entries $`\widehat{𝑪}_{ij}=𝑪_{ij}`$, as above, with $`1it`$ and $`0jt1`$.
For example, in the case $`p=9`$ and $`p^{}=31`$, where the continued fraction of $`p^{}/p`$ is $`(3,2,4)`$ and $`t_1=2`$, $`t_2=4`$ and $`t_3=8`$, we have:
$$𝑪=\left(\begin{array}{ccccccc}\mathrm{\hspace{0.33em}2}\hfill & 1\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill \\ 1\hfill & \mathrm{\hspace{0.33em}2}\hfill & 1\hfill & .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & 1\hfill & \mathrm{\hspace{0.33em}1}\hfill & \mathrm{\hspace{0.33em}1}\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}2}\hfill & 1\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}1}\hfill & \mathrm{\hspace{0.33em}1}\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}2}\hfill & 1\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}2}\hfill \end{array}\right),\widehat{𝑪}=\left(\begin{array}{ccccccc}1\hfill & \mathrm{\hspace{0.33em}2}\hfill & 1\hfill & .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & 1\hfill & \mathrm{\hspace{0.33em}1}\hfill & \mathrm{\hspace{0.33em}1}\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}2}\hfill & 1\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}1}\hfill & \mathrm{\hspace{0.33em}1}\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}2}\hfill & 1\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & 1\hfill & \mathrm{\hspace{0.33em}2}\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & .\hfill & 1\hfill \end{array}\right).$$
Since $`\widehat{𝑪}`$ is upper-triangular, its inverse is readily obtained. Given a $`t`$-dimensional vector $`𝒖`$, we then define $`Q_i\{0,1\}`$ for $`0i<t`$, by<sup>15</sup><sup>15</sup>15For $`𝒗=(v_1,v_2,\mathrm{},v_t)`$, we define $`𝒗\text{mod}\mathrm{\hspace{0.17em}2}=(v_1\text{mod}\mathrm{\hspace{0.17em}2},v_2\text{mod}\mathrm{\hspace{0.17em}2},\mathrm{},v_t\text{mod}\mathrm{\hspace{0.17em}2},)`$.
$$(Q_0,Q_1,Q_2,\mathrm{},Q_{t1})^T=\widehat{𝑪}{}_{}{}^{1}𝒖\text{mod}\mathrm{\hspace{0.17em}2}.$$
(19)
We thus define the parity vector $`𝑸(𝒖)=(Q_1,Q_2,\mathrm{},Q_{t1})`$.
Now, given a $`t`$-dimensional vector $`𝒖=(u_1,u_2,\mathrm{},u_t)`$, define the $`(t1)`$-dimensional vector $`𝒖^{\mathrm{}}=(u_1^{\mathrm{}},u_2^{\mathrm{}},\mathrm{},u_{t1}^{\mathrm{}})`$ by:
$$u_j^{\mathrm{}}=\{\begin{array}{cc}0& \text{if }t_k<jt_{k+1},k0(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ u_j& \text{if }t_k<jt_{k+1},k0(\text{mod}\mathrm{\hspace{0.17em}2}),\hfill \end{array}$$
(20)
and the $`(t1)`$-dimensional vector $`𝒖^{\mathrm{}}=(u_1^{\mathrm{}},u_2^{\mathrm{}},\mathrm{},u_{t1}^{\mathrm{}})`$ by:
$$u_j^{\mathrm{}}=\{\begin{array}{cc}u_j& \text{if }t_k<jt_{k+1},k0(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ 0& \text{if }t_k<jt_{k+1},k0(\text{mod}\mathrm{\hspace{0.17em}2}).\hfill \end{array}$$
(21)
Then, of course, $`(𝒖)_j=(𝒖^{\mathrm{}}+𝒖^{\mathrm{}})_j`$ for $`1j<t`$. For convenience, we sometimes write $`𝒖_{\mathrm{}}`$ and $`𝒖_{\mathrm{}}`$ for $`𝒖^{\mathrm{}}`$ and $`𝒖^{\mathrm{}}`$ respectively.
Finally, we define a value $`\gamma `$ that depends on $`𝚫^L`$ and $`𝚫^R`$. This value is obtained by iteratively generating the sequences $`(\beta _t,\beta _{t1},\mathrm{},\beta _0)`$, $`(\alpha _t,\alpha _{t1},\mathrm{},\alpha _0)`$, and $`(\gamma _t,\gamma _{t1},\mathrm{},\gamma _0)`$ as follows. Let $`\alpha _t=\beta _t=\gamma _t=0`$. Now, for $`j=t,t1,\mathrm{},1`$, obtain $`\alpha _{j1}`$, $`\beta _{j1}`$, and $`\gamma _{j1}`$ from $`\alpha _j`$, $`\beta _j`$, and $`\gamma _j`$ in the following three stages. Firstly, obtain:
$$(\beta _{j1}^{},\gamma _{j1}^{})=(\beta _j+(𝚫^L)_j(𝚫^R)_j,\gamma _j+2\alpha _j(𝚫^R)_j).$$
(22)
Then obtain:
$$(\alpha _{j1}^{\prime \prime },\gamma _{j1}^{\prime \prime })=(\alpha _j+\beta _{j1}^{},\gamma _{j1}^{}(\beta _{j1}^{})^2).$$
(23)
Finally, set
$$\begin{array}{c}(\alpha _{j1},\beta _{j1},\gamma _{j1})\hfill \\ =\{\begin{array}{cc}(\alpha _{j1}^{\prime \prime },\alpha _{j1}^{\prime \prime }\beta _{j1}^{},(\alpha _{j1}^{\prime \prime })^2\gamma _{j1}^{\prime \prime })& \text{if }j=t_k+1,\mathrm{\hspace{0.33em}1}kn;\hfill \\ (\alpha _{j1}^{\prime \prime },\beta _{j1}^{},\gamma _{j1}^{\prime \prime })& \text{otherwise}.\hfill \end{array}\hfill \end{array}$$
(24)
We then set $`\gamma =\gamma _0`$.
###### Theorem 8.1
If $`a,b𝒯𝒯^{}`$, define everything as above. Then:
$$\begin{array}{c}\chi _{a,b,c}^{p,p^{}}(L)=\hfill \\ \underset{𝒎𝑸(𝒖^L+𝒖^R)}{}q^{\frac{1}{4}\widehat{𝒎}^T𝑪\widehat{𝒎}\frac{1}{4}L^2\frac{1}{2}(𝒖_{\mathrm{}}^L+𝒖_{\mathrm{}}^R)𝒎+\frac{1}{4}\gamma }\underset{j=1}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}𝒖^L𝒖^R)_j}{m_j}\right]_q\hfill \\ +\{\begin{array}{cc}\chi _{a,b,c}^{z_n,y_n}(L)\hfill & \text{if }a<y_n\text{ and }b<y_n;\hfill \\ \chi _{p^{}a,p^{}b,p^{}c}^{z_n,y_n}(L)\hfill & \text{if }a>p^{}y_n\text{ and }b>p^{}y_n;\hfill \\ 0,\hfill & \text{otherwise}.\hfill \end{array}\hfill \end{array}$$
With $`𝐐(𝐮^L+𝐮^R)=(Q_1,Q_2,\mathrm{},Q_{t1})`$, the summation here is over all vectors $`𝐦=(m_1,m_2,\mathrm{},m_{t1})`$ such that $`m_j\mathrm{}_0`$ and $`m_jQ_j(\text{mod}\mathrm{\hspace{0.17em}2})`$ for $`1j<t`$. Then $`\widehat{𝐦}=(L,m_1,m_2,\mathrm{},m_{t1})`$.
The second fermionic expression for $`\chi _{a,b,c}^{p,p^{}}(L)`$ that we present, involves the modified form $`\left[\genfrac{}{}{0pt}{}{A}{B}\right]_q^{}`$ of the Gaussian polynomial defined in (2).
###### Theorem 8.2
If $`a,b𝒯𝒯^{}`$, define everything as above. Then, if $`L0`$:
$$\begin{array}{c}\chi _{a,b,c}^{p,p^{}}(L)=\hfill \\ \underset{𝒎𝑸(𝒖^L+𝒖^R)}{}q^{\frac{1}{4}\widehat{𝒎}^T𝑪\widehat{𝒎}\frac{1}{4}L^2\frac{1}{2}(𝒖_{\mathrm{}}^L+𝒖_{\mathrm{}}^R)𝒎+\frac{1}{4}\gamma }\underset{j=1}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}𝒖^L𝒖^R)_j}{m_j}\right]_q^{}.\hfill \end{array}$$
With $`𝐐(𝐮^L+𝐮^R)=(Q_1,Q_2,\mathrm{},Q_{t1})`$, the summation here is over all vectors $`𝐦=(m_1,m_2,\mathrm{},m_{t1})`$ such that $`m_j\mathrm{}_0`$ and $`m_jQ_j(\text{mod}\mathrm{\hspace{0.17em}2})`$ for $`1j<t`$. Then $`\widehat{𝐦}=(L,m_1,m_2,\mathrm{},m_{t1})`$.
### 8.2 Carrying out the induction
With $`p`$ and $`p^{}`$ fixed, employ the definitions of Section 5.1. Then, for $`0it`$, let $`k(i)`$ be such that $`t_{k(i)}i<t_{k(i)+1}`$ (i.e. $`k(i)=\zeta (i+1)`$), and define $`p_i`$ and $`p_i^{}`$ to be the positive co-prime integers for which $`p_i^{}/p_i`$ has continued fraction $`(t_{k(i)+1}+1i,c_{k(i)+1},\mathrm{},c_n)`$. Thus $`p_i^{}/p_i`$ has rank $`ti`$. As in Section 5.2, we obtain Takahashi lengths $`\{\kappa _j^{(i)}\}_{j=0}^{ti}`$ and truncated Takahashi lengths $`\{\stackrel{~}{\kappa }_j^{(i)}\}_{j=0}^{ti}`$ for $`p_i^{}/p_i`$.
###### Lemma 8.3
Let $`1it`$. If $`it_{k(i)}`$ then:
$$\begin{array}{ccc}p^{(i1)}=\hfill & p^{(i)}+p^{(i)};\hfill & \\ p^{(i1)}=\hfill & p^{(i)};\hfill & \\ \kappa _j^{(i1)}=\hfill & \kappa _{j1}^{(i)}+\stackrel{~}{\kappa }_{j1}^{(i)}\hfill & (1jt^{(i1)});\hfill \\ \stackrel{~}{\kappa }_j^{(i1)}=\hfill & \stackrel{~}{\kappa }_{j1}^{(i)}\hfill & (1jt^{(i1)}).\hfill \end{array}$$
If $`i=t_{k(i)}`$ then:
$$\begin{array}{ccc}p^{(i1)}=\hfill & 2p^{(i)}p^{(i)};\hfill & \\ p^{(i1)}=\hfill & p^{(i)}p^{(i)};\hfill & \\ \kappa _j^{(i1)}=\hfill & 2\kappa _{j1}^{(i)}\stackrel{~}{\kappa }_{j1}^{(i)}\hfill & (2jt^{(i1)});\hfill \\ \stackrel{~}{\kappa }_j^{(i1)}=\hfill & \stackrel{~}{\kappa }_{j1}^{(i)}\stackrel{~}{\kappa }_{j1}^{(i)}\hfill & (2jt^{(i1)}).\hfill \end{array}$$
Proof: If $`it_{k(i)}`$ then $`k(i1)=k(i)`$. Then $`p^{(i)}/p^{(i)}`$ and $`p^{(i1)}/p^{(i1)}`$ have continued fractions $`(t_{k(i)}+1i,c_{k(i)+1},\mathrm{},c_n)`$ and $`(t_{k(i)}+2i,c_{k(i)+1},\mathrm{},c_n)`$ respectively. That $`p^{(i1)}=p^{(i)}+p^{(i)}`$ and $`p^{(i1)}=p^{(i)}`$ follows immediately. The expressions for $`\kappa _j^{(i1)}`$ and $`\stackrel{~}{\kappa }_j^{(i1)}`$ then follow from Lemma 6.1.
If $`i=t_{k(i)}`$ then $`k(i1)=k(i)1`$. Then $`p^{(i)}/p^{(i)}`$ and $`p^{(i1)}/p^{(i1)}`$ have continued fractions $`(c_{k(i)},c_{k(i)+1},\mathrm{},c_n)`$ and $`(2,c_{k(i)},c_{k(i)+1},\mathrm{},c_n)`$ respectively. That $`p^{(i1)}=2p^{(i)}p^{(i)}`$ and $`p^{(i1)}=p^{(i)}p^{(i)}`$ follows immediately. The expressions for $`\kappa _j^{(i1)}`$ and $`\stackrel{~}{\kappa }_j^{(i1)}`$ then follow from combining Lemma 6.2 with Lemma 6.1. $`\mathrm{}`$
As above, take $`A\{R,L\}`$. If $`a^A𝒯`$, set
$$\begin{array}{cc}a_i^A=\hfill & \{\begin{array}{cc}1\hfill & \text{if }\sigma ^Ai<t;\hfill \\ \kappa _{\sigma ^Ai}^{(i)}\hfill & \text{if }0i\sigma ^A,\hfill \end{array}\hfill \\ a_i^A=\hfill & \{\begin{array}{cc}1+\delta _{i,t_{k(i)}}\hfill & \text{if }\sigma ^Ai<t;\hfill \\ \kappa _{\sigma ^Ai+1}^{(i1)}\hfill & \text{if }0i<\sigma ^A,\hfill \end{array}\hfill \end{array}$$
and if $`a^A𝒯^{}`$, set
$$\begin{array}{cc}a_i^A=\hfill & \{\begin{array}{cc}p_i^{}1\hfill & \text{if }\sigma ^Ai<t;\hfill \\ p_i^{}\kappa _{\sigma ^Ai}^{(i)}\hfill & \text{if }0i\sigma ^A.\hfill \end{array}\hfill \\ a_i^A=\hfill & \{\begin{array}{cc}p_i^{}1\delta _{i,t_{k(i)}}\hfill & \text{if }\sigma ^Ai<t;\hfill \\ p_i^{}\kappa _{\sigma ^Ai+1}^{(i1)}\hfill & \text{if }0i<\sigma ^A.\hfill \end{array}\hfill \end{array}$$
In addition, define $`k^A`$ to be such that $`t_{k^A}<\sigma ^At_{k^A+1}`$. Then, if $`a^A𝒯`$, set
$$e_i^A=\{\begin{array}{cc}0\hfill & \text{if }\sigma ^Ai<t;\hfill \\ \delta _{k,k^A}^{(2)}\hfill & \text{if }0i<\sigma ^A,\hfill \end{array}$$
and if $`a^A𝒯^{}`$, set
$$e_i^A=\{\begin{array}{cc}1\hfill & \text{if }\sigma ^Ai<t;\hfill \\ 1\delta _{k,k^A}^{(2)}\hfill & \text{if }0i<\sigma ^A.\hfill \end{array}$$
###### Lemma 8.4
Let $`1i<t`$. Then for $`A\{L,R\}`$:
$$a_i^A=\{\begin{array}{cc}a_i^A+\frac{a_i^Ap_i}{p_i^{}}+e_i^A\hfill & \text{if }it_{k(i)};\hfill \\ 2a_i^A\frac{a_i^Ap_i}{p_i^{}}e_i^A\hfill & \text{if }i=t_{k(i)}.\hfill \end{array}$$
Proof: For $`p_i^{}/p_i`$, in view of the continued fraction specified above, the analogues of the quantities defined in (13) are $`t_j^{}=t_{k(i)+j}i`$ for $`1jnk(j)+1`$. For $`i<\sigma ^A`$, the various cases are then readily proved using Lemmas 6.3 and 8.3. For $`i\sigma ^A`$, the results follow immediately. $`\mathrm{}`$
For each $`t`$-dimensional vector $`𝒖=(u_1,u_2,\mathrm{},u_t)`$, define the $`(t1)`$-dimensional vector $`𝒖^{(\mathrm{},k)}=(u_1^{(\mathrm{},k)},u_2^{(\mathrm{},k)},\mathrm{},u_{t1}^{(\mathrm{},k)})`$ by
$$u_j^{(\mathrm{},k)}=\{\begin{array}{cc}0& \text{if }t_k^{}<jt_{k^{}+1},k^{}k(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ u_j& \text{if }t_k^{}<jt_{k^{}+1},k^{}k(\text{mod}\mathrm{\hspace{0.17em}2}),\hfill \end{array}$$
(25)
and the $`(t1)`$-dimensional vector $`𝒖^{(\mathrm{},k)}=(u_1^{(\mathrm{},k)},u_2^{(\mathrm{},k)},\mathrm{},u_{t1}^{(\mathrm{},k)})`$ by
$$u_j^{(\mathrm{},k)}=\{\begin{array}{cc}u_j& \text{if }t_k^{}<jt_{k^{}+1},k^{}k(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ 0& \text{if }t_k^{}<jt_{k^{}+1},k^{}k(\text{mod}\mathrm{\hspace{0.17em}2}),\hfill \end{array}$$
(26)
For convenience, we sometimes write $`𝒖_{(\mathrm{},k)}`$ instead of $`𝒖^{(\mathrm{},k)}`$, and $`𝒖_{(\mathrm{},k)}`$ instead of $`𝒖^{(\mathrm{},k)}`$.
Now for $`0it2`$, define:
$$\begin{array}{c}F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q)=\hfill \\ (q^{\frac{1}{4}\widehat{𝒎}^{(i+1)T}𝑪\widehat{𝒎}^{(i+1)}+\frac{1}{4}m_i^2\frac{1}{2}m_im_{i+1}\frac{1}{2}(𝒖_{(\mathrm{},k(i))}^L+𝒖_{(\mathrm{},k(i))}^R)𝒎^{(i)}+\frac{1}{4}\gamma _i^{\prime \prime }}\hfill \\ \underset{j=i+1}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}^{(i)}𝒖^L𝒖^R)_j}{m_j}\right]_q),\hfill \end{array}$$
(27)
where the sum here is taken over all vectors $`(m_{i+2},m_{i+3},\mathrm{},m_{t1})(Q_{i+2},Q_{i+3},\mathrm{},Q_{t1})`$, where $`(Q_1,Q_2,\mathrm{},Q_{t1})=𝑸(𝒖^L+𝒖^R)`$. The $`(t1)`$-dimensional vector $`𝒎^{(i)}=(0,0,\mathrm{},0,m_{i+1},m_{i+2},m_{i+3},\mathrm{},m_{t1})`$ has its first $`i`$ components equal to zero. The $`t`$-dimensional vector $`\widehat{𝒎}^{(i)}=(0,0,\mathrm{},0,m_i,m_{i+1},m_{i+2},\mathrm{},m_{t1})`$ has its first $`i`$ components equal to zero.
We also define:
$$F_{a,b}^{(t1)}(𝒖^L,𝒖^R,m_{t1},m_t;q)=q^{\frac{1}{4}m_{t1}^2+\frac{1}{4}\gamma _i^{\prime \prime }}\delta _{m_t,0}.$$
(28)
For convenience, we set $`Q_t=0`$.
Since $`\left[\genfrac{}{}{0pt}{}{m+n}{m}\right]_{q^1}=q^{mn}\left[\genfrac{}{}{0pt}{}{m+n}{m}\right]_q`$, it follows that for $`0it2`$:
$$\begin{array}{c}F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q^1)=\hfill \\ (q^{\frac{1}{4}\widehat{𝒎}^{(i+1)T}𝑪\widehat{𝒎}^{(i+1)}\frac{1}{4}m_i^2\frac{1}{2}(𝒖_{(\mathrm{},k(i)1)}^L+𝒖_{(\mathrm{},k(i)1)}^R)𝒎^{(i)}\frac{1}{4}\gamma _i^{\prime \prime }}\hfill \\ \underset{j=i+1}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}^{(i)}𝒖^L𝒖^R)_j}{m_j}\right]_q),\hfill \end{array}$$
(29)
where the sum here is taken over all vectors $`(m_{i+2},m_{i+3},\mathrm{},m_{t1})(Q_{i+2},Q_{i+3},\mathrm{},Q_{t1})`$, as above. Of course, we also have:
$$F_{a,b}^{(t1)}(𝒖^L,𝒖^R,m_{t1},m_t;q^1)=q^{\frac{1}{4}m_{t1}^2\frac{1}{4}\gamma _i^{\prime \prime }}\delta _{m_t,0}.$$
(30)
###### Lemma 8.5
Let $`0i<t`$, $`m_iQ_i`$ and $`m_{i+1}Q_{i+1}`$. If
$$𝒮^{(i)}=\{\begin{array}{cc}\left\{\kappa _{t_ni}^{(i)}\right\}\hfill & \text{if }i<t_n,\sigma ^L<t_n,\sigma ^R<t_n\text{ and }a,b𝒯;\hfill \\ \left\{p_i^{}\kappa _{t_ni}^{(i)}\right\}\hfill & \text{if }i<t_n,\sigma ^L<t_n,\sigma ^R<t_n\text{ and }a,b𝒯^{};\hfill \\ \mathrm{}\hfill & \text{otherwise},\hfill \end{array}$$
then:
$$\stackrel{~}{\chi }_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}(m_i,m_{i+1})\left\{𝒮^{(i)}\right\}=F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1}).$$
(31)
In addition, $`\alpha _{a_i^L,a_i^R}^{p_i,p_i^{}}=\alpha _i^{\prime \prime }`$ and $`\beta _{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}=\beta _i^{}`$.
Proof: For $`i=t1`$, we have $`p_i^{}=3`$, $`p_i=1`$, and if $`a^A𝒯`$ then $`a_i^A=1`$, $`e_i^A=0`$, $`(𝚫^A)_t=0`$; and if $`a^A𝒯^{}`$ then $`a_i^A=2`$, $`e_i^A=1`$, $`(𝚫^A)_t=1`$. Furthermore, we have $`it_n`$. Via (22) and (23), we obtain $`\alpha _{t1}^{\prime \prime }=\beta _{t1}^{}=(𝚫^L)_t(𝚫^R)_t`$ and $`\gamma _{t1}^{\prime \prime }=((𝚫^L)_t(𝚫^R)_t)^2`$. For $`i=t1`$, the first statement of our induction proposition is now seen to hold via Lemma 2.5. The definitions of $`\alpha _{a_i^L,a_i^R}^{p_i,p_i^{}}`$ and $`\beta _{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}`$ then yield the two final statements.
Now assume the result holds for a particular $`i`$ with $`1it1`$. As above, let $`k(i)`$ be such that $`t_{k(i)}i<t_{k(i)+1}`$. First consider the case $`t_{k(i)}<i<t_{k(i)+1}`$. Equation (24) gives $`\alpha _i=\alpha _i^{\prime \prime }`$, $`\beta _i=\beta _i^{}`$ and $`\gamma _i=\gamma _i^{\prime \prime }`$. Let $`m_{i1}Q_{i1}`$. On setting $`M=m_{i1}+u_i^L+u_i^R`$, equation (19) implies that $`MQ_{i+1}`$. Then, use of the induction hypothesis, Lemmas 3.14 or Lemma 6.7 as appropriate, and Lemmas 8.3 and 8.4 yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_{i1},p_{i1}^{}}(M,m_i)\left\{𝒮^{(i1)}\right\}\hfill \\ =\underset{m_{i+1}Q_{i+1}}{}q^{\frac{1}{4}(Mm_i)^2\frac{1}{4}\beta _i^2}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(M+m_{i+1})}{m_i}\right]_qF_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1}).\hfill \end{array}$$
(32)
Here, Lemma 3.5 also gives $`\alpha _{a_i^L,a_i^R}^{p_{i1},p_{i1}^{}}=\alpha _i+\beta _i`$, and $`\beta _{a_i^L,a_i^R,e_i^L,e_i^R}^{p_{i1},p_{i1}^{}}=\beta _i`$.
That $`\left\{𝒮^{(i1)}\right\}`$ appears on the leftside here is because, via Lemma 6.3, if $`i<t_n`$ then $`\kappa _{t_ni}^{(i)}`$ is interfacial in the $`(p_i,p_i^{})`$-model, and borders the $`\stackrel{~}{\kappa }_{t_ni}^{(i)}`$th odd band, and then $`\kappa _{t_ni}^{(i)}+\stackrel{~}{\kappa }_{t_ni}^{(i)}=\kappa _{t_ni+1}^{(i1)}`$, by Lemma 8.3, and finally noting that $`it_n`$ so that if $`it_n`$ then $`i1t_n`$. (The up-down symmetry of the $`(p,p^{})`$-model implies that if $`i<t_n`$ then $`p_i^{}\kappa _{t_ni}^{(i)}`$ is interfacial in the $`(p_i,p_i^{})`$-model, and borders the $`(p_i\stackrel{~}{\kappa }_{t_ni}^{(i)})`$th odd band. Then we use $`(p_i^{}\kappa _{t_ni}^{(i)})+(p_i\stackrel{~}{\kappa }_{t_ni}^{(i)})=p_{i1}^{}\kappa _{t_ni+1}^{(i1)}`$, from Lemma 8.3.)
Since $`M=m_{i1}+u_i^L+u_i^R`$, on noting that $`t_k<i<t_{k+1}`$, we have:
$$M+m_{i+1}=2m_i(\widehat{𝑪}\widehat{𝒎}^{(i1)}𝒖^L𝒖^R)_i,$$
and
$$\begin{array}{c}\widehat{𝒎}^{(i+1)T}𝑪\widehat{𝒎}^{(i+1)}+m_i^22m_im_{i+1}+(Mm_i)^2\hfill \\ =\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}+M^22Mm_i\hfill \\ =\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}+m_{i1}^22m_im_{i1}\hfill \\ +2(m_{i1}m_i)(u_i^L+u_i^R)+(u_i^L+u_i^R)^2.\hfill \end{array}$$
(In the case $`i=t1`$, we require this expression after substituting $`m_t=0`$.) Thence,
$$\begin{array}{c}\stackrel{~}{\chi }_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L+u_i^R,m_i)\left\{𝒮^{(i1)}\right\}\hfill \\ =(q^{\frac{1}{4}\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}\frac{1}{2}m_im_{i1}+\frac{1}{2}(m_{i1}m_i)(u_i^L+u_i^R)\frac{1}{2}(𝒖_{(\mathrm{},k(i))}^L+𝒖_{(\mathrm{},k(i))}^R)𝒎^{(i)}}\hfill \\ q^{\frac{1}{4}m_{i1}^2+\frac{1}{4}(u_i^L+u_i^R)^2+\frac{1}{4}\gamma _i\frac{1}{4}\beta _i^2}\underset{j=i}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}^{(i1)}𝒖^L𝒖^R)_j}{m_j}\right]_q),\hfill \end{array}$$
where the sum is over all $`(m_{i+1},m_{i+2},\mathrm{},m_{t1})(Q_{i+1},Q_{i+2},\mathrm{},Q_{t1})`$.
If $`i=\sigma ^R`$ then $`u_i^R=1`$. In this case, by definition, we have either $`a_i^R=1`$, $`e_i^R=0`$, $`a_{i1}^R=2`$ and $`e_{i1}^R=1`$, or $`a_i^R=p_{i1}^{}1`$, $`e_i^R=1`$, $`a_{i1}^R=p_{i1}^{}2`$ and $`e_{i1}^R=0`$. It is easily checked that $`a_i^R𝒮^{(i1)}`$. Then, use of Lemma 7.6 yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_i^L,a_{i1}^R,e_i^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L,m_i)\left\{𝒮^{(i1)}\right\}\hfill \\ =q^{\frac{1}{2}u_i^R(m_{i1}+u_i^L+u_i^R)+\frac{1}{2}(𝚫^R)_i(\alpha _i+\beta _i)}\hfill \\ \stackrel{~}{\chi }_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L+u_i^R,m_i)\{𝒮^{(i1)}\}.\hfill \end{array}$$
(33)
If $`i\sigma ^R`$ then (noting that $`it_k`$) $`u_i^R=(𝚫^R)_i=0`$, $`e_{i1}^R=e_i^R`$ and $`a_{i1}^R=a_i^R`$. The preceding expression thus also holds in this case.
We also immediately obtain
$$\begin{array}{cc}\alpha _{a_i^L,a_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\alpha _i+\beta _i(𝚫^R)_i;\hfill \\ \beta _{a_i^L,a_{i1}^R,e_i^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\beta _i(𝚫^R)_i.\hfill \end{array}$$
If $`i=\sigma ^L`$ then $`u_i^L=1`$. In this case, by definition, we have either $`a_i^L=1`$, $`e_i^L=0`$, $`a_{i1}^L=2`$ and $`e_{i1}^L=1`$, or $`a_i^L=p_{i1}^{}1`$, $`e_i^L=1`$, $`a_{i1}^L=p_{i1}^{}2`$ and $`e_{i1}^L=0`$. It is easily checked that $`a_i^L𝒮^{(i1)}`$. Then, use of Lemma 7.5 yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i1)}\right\}\hfill \\ =q^{\frac{1}{2}u_i^L(m_{i1}m_i+u_i^L)\frac{1}{2}(𝚫^L)_i(\beta _i(𝚫^R)_i)}\hfill \\ \stackrel{~}{\chi }_{a_i^L,a_{i1}^R,e_i^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L,m_i)\{𝒮^{(i1)}\}.\hfill \end{array}$$
(34)
If $`i\sigma ^R`$ then (noting that $`it_k`$) $`u_i^L=(𝚫^L)_i=0`$, $`e_{i1}^L=e_i^L`$ and $`a_{i1}^L=a_i^L`$. The preceding expression thus also holds in this case.
We also obtain:
$$\begin{array}{cc}\alpha _{a_{i1}^L,a_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\alpha _i+\beta _i(𝚫^R)_i+(𝚫^L)_i;\hfill \\ \beta _{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\beta _i(𝚫^R)_i+(𝚫^L)_i.\hfill \end{array}$$
Combining all the above, and using the expression for $`\gamma _{i1}^{\prime \prime }`$ given by (22) and (23), yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i1)}\right\}\hfill \\ =(q^{\frac{1}{4}\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}+\frac{1}{4}m_{i1}^2\frac{1}{2}m_{i1}m_i\frac{1}{2}(𝒖_{(\mathrm{},k(i))}^L+𝒖_{(\mathrm{},k(i))}^R)𝒎^{(i1)}+\frac{1}{4}\gamma _{i1}^{\prime \prime }}\hfill \\ \underset{j=i}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}^{(i1)}𝒖^L𝒖^R)_j}{m_j}\right]_q)\hfill \\ =F_{a,b}^{(i1)}(𝒖^L,𝒖^R,m_{i1},m_i),\hfill \end{array}$$
which is the required result when $`it_k`$, since $`k(i)=k(i1)`$.
In this $`it_k`$ case, making use of (22), (23), we also immediately obtain:
$$\begin{array}{cc}\alpha _{a_{i1}^L,a_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\alpha _{i1}^{\prime \prime };\hfill \\ \beta _{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\beta _{i1}^{}.\hfill \end{array}$$
Now consider the case for which $`i=t_k`$. Equation (24) gives $`\alpha _i=\alpha _i^{\prime \prime }`$, $`\beta _i=\alpha _i^{\prime \prime }\beta _i^{}`$ and $`\gamma _i=\alpha _i^2\gamma _i^{\prime \prime }`$. Corollary 4.2 gives $`\alpha _{a_i^L,a_i^R}^{p_i^{}p_i,p_i^{}}=\alpha _i`$ and $`\beta _{a_i^L,a_i^R,1e_i^L,1e_i^R}^{p_i^{}p_i,p_i^{}}=\beta _i`$. Let $`m_{i1}Q_{i1}`$. On setting $`M=m_{i1}+u_i^L+u_i^R`$, equation (19) implies that $`MQ_{i+1}`$. Then, use of the induction hypothesis, Lemmas 4.6 or Lemma 6.8 as appropriate, and Lemmas 8.3 and 8.4 yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_i^L,a_i^R,1e_i^L,1e_i^R}^{p_i,p_i^{}}(M,m_i;q)\left\{𝒮^{(i)}\right\}\hfill \\ =\underset{m_{i+1}Q_{i+1}}{}(q^{\frac{1}{4}(m_i^2+(Mm_i)^2\alpha _i^2\beta _i^2)}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(M+m_im_{i+1})}{m_i}\right]_q\hfill \\ F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q^1)),\hfill \end{array}$$
(35)
where
$$𝒮^{(i)}=\{\begin{array}{cc}\left\{\kappa _{t_ni+1}^{(i1)}\right\}\hfill & \text{if }i<t_n,\sigma ^L<t_n,\sigma ^R<t_n\text{ and }a,b𝒯;\hfill \\ \left\{p_i^{}\kappa _{t_ni+1}^{(i1)}\right\}\hfill & \text{if }i<t_n,\sigma ^L<t_n,\sigma ^R<t_n\text{ and }a,b𝒯^{};\hfill \\ \mathrm{}\hfill & \text{otherwise},\hfill \end{array}$$
using a similar argument to that in the $`it_{k(i)}`$ case. Here, Lemma 3.5 also gives $`\alpha _{a_i^L,a_i^R}^{p_{i1},p_{i1}^{}}=\alpha _i+\beta _i`$, and $`\beta _{a_i^L,a_i^R,1e_i^L,1e_i^R}^{p_{i1},p_{i1}^{}}=\beta _i`$.
Now set $`M=m_{i1}+u_i^L+u_i^R`$, whence on noting that $`i=t_k`$,
$$M+m_im_{i+1}=2m_i(\widehat{𝑪}\widehat{𝒎}^{(i1)}𝒖^L𝒖^R)_i$$
(in the case $`i=t1`$, we require this expression after substituting $`m_t=0`$), and
$$\begin{array}{c}\widehat{𝒎}^{(i+1)T}𝑪\widehat{𝒎}^{(i+1)}m_i^2+m_i^2+(Mm_i)^2\hfill \\ =\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}+M^22Mm_i\hfill \\ =\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}+m_{i1}^22m_im_{i1}\hfill \\ +2(m_{i1}m_i)(u_i^L+u_i^R)+(u_i^L+u_i^R)^2.\hfill \end{array}$$
Use of (29) or (30) then gives:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_i^L,a_i^R,1e_i^L,1e_i^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L+u_i^R,m_i)\left\{𝒮^{(i)}\right\}\hfill \\ =(q^{\frac{1}{4}\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}\frac{1}{2}m_im_{i1}+\frac{1}{2}(m_{i1}m_i)(u_i^L+u_i^R)\frac{1}{2}(𝒖_{(\mathrm{},k(i)1)}^L+𝒖_{(\mathrm{},k(i)1)}^R)𝒎^{(i)}}\hfill \\ q^{\frac{1}{4}m_{i1}^2+\frac{1}{4}(u_i^L+u_i^R)^2+\frac{1}{4}\gamma _i\frac{1}{4}\beta _i^2}\underset{j=i}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}^{(i1)}𝒖^L𝒖^R)_j}{m_j}\right]_q),\hfill \end{array}$$
where the sum is over all $`(m_{i+1},m_{i+2},\mathrm{},m_{t1})(Q_{i+1},Q_{i+2},\mathrm{},Q_{t1})`$.
Now set $`𝒮^{(i)R}=𝒮^{(i)}a_i^R`$ if $`i>\sigma ^R`$ and $`𝒮^{(i)R}=𝒮^{(i)}`$ otherwise.
Since $`i=t_k`$, it follows that $`u_i^R=1`$ if $`i>\sigma ^R`$. In this case, by definition, we have either $`a_i^R=2`$, $`1e_i^R=1`$, $`a_{i1}^R=1`$ and $`e_{i1}^R=0`$, or $`a_i^R=p_{i1}^{}2`$, $`1e_i^R=0`$, $`a_{i1}^R=p_{i1}^{}1`$ and $`e_{i1}^R=1`$. Then Lemma 7.4 yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_i^L,a_{i1}^R,1e_i^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L,m_i)\left\{𝒮^{(i)R}\right\}\hfill \\ =q^{\frac{1}{2}u_i^R(m_{i1}+u_i^L+u_i^R)+\frac{1}{2}(𝚫^R)_i(\alpha _i+\beta _i)}\hfill \\ \stackrel{~}{\chi }_{a_i^L,a_i^R,1e_i^L,1e_i^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L+u_i^R,m_i)\left\{𝒮^{(i)}\right\}.\hfill \end{array}$$
(36)
In addition, the same expression clearly also holds in the case $`i\sigma ^R`$, for which $`u_i^R=(𝚫^R)_i=0`$, $`e_{i1}^R=1e_i^R`$ and $`a_{i1}^R=a_i^R`$. (In the $`i=\sigma ^R`$ case, note that $`k(i1)=k(i)1=k^R(i)`$.)
Lemma 7.4 also implies that:
$$\begin{array}{cc}\alpha _{a_i^L,a_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\alpha _i+\beta _i(𝚫^R)_i;\hfill \\ \beta _{a_i^L,a_{i1}^R,e_i^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\beta _i(𝚫^R)_i.\hfill \end{array}$$
Now set $`𝒮^{(i)L}=𝒮^{(i)R}a_i^L`$ if $`i>\sigma ^L`$ and $`𝒮^{(i)L}=𝒮^{(i)R}`$ otherwise.
Since $`i=t_k`$, it follows that $`u_i^L=1`$ if $`i>\sigma ^L`$. In this case, by definition, we have either $`a_i^L=2`$, $`1e_i^L=1`$, $`a_{i1}^L=1`$ and $`e_{i1}^L=0`$, or $`a_i^L=p_{i1}^{}2`$, $`1e_i^L=0`$, $`a_{i1}^L=p_{i1}^{}1`$ and $`e_{i1}^R=1`$. Then Lemma 7.2 yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i)L}\right\}\hfill \\ =q^{\frac{1}{2}u_i^L(m_{i1}m_i+u_i^L)\frac{1}{2}(𝚫^L)_i(\beta _i(𝚫^R)_i)}\hfill \\ \stackrel{~}{\chi }_{a_i^L,a_i^R,1e_i^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1}+u_i^L,m_i)\left\{𝒮^{(i)R}\right\}.\hfill \end{array}$$
(37)
In addition, the same expression clearly also holds in the case $`i\sigma ^L`$, for which $`u_i^L=(𝚫^L)_i=0`$, $`e_{i1}^L=1e_i^L`$ and $`a_{i1}^L=a_i^L`$. (In the $`i=\sigma ^L`$ case, note that $`k(i1)=k(i)1=k^L(i)`$.)
Lemma 7.2 also implies that:
$$\begin{array}{cc}\alpha _{a_{i1}^L,a_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\alpha _i+\beta _i(𝚫^R)_i+(𝚫^L)_i;\hfill \\ \beta _{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\beta _i(𝚫^R)_i+(𝚫^L)_i.\hfill \end{array}$$
Combining all the above cases for $`i=t_k`$ yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i)L}\right\}.\hfill \\ =(q^{\frac{1}{4}\widehat{𝒎}^{(i)T}𝑪\widehat{𝒎}^{(i)}+\frac{1}{4}m_{i1}^2\frac{1}{2}m_{i1}m_i\frac{1}{2}(𝒖_{(\mathrm{},k(i)1)}^L+𝒖_{(\mathrm{},k(i)1)}^R)𝒎^{(i1)}+\frac{1}{4}\gamma _{i1}^{\prime \prime }}\hfill \\ \underset{j=i}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}^{(i1)}𝒖^L𝒖^R)_j}{m_j}\right]_q)\hfill \\ =F_{a,b}^{(i1)}(𝒖^L,𝒖^R,m_{i1},m_i).\hfill \end{array}$$
Once it is established that
$$𝒫_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i)L}\right\}=𝒫_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i1)}\right\}.$$
we obtain the required result when $`i=t_k`$, since $`k(i)=k(i1)+1`$.
If $`i=t_n`$ then $`\{𝒮^{(i)L}\}=\{𝒮^{(i1)}\}`$ immediately. Now let $`i<t_n`$. For $`A\{L,R\}`$, if $`\sigma _i^A=1`$ then necessarily $`\sigma _{t_n}^A=1`$. In the case that $`a^A𝒯`$, this implies that $`\{2,\kappa _{t_ni}^{(i)}\}𝒮^{(i)L}`$ and $`\kappa _{t_ni}^{(i)}𝒮^{(i1)}`$. Since $`a_{i1}^A=1`$, we may drop the element $`2`$ from $`𝒮^{(i)L}`$ with no effect. Similar reasoning holds for $`a^A𝒯^{}`$ whereupon the claim is established.
In this $`i=t_k`$ case, making use of (22), (23), we also immediately obtain:
$$\begin{array}{cc}\alpha _{a_{i1}^L,a_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\alpha _{i1}^{\prime \prime };\hfill \\ \beta _{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}\hfill & =\beta _{i1}^{}.\hfill \end{array}$$
The lemma then follows by induction. $`\mathrm{}`$
Before performing a sum over $`m_1`$, we require the following result.
###### Lemma 8.6
For $`0jt`$,
$$\begin{array}{cc}\alpha _j^{\prime \prime }\hfill & Q_j(\text{mod}\mathrm{\hspace{0.17em}2});\hfill \\ \beta _j^{}\hfill & Q_jQ_{j+1}(\text{mod}\mathrm{\hspace{0.17em}2}).\hfill \end{array}$$
Proof: Since $`\alpha _t^{\prime \prime }=0`$, $`\beta _t^{}=0`$ and $`Q_t=Q_{t+1}=0`$, this result is manifest for $`j=t`$.
We now proceed by downward induction. Thus assume the result holds for a particular $`j>0`$. When $`jt_{k(j)}`$, equations (24) and (22) imply that $`\beta _{j1}^{}=\beta _j^{}+(𝒖^L)_j(𝒖^R)_j`$. Equation (19) implies that $`Q_{j1}Q_{j+1}(𝒖^L)_j(𝒖^R)_j`$. Thus the induction hypothesis immediately gives $`\beta _{j1}^{}Q_{j1}Q_j`$ in this case.
When $`j=t_{k(j)}`$, equations (24) and (22) imply that $`\beta _{j1}^{}=\alpha _j^{\prime \prime }\beta _j^{}+(𝒖^L)_j(𝒖^R)_j`$. Equation (19) implies that $`Q_{j1}Q_j+Q_{j+1}(𝒖^L)_j(𝒖^R)_j`$. Thus the induction hypothesis also gives $`\beta _{j1}^{}Q_{j1}Q_j`$ in this case.
In both cases, equations (24), (22) and (23) give $`\alpha _{j1}^{\prime \prime }=\alpha _j^{\prime \prime }+\beta _{j1}^{\prime \prime }`$, whence the induction hypothesis immediately gives $`\alpha _j^{\prime \prime }Q_{j1}`$ as required. $`\mathrm{}`$
Define:
$$\begin{array}{c}F_{a,b}(𝒖^L,𝒖^R,L;q)\hfill \\ =\underset{𝒎𝑸(𝒖^L+𝒖^R)}{}q^{\frac{1}{4}\widehat{𝒎}^T𝑪\widehat{𝒎}\frac{1}{4}L^2\frac{1}{2}(𝒖_{\mathrm{}}^L+𝒖_{\mathrm{}}^R)𝒎+\frac{1}{4}\gamma }\underset{j=1}{\overset{t1}{}}\left[\genfrac{}{}{0pt}{}{m_j\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}𝒖^L𝒖^R)_j}{m_j}\right]_q.\hfill \end{array}$$
The summation here is over all vectors $`𝒎=(m_1,m_2,\mathrm{},m_{t1})`$ such that $`m_j\mathrm{}_0`$ and $`m_jQ_j(\text{mod}\mathrm{\hspace{0.17em}2})`$ for $`1j<t`$. Then, $`\widehat{𝒎}=(m_0,m_1,m_2,\mathrm{},m_{t1})`$.
On defining
$$𝒮=\{\begin{array}{cc}\left\{\kappa _i\right\}\hfill & \text{if }\sigma ^L<t_n,\sigma ^R<t_n\text{ and }a,b𝒯;\hfill \\ \left\{p_i^{}\kappa _i\right\}\hfill & \text{if }\sigma ^L<t_n,\sigma ^R<t_n\text{ and }a,b𝒯^{};\hfill \\ \mathrm{}\hfill & \text{otherwise},\hfill \end{array}$$
we then obtain:
###### Lemma 8.7
Let $`p^{}>2p`$. If $`L\alpha _{a,b}^{p,p^{}}`$ then
$$\stackrel{~}{\chi }_{a,b,e_0^L,e_0^R}^{p,p^{}}(L)\left\{𝒮\right\}=F_{a,b}(𝒖^L,𝒖^R,L).$$
In addition, $`\delta _{b,e_0^R}^{p,p^{}}=0`$.
Proof: Lemma 8.6 implies that $`LQ_0`$. Lemma 2.3 requires the sum over all $`m_1L+\beta _{a,b,e,f}^{p,p^{}}`$ of the $`i=0`$ case of Lemma 8.5. This is applicable since for such $`m_1`$, Lemma 8.6 implies that $`m_1Q_1`$.
The lemma follows after noting that in the $`p^{}>2p`$ case, $`\widehat{𝒎}^{(1)T}C\widehat{𝒎}^{(1)}+L^22Lm_1=\widehat{𝒎}^TC\widehat{𝒎}L^2`$ and $`\gamma _0^{\prime \prime }=\gamma `$. $`\mathrm{}`$
We now transfer this result to the original weighting function of (3). To do this we require the value of $`c`$ given by (14). Then, defining $`\chi _{a,b,c}^{p,p^{}}(L)\left\{𝒮\right\}`$ in the way analogous to $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L)\left\{𝒮\right\}`$, we obtain:
###### Lemma 8.8
If $`L\alpha _{a,b}^{p,p^{}}(\text{mod}\mathrm{\hspace{0.17em}2})`$ then
$$\chi _{a,b,c}^{p,p^{}}(L)\left\{𝒮\right\}=F_{a,b}(𝒖^L,𝒖^R,L).$$
Proof: For the moment, assume that $`p^{}>2p`$. Consider $`h𝒫_{a,b,e,f}^{p,p^{}}(L)`$ and $`h^{}𝒫_{a,b,c^{}}^{p,p^{}}(L)`$ given by $`h_i^{}=h_i`$ for $`0iL`$. If $`\delta _{b,f}^{p,p^{}}=0`$ and $`c^{}=b+(1)^f`$ then, as noted in Section 2, $`\stackrel{~}{wt}(h)=wt(h^{})`$. Consequently, $`\stackrel{~}{\chi }_{a,b,e,f}^{p,p^{}}(L)\left\{𝒮\right\}=\chi _{a,b,c^{}}^{p,p^{}}(L)\left\{𝒮\right\}`$. However, if $`b`$ is interfacial then the same is true for $`c^{}=b\pm 1`$. As noted at the end of Section 6.1, $`b`$ is interfacial if $`\sigma ^Rt_1`$. Otherwise, the current lemma follows from noting that for the $`c`$ defined above, $`c=b+(1)^{e_0^R}`$.
Now given $`h𝒫_{a,b,c}^{p,p^{}}(L)`$, define $`\widehat{h}𝒫_{a,b,c}^{p^{}p,p^{}}(L)`$ by by $`\widehat{h}_i=h_i`$ for $`0iL`$. As in Lemma 4.1, $`\mathrm{wt}(\widehat{h})=\frac{1}{4}(L^2\alpha ^2)\mathrm{wt}(h)`$, where $`\alpha =\alpha _{a,b}^{p,p^{}}`$. Therefore $`\chi _{a,b,c}^{p,p^{}}(L)\left\{𝒮\right\}=q^{\frac{1}{4}(L^2\alpha ^2)}\chi _{a,b,c}^{p,p^{}}(L;q^1)\left\{𝒮\right\}`$. Since $`\alpha _{a,b}^{p,p^{}}=\alpha _0^{\prime \prime }`$ by Lemma 8.5, and $`\gamma _0=(\alpha _0^{\prime \prime })^2\gamma _0^{\prime \prime }`$ by (24), the $`p^{}<2p`$ case follows from the $`p^{}>2p`$ case obtained above after using $`\left[\genfrac{}{}{0pt}{}{m+n}{m}\right]_{q^1}=q^{mn}\left[\genfrac{}{}{0pt}{}{m+n}{m}\right]_q`$, and noting the change in the definition of $`𝑪`$. $`\mathrm{}`$
Proof of Theorem 8.1: First consider the case where $`a<y_n`$ and $`b<y_n`$. Then necessarily $`a,b𝒯`$. Since $`y_n=\kappa _{t_n}`$, we have $`\sigma ^L<t_n`$ and $`\sigma ^R<t_n`$. Thereupon, $`𝒮=\{y_n\}`$. Let $`h𝒫_{a,b,c}^{p,p^{}}(L)\backslash 𝒫_{a,b,c}^{p,p^{}}(L)\{y_n\}`$. Then $`1h_i<y_n`$ for $`0iL`$. Since, by Lemma 6.4, the lowermost $`y_n2`$ bands of the $`(p,p^{})`$-model have exactly the same parities as the corresponding bands of the $`(z_n,y_n)`$-model, we see that if $`h^{}𝒫_{a,b,c}^{z_n,y_n}(L)`$ is defined by $`h_i^{}=h_i`$ for $`0iL`$ then $`wt(h^{})=wt(h)`$. Since all of $`𝒫_{a,b,c}^{z_n,y_n}(L)`$ arises in this way, we have $`\chi _{a,b,c}^{p,p^{}}(L)=\chi _{a,b,c}^{p,p^{}}(L)\{y_n\}+\chi _{a,b,c}^{z_n,y_n}(L)`$. This proves the first case of Theorem 8.1.
The second case arises if $`a>p^{}y_n`$ and $`b>p^{}y_n`$. Here, necessarily $`a,b𝒯^{}`$, whence again $`\sigma ^L<t_n`$ and $`\sigma ^R<t_n`$. The argument proceeds as above, noting that both the $`(p,p^{})`$\- and $`(z_n,y_n)`$-models are up-down symmetric.
The other cases are immediate since $`𝒮=\mathrm{}`$. $`\mathrm{}`$
### 8.3 The $`𝒎`$$`𝒏`$-system
Each term in the fermionic expressions given by Theorem 8.1 or Theorem 8.2 corresponds to a vector $`𝒎=(m_1,m_2,\mathrm{},m_{t1})`$ where $`𝒎𝑸(𝒖^L+𝒖^R)`$. As usual, we set $`\widehat{𝒎}=(L,m_1,m_2,\mathrm{},m_{t1})`$. Now, for each $`𝒎`$, define a vector $`𝒏=(n_1,n_2,\mathrm{},n_t)`$ by
$$𝒏=\frac{1}{2}(\widehat{𝑪}\widehat{𝒎}+𝒖).$$
(38)
In view of (19), we see that $`n_j\mathrm{}`$ for $`1jt`$. Then since
$$\frac{1}{2}(𝑪\widehat{𝒎}𝒖^L𝒖^R)_j=n_j,$$
(39)
in those terms that provide a non-zero contribution to the fermionic expression of Theorem 8.1, $`n_j0`$ for $`1jt`$.
On examining the proof of Lemma 8.5, we see that $`n_i`$ is the number of particles added at the $`i`$th induction step to pass from $`𝒫_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}(m_i,m_{i+1})\left\{𝒮^{(i)}\right\}`$ to $`𝒫_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{𝒮^{(i1)}\right\}`$.
The set of equations that link the two vectors $`\widehat{𝒎}`$ and $`𝒏`$ is known as the $`𝒎`$$`𝒏`$-system. On account of (18), the equations are more explicitly given by, for $`1jt`$:
$`m_{j1}m_{j+1}=m_j+2n_ju_j\text{if }j=t_k,k=1,2,\mathrm{},n\text{;}`$ (40)
$`m_{j1}+m_{j+1}=2m_j+2n_ju_j\text{ otherwise,}`$ (41)
where we set $`m_t=m_{t+1}=0`$.
Using these two expressions, and setting $`m_0=L`$, it may be readily shown that:
$$\underset{i=1}{\overset{t}{}}l_in_i=\frac{1}{2}\left(L+\underset{i=1}{\overset{t}{}}l_iu_i\right).$$
(42)
Thereupon, the summands in the expression for $`F_{a,b}(𝒖^L,𝒖^R,L)`$ given in Theorem 8.1 correspond to solutions of (42) with each $`n_i`$ a non-negative integer.
### 8.4 The second fermionic form
The proof of Theorem 8.2 follows the same lines as that of Theorem 8.1. We will not give the full description, but indicate how the proof of Lemma 8.5 is affected by the use of the modified Gaussians. We first define $`F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q)`$ for $`0i<t`$ in the same way as $`F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q)`$ in (27) and (28), except employing the modified Gaussians instead of the classical Gaussians. Note that this modified form of the Gaussian differs from the form defined in (1) if and only if $`A<0`$ and $`B0`$. In this case, $`\left[\genfrac{}{}{0pt}{}{A}{B}\right]=0`$. In addition, since $`\left[\genfrac{}{}{0pt}{}{m+n}{m}\right]_{q^1}^{}=q^{mn}\left[\genfrac{}{}{0pt}{}{m+n}{m}\right]_q^{}`$, it follows that the analogues of (29) and (30) hold.
###### Lemma 8.9
Let $`0i<t`$, $`m_iQ_i`$ and $`m_{i+1}Q_{i+1}`$. If $`m_i0`$ then:
$$\stackrel{~}{\chi }_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}(m_i,m_{i+1})=F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1}).$$
(43)
In addition, $`\alpha _{a_i^L,a_i^R}^{p_i,p_i^{}}=\alpha _i^{\prime \prime }`$ and $`\beta _{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}=\beta _i^{}`$.
Proof: The proof proceeds much as in the proof of 8.5. However, we must certainly check that using the modified Gaussians does not introduce unwanted terms.
Consider the $`it_{k(i)}`$ case. Combining the analogues of (32), (33) and (34) yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\hfill \\ =\underset{\genfrac{}{}{0pt}{}{m_{i+1}Q_{i+1}}{0m_{i+1}m_i+1}}{}q^{\frac{1}{2}\left(m_iu_i^Lm_{i1}(u_i^L+u_i^R)u_i^Lu_i^R2+\beta _i((𝚫^R)_i(𝚫^L)_i)+\alpha _i(𝚫^R)_i+(𝚫^L)_i(𝚫^R)_i\right)}\hfill \\ \times q^{\frac{1}{4}(Mm_i)^2\frac{1}{4}\beta _i^2}\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(M+m_{i+1})}{m_i}\right]_qF_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1}),\hfill \end{array}$$
where $`M=m_{i1}+u_i^L+u_i^R`$. Since $`m_{i1},m_{i+1}0`$, and $`u_i^L,u_i^R0`$ (because $`it_{k(i)}`$), we have
$$\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_{i1}+m_{i+1}+u_i^L+u_i^R)}{m_i}\right]_q^{}=\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_{i1}+m_{i+1}+u_i^L+u_i^R)}{m_i}\right]_q.$$
(44)
The induction step for $`it_{k(i)}`$ then proceeds exactly as in the proof of Lemma 8.5.
For the $`i=t_{k(i)}`$ case, combining the analogues of (35), (36) and (37) yields:
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{\stackrel{~}{𝒮}\right\}\hfill \\ =\underset{\genfrac{}{}{0pt}{}{m_{i+1}Q_{i+1}}{0m_{i+1}m_i+1}}{}q^{\frac{1}{2}\left(m_iu_i^Lm_{i1}(u_i^L+u_i^R)u_i^Lu_i^R+(𝚫^L)_i(𝚫^R)_i\right)1}\hfill \\ \times q^{\frac{1}{2}\left(\beta _i((𝚫^R)_i(𝚫^L)_i)+\alpha _i(𝚫^R)_i\right)+\frac{1}{4}\left(m_i^2+(Mm_i)^2\alpha _i^2\beta _i^2\right)}\hfill \\ \times \left[\genfrac{}{}{0pt}{}{\frac{1}{2}(M+m_im_{i+1})}{m_i}\right]_qF_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q^1),\hfill \end{array}$$
(45)
where $`M=m_{i1}+u_i^L+u_i^R`$, and $`2\stackrel{~}{𝒮}`$ if and only if either $`a_i^L=1`$ or $`a_i^R=1`$; $`p^{}2\stackrel{~}{𝒮}`$ if and only if either $`a_i^L=p^{}1`$ or $`a_i^R=p^{}1`$; and $`\stackrel{~}{𝒮}`$ contains no other values.
We must check that (45) holds if the Gaussian is replaced by its modified form, and the ‘$`\{\stackrel{~}{𝒮}\}`$’ is removed.
If $`u_i^L=u_i^R=0`$ then $`\stackrel{~}{𝒮}=\mathrm{}`$. In addition $`m_{i+1}m_i+1`$ implies that:
$$\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_{i1}+m_im_{i+1}+u_i^L+u_i^R)}{m_i}\right]_q^{}=\left[\genfrac{}{}{0pt}{}{\frac{1}{2}(m_{i1}+m_im_{i+1}+u_i^L+u_i^R)}{m_i}\right]_q.$$
(46)
Thereupon, the induction step for this subcase of $`i=t_{k(i)}`$ follows as in the proof of Lemma 8.5.
Now consider $`u_i^Lu_i^R`$. We tackle the case $`u_i^L=0`$ and $`u_i^R=1`$ (the case $`u_i^L=1`$ and $`u_i^R=0`$ is similar). This implies that $`\sigma ^Lt_{k(i)}`$ and $`\sigma ^R<t_{k(i)}`$. Then either $`a_{i1}^R=1`$ and $`\stackrel{~}{𝒮}=\{2\}`$, or $`a_{i1}^R=p^{}1`$ and $`\stackrel{~}{𝒮}=\{p^{}2\}`$. In addition, $`2a_{i1}^Lp^{}2`$. We immediately see that
$$\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i)\left\{\stackrel{~}{𝒮}\right\}=\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(m_{i1},m_i).$$
(47)
On the other hand, since $`m_{i+1}m_i+1`$, (46) is valid here unless $`m_{i1}=m_i=0`$ and $`m_{i+1}=1`$. Now $`\sigma ^Lt_{k(i)}`$ implies that if $`a_i^L=a_i^R`$ then $`\sigma ^L=t_{k(i)}`$ and $`e_i^L=e_i^R`$ whereupon $`F_{a,b}^{(i)}(𝒖^L,𝒖^R,0,1;q^1)=0`$. In this case, since $`a_{i1}^La_{i1}^R`$, then $`\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(0,0)=0`$. Therefore, the induction step holds in this $`u_i^Lu_i^R`$ case.
Now consider $`u_i^L=u_i^R=1`$, so that $`\sigma ^L<t_{k(i)}`$ and $`\sigma ^R<t_{k(i)}`$. If $`a^A𝒯`$ then $`a_{i1}^A=1`$, and if $`a^A𝒯^{}`$ then $`a_{i1}^A=p^{}1`$. Thereupon, (47) holds unless $`m_{i1}=m_i=0`$ and either both $`a,b𝒯`$ or both $`a,b𝒯^{}`$. In these cases,
$$\begin{array}{c}\stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(0,0)\left\{\stackrel{~}{𝒮}\right\}=0;\hfill \\ \stackrel{~}{\chi }_{a_{i1}^L,a_{i1}^R,e_{i1}^L,e_{i1}^R}^{p_{i1},p_{i1}^{}}(0,0)=1,\hfill \end{array}$$
(48)
by direct enumeration. On the other hand, (46) is valid here unless $`m_{i1}+m_im_{i+1}=0`$, and $`m_i=0`$. If $`m_{i1}=m_i=0`$ then since $`\left[\genfrac{}{}{0pt}{}{1}{0}\right]^{}=1`$, and $`\alpha _i=\beta _i=0`$, the required analogue of (45) holds in this case. If $`m_{i1}=1`$ and $`m_i=0`$ then both sides of the analogue of (45) are easily seen to be zero.
The induction step is now complete, whence the lemma follows. $`\mathrm{}`$
Note that, at the $`i`$th step in the induction, an extra term arises due to the modified Gaussian only if $`i=t_{k(i)}`$, $`\sigma ^L<i`$, $`\sigma ^R<i`$ and either both $`a,b𝒯`$ or both $`a,b𝒯^{}`$. In this case, consider the term $`F_{a,b}^{(i)}(𝒖^L,𝒖^R,m_i,m_{i+1};q^1)`$, in (45) which enumerates the elements of $`𝒫_{a_i^L,a_i^R,e_i^L,e_i^R}^{p_i,p_i^{}}(m_i,m_{i+1})`$. In the case where the extra term arises, $`m_i=m_{i+1}=0`$ and either both $`a_i^L=a_i^R=1`$ and $`e_i^L=e_i^R=0`$, or both $`a_i^L=a_i^R=p^{}1`$ and $`e_i^L=e_i^R=1`$. Thus there is precisely one path $`\stackrel{~}{h}`$ of zero length.
Equation (45) encapsulates the action of a $`𝒟`$-transform followed by a $`(k,\lambda )`$-transform on $`\stackrel{~}{h}`$, followed by extending the result on both sides (since $`u_i^L=u_i^R=1`$). We thus obtain a path of length $`m_{i1}=2k+2`$ in the $`(p_{i1},p_{i1}^{})`$-model. This path has the form given in Fig. 9.
That this path contains $`n_i=k`$ particles, is also encoded in (40).
When the classical Gaussians are employed, equation (45) thus fails to account for the case of a zero length path. Use of the modified Gaussian remedies this, by permitting the case $`n_i=1`$. This may be viewed as an annihilation of the $`k=0`$ case of Fig. 9, which although appearing to be a particle (c.f. Lemma 3.12), arises through solely the action of the $`B_1`$-transform followed by path extension.
Acknowledgments: We would like to thank Professor Y. Pugai for collaboration on an earlier stage of this work, and on related works, and for many useful discussions. His contributions to this work are gratefully acknowledged. We also wish to thank Professors A. Berkovich, B. McCoy and A. Schilling for many informative discussions on . Finally, we wish to thank Professors M. Kashiwara and T. Miwa for the invitation to attend ‘Physical Combinatorics’ where a preliminary version of this work was presented, and for their excellent hospitality. |
warning/0002/astro-ph0002145.html | ar5iv | text | # On radial gas flows, the Galactic Bar and chemical evolution in the Galactic Disc
## 1 Introduction
The observed chemical and spectro–photometric properties of galaxies are one of the main sources of information for our understanding of galaxy formation and evolution. The corresponding theoretical modelling involves star formation (SF) as a basic ingredient. Unfortunately, this process is rather poorly known on the large scales relevant to galaxy evolution. Portinari & Chiosi (1999, hereinafter PC99) analysed the effects of adopting different SF laws in a chemical model for the Galactic Disc, a system which we can study in great detail. In this paper we address another phenomenon which can bear interesting effects on the chemical evolution of galaxies: radial gas flows. A few papers in literature (Section 2) demonstrate that radial gas flows influence chemical models for the Disc, especially in their predictions on the metallicity gradient. It is therefore interesting to discuss the radial profile of the Disc with models including also radial flows, in addition to various options for the SF law. In particular, radial flows can help to overcome some difficulties that “static” models find in reproducing, at the same time, the metallicity gradient and the radial gas profile of the Disc (PC99).
We develop a chemical model with radial gas flows as a multi–dimensional extension of the model of Portinari et al. (1998, hereinafter PCB98). The model is described in Section 3 and in the appendices. In Section 4 we discuss the general qualitative effects of superposing radial flows upon a chemical model. In Section 5 we present models for the Galactic Disc with radial gas flows and different SF laws, showing that radial flows provide an alternative or additional dynamical effect to the “inside–out” formation scenario to explain the metallicity gradient. Section 6 is dedicated to qualitative simulations of the dynamical effects of the Galactic Bar upon the gas distribution, with the aim to reproduce the molecular ring around 4 kpc, which static models cannot account for (PC99). Section 7 contains a final summary and conclusions.
## 2 Radial flows: previous literature
The possibility that radial flows play a role in establishing the radial metallicity gradients in galactic discs was first suggested by Tinsley & Larson (1978). Following Lacey & Fall (1985), we mention that radial gas flows in a disc can be driven by three main mechanisms:
1. the infalling gas has a lower angular momentum than the circular motions in the disc, and mixing with the gas in the disc induces a net radial inflow with a velocity up to a few km sec<sup>-1</sup>;
2. viscosity in the gas layer induces radial inflows in the inner parts of the disc and outflows in the outer parts, with velocities of $``$0.1 km sec<sup>-1</sup>;
3. gravitational interactions between gas and spiral density waves lead to large-scale shocks, dissipation and therefore radial inflows of gas (or outflows in the outer parts) with typical velocities of $``$0.3 km sec<sup>-1</sup> (e.g. Bertin & Lin 1996 and references therein); much larger velocities can be achieved in the inner few kpc in the presence of a barred potential.
In summary, radial flows are plausible with velocities of $``$0.1–1 km sec<sup>-1</sup>, and they are expected to be inflows over most of the disc. Observational upper limits permit radial inflows in the Galactic Disc with velocities up to 5 km sec<sup>-1</sup> at the present time. For further details, see Lacey & Fall (1985) and references therein.
The first of the above mentioned mechanisms was modelled in detail by Mayor & Vigroux (1980), and later by Pitts & Tayler (1989, 1996), Chamcham & Tayler (1994). The effects of a generic inflow velocity profile on chemical evolution models has been explored by Lacey & Fall (1985), Tosi(1988), Götz & Köppen (1992), Köppen (1994), Edmunds & Greenhow (1995). A different approach is that of viscous disc models which, rather than imposing arbitrary radial velocity patterns, describe the evolution of the gas distribution in the disc self–consistently, following the model suggested by Lin & Pringle (1987). Viscous chemical models have been developed by Clarke (1989), Yoshii & Sommer-Larsen (1989) and Sommer-Larsen & Yoshii (1990), Thon & Meusinger (1998). All these studies show how radial inflows can steepen the metallicity gradients with respect to static models, especially if an outer cut–off of SF is assumed.
## 3 Modelling radial flows
We formulate our chemical model with radial flows as a multi-dimensional extension of the static model of PCB98 and PC99, an open model where the disc forms gradually by accretion of protogalactic gas. The disc is divided in $`N`$ concentric rings or shells; in each ring $`k`$ the gaseous component and its chemical abundances evolve due to:
1. depletion by SF, which locks up gas into stars;
2. stellar ejecta which shed back enriched material to the interstellar medium (ISM);
3. infall of primordial protogalactic gas;
4. gas exchange with the neighbouring rings because of radial flows.
The set of equations driving the chemical evolution of the $`k`$-th shell is:
$$\begin{array}{cc}\frac{d}{dt}G_i(r_k,t)=\hfill & X_i(r_k,t)\mathrm{\Psi }(r_k,t)+\hfill \\ & \\ & +_{M_l}^{M_u}\mathrm{\Psi }(r_k,t\tau _M)R_i(M)\mathrm{\Phi }(M)𝑑M+\hfill \\ & \\ & +\left[\frac{d}{dt}G_i(r_k,t)\right]_{inf}+\hfill \\ & \\ & +\left[\frac{d}{dt}G_i(r_k,t)\right]_{rf}\hfill \end{array}$$
(1)
where the various symbols are defined here below.
Primordial gas is accreted at an exponentially decreasing rate with time-scale $`\tau `$:
$$\dot{\sigma }_{inf}(r_k,t)=A(r_k)e^{\frac{t}{\tau (r_k)}}$$
(2)
$`A(r_k)`$ is obtained by imposing that the integrated contribution of infall up to the present Galactic age $`t_G=`$15 Gyr, corresponds to an assumed exponential profile $`\sigma _A(r_k)`$:
$$A(r_k)=\frac{\sigma _A(r_k)}{\tau (r_k)(1e^{t_G/\tau (r_k)})}=\frac{\sigma _A(r_{})e^{\frac{r_kr_{}}{r_d}}}{\tau (r_k)(1e^{t_G/\tau (r_k)})}$$
(3)
Indicating with $`\sigma _g(r_k,t)`$ the surface gas density, we define the gas fraction:
$$G(r_k,t)=\frac{\sigma _g(r_k,t)}{\sigma _A(r_k)}$$
(4)
and the normalized surface gas density for each chemical species $`i`$:
$$G_i(r_k,t)=X_i(r_k,t)G(r_k,t)$$
(5)
where $`X_i`$ is the fractionary abundance by mass of $`i`$.
The 1<sup>st</sup> term on the right-hand side of Eq. (1) represents the depletion of species $`i`$ from the ISM due to star formation; see PC99 for the various options concerning the SF rate, $`\mathrm{\Psi }(r,t)`$. The 2<sup>nd</sup> term is the amount of species $`i`$ ejected back to the ISM by dying stars; the returned fractions $`R_i(M)`$ are calculated on the base of the detailed stellar yields from PCB98 and keep track of finite stellar lifetimes (no instantaneous recycling approximation IRA). The 3<sup>rd</sup> term is the contribution of infall, while the 4<sup>th</sup> term describes the effect of radial flows. Full details on the first three terms can be found in the original static model by PCB98 and PC99. The novelty in Eq. (1) is the radial flow term, which we develop here below. We will adopt the simplified notation $`\sigma _{gk}\sigma _g(r_k,t)`$ and the like.
Let the $`k`$-th shell be defined by the galactocentric radius $`r_k`$, its inner and outer edge being labelled as $`r_{k\frac{1}{2}}`$ and $`r_{k+\frac{1}{2}}`$. Through these edges, gas flows with velocity $`v_{k\frac{1}{2}}`$ and $`v_{k+\frac{1}{2}}`$, respectively (Fig. 1). Flow velocities are taken positive outward; the case of inflow is correspondingly described by negative velocities.
Radial flows through the borders, with a flux $`F(r)`$, contribute to alter the gas surface density in the $`k`$-th shell according to:
$$\left[\frac{d\sigma _{gk}}{dt}\right]_{rf}=\frac{1}{\pi \left(r_{k+\frac{1}{2}}^2r_{k\frac{1}{2}}^2\right)}\left[F(r_{k+\frac{1}{2}})F(r_{k\frac{1}{2}})\right]$$
(6)
The gas flux at $`r_{k+\frac{1}{2}}`$ can be written:
$$F\left(r_{k+\frac{1}{2}}\right)=2\pi r_{k+\frac{1}{2}}v_{k+\frac{1}{2}}\left[\chi \left(v_{k+\frac{1}{2}}\right)\sigma _{gk}+\chi \left(v_{k+\frac{1}{2}}\right)\sigma _{g\left(k+1\right)}\right]$$
(7)
where $`\chi (x)`$ is the step function: $`\chi (x)=1`$ or 0 for $`x>`$ or $`0`$, respectively. Eq. (7) is a sort of “upwind approximation” for the advection term to be included in the model equations (e.g. Press et al. 1986), describing either inflow or outflow depending on the sign of $`v_{k+\frac{1}{2}}`$. An analogous expression holds for $`F(r_{k\frac{1}{2}})`$.
Let’s take the inner edge $`r_{k\frac{1}{2}}`$ at the midpoint between $`r_{k1}`$ and $`r_k`$, and similarly for $`r_{k+\frac{1}{2}}`$ (Fig. 1). Writing Eq. (6) separately for each chemical species $`i`$, in terms of the $`G_i`$’s we obtain the radial flow term of Eq. (1) as:
$$\begin{array}{cc}\left[\frac{d}{dt}G_i(r_k,t)\right]_{rf}=\hfill & \alpha _kG_i(r_{k1},t)\beta _kG_i(r_k,t)+\hfill \\ & +\gamma _kG_i(r_{k+1},t)\hfill \end{array}$$
(8)
where:
$$\begin{array}{cc}\alpha _k=\hfill & \frac{2}{r_k+\frac{r_{k1}+r_{k+1}}{2}}\left[\chi (v_{k\frac{1}{2}})v_{k\frac{1}{2}}\frac{r_{k1}+r_k}{r_{k+1}r_{k1}}\right]\frac{\sigma _{A(k1)}}{\sigma _{Ak}}\hfill \\ & \\ \beta _k=\hfill & \frac{2}{r_k+\frac{r_{k1}+r_{k+1}}{2}}\times \hfill \\ \multicolumn{2}{c}{\times \left[\chi (v_{k\frac{1}{2}})v_{k\frac{1}{2}}\frac{r_{k1}+r_k}{r_{k+1}+r_{k1}}\chi (v_{k+\frac{1}{2}})v_{k+\frac{1}{2}}\frac{r_k+r_{k+1}}{r_{k+1}r_{k1}}\right]}\\ & \\ \gamma _k=\hfill & \frac{2}{r_k+\frac{r_{k1}+r_{k+1}}{2}}\left[\chi (v_{k+\frac{1}{2}})v_{k+\frac{1}{2}}\frac{r_k+r_{k+1}}{r_{k+1}r_{k1}}\right]\frac{\sigma _{A(k+1)}}{\sigma _{Ak}}\hfill \end{array}$$
(9)
The terms on the right-hand side of Eq. (8) evidence the contribution of the 3 contiguous shells involved: the first term represents the gas being gained in shell $`k`$ from $`k`$–1, the second term is the gas being lost from $`k`$ to $`k`$–1 and $`k`$+1, and the third term is the gas being gained in $`k`$ from $`k`$+1. The coefficients (9) are all $`0`$ and depend only on the shell $`k`$, not on the chemical species $`i`$ considered. If the velocity pattern is constant in time, $`\alpha _k`$, $`\beta _k`$ and $`\gamma _k`$ are also constant in time.
Notice that in the case of static models the final surface mass density is completely determined by the assumed accretion profile, namely $`\sigma (r_k,t_G)\sigma _A(r_k)`$. Therefore, in static models the radial profile for accretion can be directly chosen so as to match the observed present–day surface density in the Disc (see PCB98 and PC99). The inclusion of the term of radial gas flows alters the expected final density profile and $`\sigma (r_k,t_G)\sigma _A(r_k)`$. Hence, $`\sigma (r,t_G)`$ cannot be assumed in advance and is only known a posteriori (see §4); at the end of each simulation we need to check how much radial flows have altered the actual density profile $`\sigma (r_k,t_G)`$ with respect to the pure accretion profile $`\sigma _A(r_k)`$. With the slow flow speeds considered ($`v\stackrel{<}{}1`$ km sec<sup>-1</sup>), the two profiles will not be too dissimilar anyways.
### 3.1 Boundary conditions
Eq. (8) needs to be slightly modified in the case of the innermost and the outermost shell, since the shell $`k`$–1 or $`k`$+1 are not defined in these two respective cases.
#### 3.1.1 The innermost shell
Our disc models will extend down to where the Bulge becomes the dominating Galactic component ($`r_1=2.5`$ kpc). As to the innermost edge, we assume that the first shell is symmetric with respect to $`r_1`$:
$$r_{\frac{1}{2}}=\frac{3r_1r_2}{2}$$
and that $`v_{\frac{1}{2}}0`$ always, since we cannot account for outflows from still inner shells, not included in the model. For $`k=1`$, Eq. (8) then becomes:
$$\left[\frac{d}{dt}G_i(r_1,t)\right]_{rf}=\beta _1G_i(r_1,t)+\gamma _1G_i(r_2,t)$$
(10)
with:
$$\beta _1=\frac{1}{2r_1}\left[v_{\frac{1}{2}}\frac{3r_1r_2}{r_2r_1}\chi (v_{\frac{3}{2}})v_{\frac{3}{2}}\frac{r_1+r_2}{r_2r_1}\right]$$
$$\gamma _1=\chi (v_{\frac{3}{2}})v_{\frac{3}{2}}\frac{1}{2r_1}\frac{r_1+r_2}{r_2r_1}\frac{\sigma _{A2}}{\sigma _{A1}}$$
(11)
#### 3.1.2 Boundary conditions at the disc edge
As to the outermost shell ($`k=N`$), we need a boundary condition for the gas inflowing from the outer disc. We assume a SF cut-off in the outer disc, while the gaseous layer extends much further. In fact, in external spirals HI discs are observed to extend much beyond the optical disc, out to 2 or even 3 optical radii. A threshold preventing SF beyond a certain radius is expected from gravitational stability in fluid discs (Toomre 1964, Quirk 1972) and has observational support as well (Kennicutt 1989). For the Galactic Disc we assume no SF beyond the last shell at $`r_N=20`$ kpc, which is the empirical limit for the optical disc and for HII regions and bright blue stars tracing active SF. Gas can though flow in from the outer disc; extended gas discs might actually provide a much larger gas reservoir for star–forming spirals than vertical infall, at least at the present time when the gravitational settling of the protogalactic cloud is basically over.
If no SF occurs in the outer disc, the evolution of the gas (and total) surface density can be expressed as:
$$\frac{\sigma }{t}(r,t)=A(r)e^{\frac{t}{\tau (r)}}\frac{1}{r}\frac{}{r}(rv\sigma )r>r_{N+\frac{1}{2}}$$
(12)
(e.g. Lacey & Fall 1985, their model equation with the SF term dropped). Here, with no SF, $`\sigma \sigma _g`$ and abundances always remain the primordial ones ($`X_{i,inf}`$). Let’s assume the following simplifying conditions for the outer disc:
1. the infall time-scale is uniform:
$$\tau (r)\tau (r_N)r>r_{N+\frac{1}{2}}$$
2. the inflow velocity is uniform and constant:
$$v(r,t)v_{N+\frac{1}{2}}r>r_{N+\frac{1}{2}},t$$
3. the infall profile is flat:
$$A(r)A_{ext}r>r_{N+\frac{1}{2}}$$
in accordance with observed extended gas discs in spirals, showing a much longer scale-length than the stellar component.
With these assumptions, Eq. (12) becomes:
$$\frac{\sigma }{t}+v\frac{\sigma }{r}=Ae^{\frac{t}{\tau }}\frac{v}{r}\sigma $$
(13)
where we indicate $`\tau \tau (r_N)`$, $`vv_{N+\frac{1}{2}}`$ and $`AA_{ext}`$ to alleviate the notation. Eq. (13) has a straightforward analytical solution (Appendix B):
$$\begin{array}{cc}\sigma (r,t)=\hfill & A\tau \times \hfill \\ & \\ \multicolumn{2}{c}{\times \left[\left(1e^{\frac{t}{\tau }}\right)+\frac{v}{r}\left(\tau \left(e^{\frac{T_{rf}}{\tau }}e^{\frac{t}{\tau }}\right)(tT_{rf})\right)\right]}\end{array}$$
(14)
where $`T_{rf}0`$ is the time when radial inflows are assumed to activate. Eq. (14) is our boundary condition at the outermost edge.
Notice that (14) is the solution of (13) in the idealized case of an infinite, flat gas layer extending boundless to any $`r>r_N`$ (see also Appendix B). Of course, this does not correspond to gaseous discs surrounding real spirals; but since we will consider only slow inflow velocities ($`v\stackrel{<}{}1`$ km sec<sup>-1</sup>), with typical values of $`r_N=20`$ kpc and $`t_G=15`$ Gyr, the gas actually drifting into the model disc shells will be just the gas originally accreted within $`r35`$ kpc. Therefore, the boundary condition (14) remains valid as long as the gas layer stretches out to $`35`$ kpc, a very plausible assumption since observed gaseous discs extend over a few tens or even $`100`$ kpc.
#### 3.1.3 The outermost shell
We take a reference external radius $`r_{ext}>r_N`$ in the outer disc where the (total and gas) surface density $`\sigma (r_{ext},t)\sigma _{ext}(t)`$ is given by the boundary condition (14); typically, $`r_N=20`$ kpc and $`r_{ext}21`$ kpc. We take the outer edge of the shell at the midpoint:
$$r_{N+\frac{1}{2}}=\frac{r_N+r_{ext}}{2}$$
and replace
$$X_{i(k+1)}\sigma _{g(k+1)}X_{i,inf}\sigma _{ext}$$
in Eqs. (6) and (7), since the primordial abundances $`X_{i,inf}`$ remain unaltered in the outer disc, in the absence of SF. We thus write the radial flow term for the $`N`$-th shell as:
$$\begin{array}{cc}\left[\frac{d}{dt}G_i(r_N,t)\right]_{rf}=\hfill & \alpha _NG_i(r_{N1},t)\beta _NG_i(r_N,t)+\hfill \\ & +\omega _i(t)\hfill \end{array}$$
(15)
where:
$$\begin{array}{cc}\omega _i=\hfill & X_{i,inf}\chi (v_{N+\frac{1}{2}})v_{N+\frac{1}{2}}\times \hfill \\ & \times \frac{4}{r_{N1}+2r_N+r_{ext}}\frac{r_N+r_{ext}}{r_{ext}r_{N1}}\frac{\sigma _{ext}(t)}{\sigma _{AN}}\hfill \end{array}$$
(16)
### 3.2 The numerical solution
Using (8), (10) and (15), the basic set of equations (1) can be written as:
$$\{\begin{array}{cc}\frac{d}{dt}G_i(r_1,t)=\hfill & \vartheta _1(t)G_i(r_1,t)+\gamma _1G_i(r_2,t)+W_i(r_1,t)\hfill \\ & \\ \frac{d}{dt}G_i(r_k,t)=\hfill & \alpha _kG_i(r_{k1},t)+\vartheta _k(t)G_i(r_k,t)+\hfill \\ & +\gamma _kG_i(r_k,t)+W_i(r_k,t)\hfill \\ \multicolumn{2}{c}{k=2,\mathrm{}N1}\\ & \\ \frac{d}{dt}G_i(r_N,t)=\hfill & \alpha _NG_i(r_{N1},t)+\vartheta _N(t)G_i(r_N,t)+\hfill \\ & +W_i(r_N,t)+\omega _i(t)\hfill \end{array}$$
where we have introduced:
$$\begin{array}{ccc}\vartheta _k(t)\hfill & & \left(\eta (r_k,t)+\beta _k\right)0\hfill \\ & & \\ \eta (r_k,t)\hfill & & \frac{\mathrm{\Psi }}{G}(r_k,t)\hfill \\ & & \\ W_i(r_k,t)\hfill & & _{M_l}^{M_u}\mathrm{\Psi }(r_k,t\tau _M)R_i(M)\mathrm{\Phi }(M)𝑑M+\hfill \\ & & +\left[\frac{d}{dt}G_i(r_k,t)\right]_{inf}\hfill \end{array}$$
We refer to PCB98 for further details on the quantities $`\eta `$ and $`W_i`$, appearing also in the original static model. Neglecting, for the time being, that the $`\eta `$’s and the $`W_i`$’s contain the $`G_i`$’s themselves, we are dealing with a linear, first order, non homogeneous system of differential equations with non constant coefficients, of the kind:
$$\frac{dG_i}{dt}=𝒜(t)G_i(t)+W_i(t)$$
(17)
There is a system (17) for each chemical species $`i`$, but the matrix of the coefficients $`𝒜(t)`$ is independent of $`i`$.
We solve the system by the same numerical method used for the original equation of the static model — see Talbot & Arnett (1971) and PCB98 for details. We just need to extend the method to the present multi–dimensional case (17). If we consider the evolution of the $`G_i`$’s over a short enough timestep $`t_1t_0=\mathrm{\Delta }t`$, the various quantities $`\eta (r_k,t)`$, $`\vartheta _k(t)`$ and $`W_i(r_k,t)`$ will remain roughly constant within $`\mathrm{\Delta }t`$; similarly to the method for the static model (see PCB98), within $`\mathrm{\Delta }t`$ we approximate them with the values $`\overline{\eta }_k`$, $`\overline{\vartheta }_k`$ and $`\overline{W}_i(r_k)`$ they assume at the midstep $`t_{\frac{1}{2}}=t_0+\frac{1}{2}\mathrm{\Delta }t`$. Over $`\mathrm{\Delta }t`$, (17) can then be considered a system with constant coefficients $`𝒜𝒜(t_{\frac{1}{2}})`$, and $`W_i(t)`$ becomes a constant vector, which allows for the analytical solution:
$$G_i(t_1)=e^{\mathrm{\Delta }t𝒜}G_i(t_0)+\left[_{t_0}^{t_1}e^{(t_1t)𝒜}𝑑t\right]W_i$$
(18)
where $`e^{t𝒜}`$ indicates the matrix:
$$e^{t𝒜}=\left(\begin{array}{ccc}e^{\lambda _1t}u_1& \mathrm{}& e^{\lambda _Nt}u_N\end{array}\right)\left(\begin{array}{ccc}u_1& \mathrm{}& u_N\end{array}\right)^1$$
(19)
with $`\lambda _k`$ the eigenvalues of $`𝒜`$ and $`u_k`$ the corresponding eigenvectors. The matrix $`e^{t𝒜}`$ and the explicit expression of the solution (18) for the $`G_i`$’s are calculated in Appendix A, resulting in:
$$\begin{array}{cc}G_i(r_k,t_1)=\hfill & \alpha _k\times \hfill \\ \multicolumn{2}{c}{\times \frac{e^{(\overline{\eta }_k+\beta _k)\mathrm{\Delta }t}e^{(\overline{\eta }_{k1}+\beta _{k1})\mathrm{\Delta }t}}{(\overline{\eta }_{k1}+\beta _{k1})(\overline{\eta }_k+\beta _k)}G_i(r_{k1},t_0)+}\\ & +e^{(\overline{\eta }_k+\beta _k)\mathrm{\Delta }t}G_i(r_k,t_0)+\hfill \\ \multicolumn{2}{c}{+\gamma _k\frac{e^{(\overline{\eta }_{k+1}+\beta _{k+1})\mathrm{\Delta }t}e^{(\overline{\eta }_k+\beta _k)\mathrm{\Delta }t}}{(\overline{\eta }_k+\beta _k)(\overline{\eta }_{k+1}+\beta _{k+1})}G_i(r_{k+1},t_0)+}\\ & +\alpha _k\frac{\overline{W}_i(r_{k1})}{(\overline{\eta }_{k1}+\beta _{k1})(\overline{\eta }_k+\beta _k)}\times \hfill \\ \multicolumn{2}{c}{\times \left(\frac{1e^{(\overline{\eta }_k+\beta _k)\mathrm{\Delta }t}}{\overline{\eta }_k+\beta _k}\frac{1e^{(\overline{\eta }_{k1}+\beta _{k1})\mathrm{\Delta }t}}{\overline{\eta }_{k1}+\beta _{k1}}\right)+}\\ & +\overline{W}_i(r_k)\frac{1e^{(\overline{\eta }_k+\beta _k)\mathrm{\Delta }t}}{\overline{\eta }_k+\beta _k}+\hfill \\ & +\gamma _k\frac{\overline{W}_i(r_{k+1})}{(\overline{\eta }_k+\beta _k)(\overline{\eta }_{k+1}+\beta _{k+1})}\times \hfill \\ & \times \left(\frac{1e^{(\overline{\eta }_{k+1}+\beta _{k+1})\mathrm{\Delta }t}}{\overline{\eta }_{k+1}+\beta _{k+1}}\frac{1e^{(\overline{\eta }_k+\beta _k)\mathrm{\Delta }t}}{\overline{\eta }_k+\beta _k}\right)\hfill \end{array}$$
(20)
where we intend $`\alpha _10`$, $`\gamma _N0`$, and for the outermost shell one should replace:
$$\overline{W}_i(r_N)\overline{W}_i(r_N)+\overline{\omega }_i=\overline{W}_i(r_N)+\frac{1}{\mathrm{\Delta }t}_{t_0}^{t_1}\omega _i(t)𝑑t$$
$`\omega _i`$ being defined by (16).
Similarly to the case of an isolated shell (see PCB98) the system (20) does not provide the final solution for $`G_i(r_k,t_1)`$, since the $`\overline{\eta }_k`$’s and the $`\overline{W}_i`$’s on the right–hand side actually depend on the $`G_i`$’s; it just represents a set of implicit non–linear expressions in $`G_i(r_k,t_1)`$. As in the original static model, we can neglect the dependence of $`\overline{W}_i(r_k)`$ on the $`G_i`$’s and consider only that of $`\overline{\eta }_k`$ (see PCB98 for a detailed discussion). We must finally find the roots of the system (20) by applying the Newton-Raphson method, generalized to many dimensions (cfr. Press et al., 1986). Such a system holds for each of the chemical species $`i`$ considered, so at each iteration we actually need to solve as many systems as the species included in the model. Full details on the mathematical development of the model and its numerical solution can be found in the appendices and in Portinari (1998).
We tested the code against suitable analytical counterparts and obtained the following conditions for model consistency (Appendix B).
1. Rather small timesteps are needed for the numerical model to keep stable; the required timesteps get smaller and smaller the higher the flow velocities considered, and the thinner the shells.
2. To describe gas flows in a disc with an exponential density profile, the shells should be equispaced in the logarithmic, rather than linear, scale (so that they roughly have the same mass, rather than the same width).
We modelled the Galactic Disc using 35 shells from 2.5 to 20 kpc, equally spaced in the logarithmic scale, their width ranging from $`0.2`$ kpc for the inner shells to $`1`$ kpc for the outermost ones. With such a grid spacing, and velocities up to $`1`$ km sec <sup>-1</sup>, suitable timesteps are of $`10^4`$ Gyr (Appendix B; see also Thon & Meusinger 1998). This means that roughly $`1.5\times 10^5`$ timesteps, times 35 shells, are needed to complete each model, which would translate in excessive computational times. This drawback was avoided by separating the time-scales in the code.
1. The timestep $`\mathrm{\Delta }t`$ used to update the “chemical” variables ($`\eta `$, $`W_i`$, etc.) is the minimum among: $`\mathrm{\Delta }t_1`$ which guarantees that the relative variations of the $`G_i`$’s are lower than a fixed $`ϵ`$; $`\mathrm{\Delta }t_2`$ which guarantees that the total surface mass density $`\sigma (t)`$ increases by no more than 5%; $`\mathrm{\Delta }t_3`$ which is twice the previous timestep of the model, to speed up the computation when possible; $`\mathrm{\Delta }t_4`$ which guarantees the Courant condition $`\mathrm{\Delta }t<v\mathrm{\Delta }r`$, indispensable for the stability of a numerical algorithm describing flows. So, $`\mathrm{\Delta }t`$ is basically set by the requirement that the chemical quantities do not vary too much within it, and it can get relatively large (up to 0.2 Gyr), especially at late ages when the various chemical variables evolve slowly.
2. It is only the numerical solution (20) which needs very short timesteps to keep stable. Therefore, once the chemical variables are upgraded, the main timestep $`\mathrm{\Delta }t`$ is subdivided in much shorter timesteps $`\delta t=10^4`$ Gyr, upon which the solution (20) and its Newton-Raphson iteration are successively applied to cover the whole $`\mathrm{\Delta }t`$. Only then a new upgrade of all the $`\overline{\eta }_k`$’s and $`\overline{W}_i`$’s is performed.
This trick keeps the code roughly as fast as if it would evolve with a single time-scale $`\mathrm{\Delta }t`$, and yet it gives the same results as the “slow” version in which all quantities are upgraded at each $`\delta t=10^4`$ Gyr.
## 4 Effects of radial inflows on the disc
After presenting our new model with radial flows, we investigate their general effects on the chemical evolution of the disc. Here we first analyse how the predicted metallicity gradient and gas and total surface density profiles of a static model are altered by a superimposed uniform gas inflow, while in Section 5 we will suitably combine radial flows with various SF laws to match the observed radial profile of the Disc. The characteristics and parameters of the various models presented in this and in the next section are summarized in Table 1.
To gain a qualitative understanding of the effects of radial gas inflows on chemical evolution, we impose a uniform inflow pattern $`v=1`$ km sec<sup>-1</sup> upon a static chemical model and compare the new outcome with the original static case. Radial flows over the disc are expected to be mainly inflows, with velocities from 0.1 to $``$1 km sec<sup>-1</sup> (see §2); imposing a uniform inflow of 1 km sec<sup>-1</sup> is therefore a sort of “extreme case”, considered here for the sake of qualitative analysis. Anyways, previous studies in literature suggest that the effects of radial flows saturate for much higher velocities (Köppen 1994). We will consider both the case of inflow from the outer gaseous disc and not ($`v_{N+\frac{1}{2}}=1`$ or 0, respectively).
All models with radial flows are rescaled so that the final surface density at the Solar ring ($`r_{}=8.5`$ kpc) corresponds to 50 $`M_{}`$ pc<sup>-2</sup>. Namely, with radial flows $`\sigma (r_{},t_G)\sigma _A(r_{})`$ and it cannot be imposed as an input datum (see §3), but the zero–point $`A(r_{})`$ of the exponential accretion profile (3)
$$A(r)=A(r_{})e^{\frac{rr_{}}{r_d}}$$
can be rescaled so that at the end of the simulation $`\sigma (r_{},t_G)=50`$ $`M_{}`$ pc<sup>-2</sup>. This zero–point does not influence the profile, nor the chemical evolution, so it can always be rescaled a posteriori.
For our example, we take as the reference static case a model adopting a Schmidt SF law with $`\kappa =1.5`$ (Kennicutt 1998) and a uniform infall time-scale of 3 Gyr (model S15a of PC99, see also Table 1).
### 4.1 Models with inflow from the outer disc
In model S15RFa a uniform radial inflow pattern with $`v=1`$ km sec<sup>-1</sup> and inflow from the outer disc is imposed upon the static model S15a, with no further change in the model parameters (Table 1). Fig. 2 shows the effects of inflow on the total surface density and gas density distribution, and on the oxygen gradient (details on the observational data in the plots can be found in PC99). By comparing the solid line and the dashed line, we notice the following main effects.
* Since matter flows inward and accumulates toward the Galactic Centre, the density profile gets steeper in the inner parts, while in the outer parts it remains rather flat because gas is continuously poured in by the flat outer gaseous disc. The gas density distribution shows a similar behaviour.
* The overall metallicity gets much lower because of the dilution by primordial gas inflowing from the outer disc. The gradient becomes steeper, especially in the outermost shells, since the discontinuity in metallicity between the star–forming disc and the outer gaseous layer is smeared inward by radial inflows.
To compensate for the steepening of the density distribution induced by radial inflows, we must adopt a shallower initial accretion profile; our simulations show that a scale length $`r_d5`$ kpc for the infall profile reduces in the end to a density profile matching the desired scale length of $``$4 kpc at the Solar Neighbourhood. At the same time, the chemical enrichment must get more efficient for the overall metallicity to increase to the observed levels; the predicted metallicity is improved by adopting an IMF more weighted towards massive stars, i.e. by increasing the “IMF scaling fraction” $`\zeta `$ (see PCB98 and PC99 for a description of our model parameters). Together with the SF efficiency $`\nu `$, $`\zeta `$ is re-calibrated to match the observed gas surface density and metallicity at the Solar Neighbourhood (as for the models of PC99). In this way we calibrate model S15RFb with respect to the Solar Neighbourhood; see Table 1 for details on the adopted parameters. Comparing now this re-calibrated model with radial flows (dotted line) to the original model S15a, the metallicity gradient is increased with respect to the static case, but still a bit flat with respect to observations. The gas density distribution peaks in the inner regions, as expected, and remains quite high (much higher than observed) in the outer regions due to substantial replenishment from the outer disc.
### 4.2 Models with no inflow from the outer disc
Let’s now consider the case when radial flows are limited within the star forming disc and there is no inflow from the external gaseous disc ($`v_{N+\frac{1}{2}}=0`$). If radial flows are mainly driven by shocks in spiral arms, for instance, they might indeed occur only within the stellar disc, while the outer gas layer remains unperturbed.
Fig. 3 (left panels) shows how the final density profile becomes much steeper than the reference accretion profile (model S15RFc, dashed line), as expected since matter is efficiently drifting inward with no replenishment from the outer disc. For the final density profile to match the observed one, we must assume a much shallower initial accretion profile ($`r_d67`$ kpc). Then the resulting local density profile is close to the observed one, while in the outer parts the profile remains steeper (model S15RFd, dotted line).
The gas density profile shows a similar behaviour, strongly peaked toward the centre while dropping quickly (much more quickly than observed) outside the Solar ring (upper right panel).
The overall metallicity is reduced with respect to the original model S15a (lower right panel, model S15RFc) and again we need to increase the IMF scaling fraction $`\zeta `$ to raise the chemical enrichment to the observed values (model S15RFd). The resulting gradient is roughly comparable to the observed one.
### 4.3 Concluding remarks
We can summarize the general effects of radial flows as follows.
* Since matter flows inward and accumulates toward the Galactic Centre, the density profile in the inner parts gets steeper. A shallower intrinsic accretion profile is to be adopted, in order to recover the observed local scale–length at the end of the simulation.
* In the outer parts, the profile declines more or less sharply, depending on whether the inflow occurs only within the star–forming disc or there is also substantial inflow from the outer, purely gaseous disc.
* The gas distribution shows a similar behaviour: it remains quite flat in the outer regions if gas is poured in by the flat outer gas disc, while it tends to drop sharply (more than observed) otherwise.
* The inner gas profile is very steep in the case of a Schmidt SF law (models S15RF) — while with other SF laws with radially decaying efficiency (see PC99 and §5) the effect would be somewhat compensated by a larger gas consumption by SF in the inner regions.
* The overall metallicity gets much lower because of the dilution by gas inflowing from metal poor outer shells (and possibly the primordial outer disc). A higher fraction of stars contributing to the chemical enrichment is needed in the model to match the observed metallicity.
* The metallicity gradient tends to steepen, in agreement with results by other authors (see references in §2).
## 5 Some “successful” models
In §4 we have illustrated the qualitative effects of radial inflows using, for the sake of example, models with a Schmidt SF law. We will now consider models with various SF laws (see PC99) and tune the inflow velocity pattern, in each case, so as to match the observational data on the radial profile of the Galactic Disc. The various “successful” models presented here should not be taken as detailed, unique recipes to reproduce the Disc. Rather, they are meant as examples of how inclusion of radial flows in the chemical model can improve the match with the data.
### 5.1 Models with Schmidt law
Let’s first consider again the Schmidt SF law. Inspection of models S15RFb and S15RFd (with and without radial inflows from the outer disc) suggests to proceed as follows.
* Some inflow from the outer gaseous layer is needed to reproduce the shallow decline of the gas distribution observed out of the Solar ring, though the inflow speed from the outer disc should be slower than –1 km sec<sup>-1</sup> otherwise the predicted gas density is too high in the outer parts (as in model S15RFb).
* Models with radial inflows can predict metallicity gradients close to the observed ones even with a Schmidt SF law, provided drift velocities within the star–forming edge are relatively high (of the order of –1 km sec<sup>-1</sup>). Inflow patterns decelerating inward are particularly efficient in building the metallicity gradient (Götz & Köppen 1992).
An example of a drift velocity profile shaped according to these considerations is shown in Fig. 4; the corresponding model gives indeed a reasonable fit to the data (model S15RFe, Fig. 5). The gas profile keeps increasing inward, yet the model does not reproduce the peak of the molecular ring, for which we need to include the peculiar radial flows induced by the Bar (see the discussion in PC99 and §6). Model S15RFe shows how radial inflows can in fact allow for metallicity gradients comparable with the observed ones, even with a Schmidt SF law which would be excluded on the sole base of static models (PC99).
### 5.2 Models with spiral–triggered SF laws
Let’s consider now models adopting an Oort–type SF law, with a SF efficiency inversely proportional to the galactocentric radius and a Schmidt–like exponent $`\kappa =1.0`$ (Kennicutt 1998, PC99). The static and the “successful” model with this SF law are model O10a from PC99 and model O10RFe, respectively (Table 1).
As in the static case, (see PC99), these models behave somewhat similarly to models adopting the Schmidt SF law. Fig. 6 shows that a good fit is obtained with model O10RFe, where we applied the velocity pattern of Fig. 4 with inflows becoming slower inwar, very close to that used for model S15RFe. Notice that in this case the higher SF efficiency in the inner region and the slowdown of inflows conspire to accumulate gas around $`r=3`$ kpc while consuming it at inner radii, generating a peak in the predicted gas distribution which closely reminds the observed molecular ring.
### 5.3 Models with gravitational self–regulating SF laws
We now investigate models adopting a self–regulating SF process driven by gravitational settling and feed–back from massive stars, implying a SF efficiency exponentially decaying outward in radius, such as the law by Dopita & Ryder (1994; see PC99). The reference static model in this case is model DRa of PC99, while the corresponding “successful” model with radial flows is DRRFe (see Table 1 for the relevant parameters).
To reproduce the observed metallicity gradient with this SF law, negligible inflow is required in the outer regions where static models already predict a gradient with the right slope (see model DRa), while a moderate inflow velocity ($``$0.3 km sec<sup>-1</sup>) in the inner parts is needed to maintain the observed slope also where the predicted gradient would otherwise flatten (see PC99). Such a model is, for example, model DRRFe, obtained with the inflow velocity pattern shown in Fig. 7. In Fig. 8 the “successful” model DRRFe is compared to the original model DRa of PC99 with no radial flows.
### 5.4 Concluding remarks
The aim of the previous “successful” examples is to remark how, for each SF law considered, it is possible to tune the inflow velocity pattern so as to get a good overall agreement with the observational data. As an indication, suitable inflow profiles have the following characteristics.
* With a Schmidt SF law, relatively fast (but still plausible) flows are needed within the star–forming disc, with velocities of the order of 1 km sec<sup>-1</sup>. The required metallicity gradient is easily obtained, especially if the inflow velocity is slowly decreasing inward.
Inflow from the outer gaseous disc must be less prominent (though present), otherwise the predicted gas density in the outer regions is too high.
* In models adopting an Oort–type SF law with $`\kappa =1.0`$, a good fit to the data is obtained with an inflow profile of the same kind: inflow at $`0.5`$ km sec<sup>-1</sup> from the outer gaseous disc, drift velocities raising at $`1`$ km sec<sup>-1</sup> at the stellar disc edge and declining inward. Curiously enough, in this case the higher SF efficiency in the inner regions combines with the slowdown of inflows to predict a peak in the gas distribution around $`r=3`$ kpc, quite reminiscent of the observed molecular ring.
* Models adopting the SF law by Dopita & Ryder give a good description of the outer regions of the Disc already in the absence of flows (PC99). Mild radial inflows can be assumed in the inner regions, where the radial gradient would otherwise flatten, to obtain the right slope throughout the disc. The required drift velocities are around 0.3 km sec<sup>-1</sup>, which can be reasonably provided, for instance, by the shocks occurring within the spiral arms (see §2).
The various “successful models” presented here are not meant to be definitive recipes to reproduce the radial profile of the Galactic Disc. In fact, the adopted inflow velocity profiles are quite arbitrary and ad hoc. These models are rather meant to show how radial flows can be a viable mechanism to interpret the properties of the Disc. In PC99 we showed how none of the various SF laws investigated is able, by itself, to reproduce the observed metallicity gradient throughout the whole extent of the Disc, unless some additional “dynamical” assumption is included. In PC99 we considered the classical case of an inside–out disc formation, namely of an infall time-scale increasing outward. Here, we just show that radial inflows can provide another viable “dynamical” assumption to be combined with any SF law to reproduce the observed gradient. In particular, if we adopt a Schmidt SF law with $`\kappa =1.5`$ or an Oort–type SF law with $`\kappa =1.0`$, as recent empirical evidence seems to support (Kennicutt 1998), the required variation of the accretion time-scale is too extreme in the pure inside–out assumption (PC99), and radial inflows are then necessary to explain the metallicity gradient. Of course, the two effects (inside–out formation and radial flows) can also play at the same time. We do not address here their combined outcome, because in our models this would merely translate into increasing the number of parameters which one can tune to fit the observational data. No further insight in the problem would be gained. With this kind of models we can just learn how the various players (different SF laws, radially varying accretion time-scales, radial gas inflows) enter the game of reconstructing the general picture, and analyse their effects one by one.
It is also worth stressing that even slow radial gas flows, with velocities well plausible in terms of the triggering physical mechanisms and within the observational limits (§2), have non–negligible influence on chemical models, especially on the gas density distribution. It is therefore misleading to seek for a one–to–one relation between gas content and metallicity, or between gas profile and metallicity gradient, like the one predicted by simple models (e.g. Tinsley 1980). When studying the chemical features of galaxies, it is very dangerous to assume that such a relation must hold, since even mild flows can easily alter the overall distribution.
The models with smooth radial inflows presented in this section (with the possible exception of model O10RFe) are still unable to reproduce the gas density peak corresponding to the molecular ring around 4 kpc, since that needs to take into account the peculiar dynamical influence of the Galactic Bar on gas flows. This will be the issue of the next section.
## 6 The role of the Galactic Bar
There is by now substantial evidence that the Milky Way hosts a small Bar in its inner 3 kpc or so. The idea was originally suggested to explain the kinematics of the atomic and molecular gas near the Galactic Centre (de Vaucouleurs 1964, Peters 1975, Liszt & Burton 1980, Mulder & Liem 1986). In recent years further evidence for a Galactic Bar has piled up from several tracers: gas dynamics (Binney et al. 1991; Weiner & Sellwood 1996, 1999; Yuan 1993; Wada et al. 1994; Englmaier & Gerhard 1999; Fux 1999), IR photometry and star counts (Blitz & Spergel 1991, Weinberg 1992, Nikolaev & Weinberg 1997, Dwek et al. 1995, Binney et al. 1997, Unavane & Gilmore 1998), stellar kinematics (Zhao et al. 1994; Weinberg 1994; Fux 1997; Sevenster 1997, 1999; Sevenster et al. 1999; Raboud et al. 1998), OGLE data (Stanek et al. 1994, 1997; Paczynski et al. 1994; Evans 1994; Zhao et al. 1995, 1996; Zhao & De Zeeuw 1998, Ng et al. 1996). For a review on the Galactic Bar see e.g. Gerhard (1996, 1999). The determination of the characteristic parameters (size, axis ratio, rotation speed, orientation and so forth) is even more difficult for the Bar of our own Milky Way than for external galaxies. Broadly speaking, the various studies mentioned above indicate a Bar with a major axis of 2–4 kpc viewed at an angle of 15–45<sup>o</sup> in the first longitude quadrant, an axis ratio around 3:1 and a pattern speed $`\mathrm{\Omega }_p60`$ km sec<sup>-1</sup> kpc<sup>-1</sup>.
PC99 underlined that the dynamical influence of the Galactic Bar is likely to account for the peak at 4 kpc displayed by the gas profile in the Disc, which static chemical models are unable to reproduce if other constraints, like the observed metallicity gradient, are to be matched as well. In fact, gravitational torques in a barred, or non–axisymmetric, potential are thought to induce gas accumulation and formation of rings at the corresponding Lindblad resonances (e.g. Combes & Gerin 1985; Schwarz 1981, 1984). In brief, bar–induced flows sweep gas away from the co–rotation (CR) radius, where the bar roughly ends, toward the inner and outer Lindblad resonances (ILR, OLR). In fact, we developed our new chemical model with radial flows also with the aim to mimic such effects of the Bar upon the gas distribution, by simulating suitable flow velocity profiles.
As mentioned above, though the existence and gross features of the Galactic Bar are by now established, there is no general agreement on details like its size and pattern speed, and on the corresponding radii for its CR and ILR, OLR. In this paper, with the aim to reproduce the molecular ring at 4–6 kpc, we will consider two Bar models covering the range of scenarios suggested in literature:
Bar’s CR around 3.5 kpc, so that the dip in the gas distribution between 1.5 and 3.5 kpc is interpreted as a fast inward drift of gas from CR to the ILR and the nuclear ring;
Bar’s CR around 2.5 kpc and OLR around 4.5 kpc, so that the molecular ring is interpreted as accumulation of gas from CR toward the OLR.
In any case, we will assume here that the Bar influences only the inner 5–6 kpc of the Galactic Disc, where its OLR is supposed to lie at the outermost according to current understanding, while leaving regions outside the OLR unaffected (see also Gerhard 1999). Actually, the possible influence of the Bar over a larger Disc region than its formal extent is still an open problem, as we will comment upon in the final conclusions (§7).
In the framework of our chemical model, the inclusion of the Bar translates into imposing a suitable velocity profile for the radial flows in the inner regions of the disc. Namely, we run the “successful models” with radial flows presented in §5, but at a suitable age the Bar is assumed to form and the radial velocity profile is altered correspondingly, by modifying the coefficients $`\alpha _k`$, $`\beta _k`$, $`\gamma _k`$ describing the radial flow pattern (§3). In case A, we will impose fast radial inflow velocities within CR at 3–4 kpc to mimic the rapid drift of the gas toward the ILR (§6.1). In case B, we will impose outflows from CR to the OLR around 4.5 kpc (§6.2).
As to the age of the Bar, an upper limit is set by the typical age of its stellar population, 8–9 Gyr (Ng et al. 1996); an age of 5 to 8 Gyr has been suggested by Sevenster (1997, 1999). The results from our simulations turned out to be quite insensitive to a change in the Bar’s age from 9 to 5 Gyr; therefore, we will present simulations with a 5 Gyr old Bar ($`T_{bar}=5`$) as representative of the generic case of age $`\stackrel{>}{}5`$ Gyr. We will also consider, for the sake of completeness, the case of a much younger Bar of 1 Gyr of age ($`T_{bar}=1`$).
### 6.1 Modelling the effects of the Bar: case A
To mimic case A, where the Bar induces fast inflows from CR to its ILR, the inflow velocity is typically increased for $`r<3.54`$ kpc (the Bar’s CR radius) with respect to the drift pattern adopted before the onset of the Bar. Models corresponding to case A give good results when combined with a Schmidt SF law with radial inflows. Starting from the corresponding “successful” model S15RFe of §5 and changing the velocity law as in Fig. 9 at the onset of the Bar, we obtained the models shown in Fig. 10, compared to the original model S15RFe with no Bar’s effects included. Obviously, if the Bar is younger ($`T_{bar}=1`$ Gyr) faster induced inflows need to be assumed in the inner regions so that the observed sharp dip in the gas profile at $`r\stackrel{<}{}3.5`$ kpc is obtained in a shorter time (Fig. 9). Anyways, no extreme speeds need to be induced by the Bar ($`|v|<5`$ km sec) at these radii yet, to resemble the observed gas profile, so the models remain plausible.
This type of solution actually corresponds to the one originally suggested by Lacey & Fall (1985), who in fact assumed a Schmidt SF law, radial inflows of the order of –1 km sec<sup>-1</sup> to reproduce the metallicity gradient, and for $`r4`$ kpc a raise in the inflow speed up to –10 km sec<sup>-1</sup> to reproduce the gas profile. Without such a spike in the inflow velocity, the gas distribution keeps rising inward, with no depression (cfr. model S15RFe). This picture also corresponds to the situation for viscous models as suggested by Thon & Meusinger (1998): the detailed gas profile of the Disc can be reproduced only by artificially increasing the viscosity in the inner regions, so as to mimic the influence of the Galactic Bar.
Case A can also well combine with model O10RFe, namely with an Oort–type SF law with $`\kappa =1`$ plus the corresponding “successful” inflow pattern. In fact, also in model O10RFe (which is actually the only model able to predict a peak in the gas profile reminding of the molecular ring even without including the Bar, see §5) the gas profile rises inward, and by increasing the inflow velocities inside 4 kpc one can fit the observed gas depression. As an example, we show models O10RF with inflow patterns as in Fig. 11; the corresponding results are plotted in Fig. 12 together with the base model O10RFe.
Notice how in all these “case A” models the metallicity gradient is only negligibly affected by the switch of the inner inflow pattern at the time of onset of the Bar, $`t_GT_{bar}`$ ($`t_G=15`$ Gyr is the present Galactic age). The effects are at most limited to $`r\stackrel{<}{}3`$ kpc, where they are hard to check from observations, since in that region the gas distribution is depressed and the metallicity tracers are missing as well, since they are young objects strongly correlated to present–day SF activity and therefore to the presence of gas. We remark that abundance data at $`r=0`$ in the plots for the abundance gradient are to be disregarded as a constraint for the model, since they refer to the Galactic Centre population, not to the Disc (see PC99).
### 6.2 Modelling the effects of the Bar: case B
According to what we labelled above as “case B”, the Galactic Bar ends in correspondence to its CR around 2.5 kpc, and it has an OLR around 4.5 kpc. Therefore, the gas is expected to drift from CR outward and accumulate at the OLR, while gas inflowing from outer regions slows down and accumulates as well at the level of the resonance at 4.5 kpc.
Models for case B will adopt, starting from $`t_GT_{bar}`$, a radial flow pattern with positive velocities (outflow) from $`r=2.5`$ kpc to $`r4.5`$ kpc, and negative velocities (inflow) from outside dropping to zero at $`r4.5`$ kpc. Case B can be combined with model DRRFe, namely with the SF law by Dopita & Ryder (1994) and the corresponding “successful” overall inflow pattern (§5). The relevant models with their detailed velocity patterns, for the usual two values of $`T_{bar}`$, are shown in Figs. 13 and 14. Notice once again that rather low drift velocities ($`|v|<0.5`$ km sec) suffice to reproduce the gas peak, which reinforces the plausibility of the models.
In case B models, the metallicity gradient is clearly affected in the inner regions by the switch in the gas flow pattern, at least if this lasted for some time (Fig. 14, case $`T_{bar}=5`$). This feature, though, is again hard to check from observations since there are no tracers of Disc metallicity for $`r<4`$ kpc (see §6.1).
Case B can also reproduce the observed gas profile combined with an Oort–type SF law with $`\kappa =1`$ (model O10RFe), especially provided the Bar formed recently, as displayed in Fig. 16 with $`T_{bar}=1`$. If the Bar–induced flows activate much earlier, the predicted peak is very sharp and narrow (Fig. 16, case $`T_{bar}=5`$ Gyr), due to the milder sensitivity of SF to the gas density in this case (Schmidt-like exponent $`\kappa =1`$): with respect to other SF laws, this SF is relatively less efficient where the gas density is high, namely where the gas accumulates around 4.5 kpc, while it is relatively more efficient where the density drops, below 4 kpc. Such a sharp peak seems in contrast with the observed, quite broad distribution (the observational uncertainty on the gas density profile in the inner Galactic region is less than a factor of 2, Dame 1993). Models O10RF with an “old” Bar are therefore less appealing in case B.
### 6.3 Concluding remarks
We referred to the current understanding of the structure and features of the Galactic Bar to simulate its dynamical influence on the surrounding gas flows and distribution. Chemical models accounting for the effect of the Bar make it finally possible to reproduce the gas profile properly, which could not be accomplished by static models (PC99).
Broadly speaking, two main scenarios are presented.
1. The Bar stretches to 3.5 kpc and the gas peak is due to a rapid depletion of gas drifting from the Bar’s CR to the Galactic Centre (case A). To reproduce the gas profile, we need a gas distribution which keeps increasing from the outer disc inward, down to $`4`$ kpc where the Bar’s influence produces the sharp depletion. This can be achieved only by means of efficient radial inflows over the whole disc, and in this case radial inflows are also the major source for the observed metallicity gradients (models S15RFe and O10RFe).
2. The Bar ends around 2.5 kpc, where its CR is set, and the gas peak is due to the accumulation of gas from CR outward to its OLR at $`4.5`$ kpc (case B). A rather limited contribution of radial inflows from the outer regions suffices to obtain the peak. In this case, the metallicity gradients in the outer regions of the disc may be mainly due to the SF process itself, and to the intrinsic variation of the SF efficiency with radius (model DRRFe).
Within these simplified models it is unfeasible to discuss any further on the scenarios for Bar structure and age, and related gas flows. Only detailed dynamical simulations for Bar formation, evolution and potential can tell how the molecular ring consequently formed, which is beyond the goals of this paper (see also §7). Here we were just interested in showing how even simple qualitative models for Bar–induced gas flows in the inner disc solve in fact the puzzle encountered in PC99. Namely, they can reproduce at the same time both the metallicity gradient and the gas distribution, in particular the peak corresponding to the molecular ring around 4 kpc. This can be accomplished already with quite slow, and largely plausible, flow velocities. The present modelling therefore provides a simple tool for qualitative understanding of possible behaviours.
Regardless of details, however, one condition is necessary for the scheme to work: there must be enough gas in the inner regions of the Disc, which the Bar can then “shape” to resemble the detailed observed distribution. This favours chemical models with radial inflows where the metallicity gradient can coexist with high gas fractions in the inner shells. Our simulations showed that no Bar–induced gas flow superimposed on otherwise static models (as those by PC99) can produce the observed gas peak: whatever the assumed age of the Bar or gas velocity profile, there is not enough gas left in the inner Galactic regions if we are to reproduce the metallicity gradient as well. Gas must be continuously replenished by inflows from outer regions; that’s why we presented here “barred” models based only on the “successful” models with radial inflows from §5 and never on the static models of PC99.
## 7 Summary and conclusions
From the results of static chemical models, PC99 underlined the need to introduce radial flows to explain some features of the Galactic Disc. In fact, static models are unable to reproduce, at the same time, both the metallicity gradient and the radial gas profile; in particular, the peak corresponding to the molecular ring at 4–6 kpc is likely to be a consequence of gas drifts induced by the dynamical influence of the Galactic Bar. Therefore, in the present paper we introduced a new chemical model including radial gas flows, developed as a multi–dimensional generalization of the original static model (§3). Our model is conceived so as to adapt to any imposed radial velocity profile, describing both inflows and outflows in any part of the disc. The model is carefully tested against instability problems and spatial resolution, by comparing it to suitable exact analytical cases (Appendix B). In this paper we applied the model to the Galactic Disc; more in general, such models allowing for gas drifts are meant to be used as fast and handy interfaces between detailed dynamical galaxy models (predicting the velocity profiles) and parametric chemical and spectro–photometric models.
An overview of the behaviour of chemical models with radial inflows of gas shows that these provide an alternative “dynamical” assumption to the inside–out disc formation scenario to explain the metallicity gradient (§4). With radial gas flows, the model can reproduce the metallicity gradient even in the case of a Schmidt or an Oort–type SF law, which were excluded in the case of static models (see PC99). In addition, it appears that even low radial flow velocities, well within observational limits and theoretical expectations (see §2), have non–negligible effects upon model predictions on the metallicity gradient and moreover on the gas distribution. In particular, if radial gas inflows are allowed for, a metallicity gradient can coexist with a high gas fraction in the inner regions, at odds with simple static models. This is indispensable to reproduce the observed gas distribution in the inner Galaxy (see point 2 below). The remarkable effects of even slow radial flows upon observable quantities, mainly upon the gas distribution, should be kept in mind as a caveat when comparing real galaxies to simple analytical models which predict a one–to–one relationship between metallicity and gas fraction (e.g. Tinsley 1980). Our models show that small dynamical effects, like slow gas flows, can easily make real systems depart from the behaviour of simple models.
With our model it is possible to mimic the dynamical influence of the Galactic Bar and reproduce the peak in the gas distribution around 4 kpc (§6). Two scenarios, related to two different models for the structure of the Bar, are qualitatively suggested.
With a Schmidt or an Oort–type SF law, slow radial inflows in the disc pile up gas inward down to $`r=3.54`$ kpc. Here, the Bar CR radius is found and the gas is quickly swept inward from CR toward an ILR, which causes the drop in the gas profile at 3.5 kpc.
With a SF law like that by Dopita & Ryder (1994), smaller inflow rates suffice to reproduce the metallicity gradient, leading to a lower concentration of gas in the inner regions than in the previous case. The peak corresponding to the molecular ring can be reproduced with a Bar CR around 2.5 kpc and its OLR around 4.5 kpc, so that all the gas external to $`r=2.5`$ kpc tends to pile up around the OLR.
Though these models are just qualitative and cannot describe the detailed dynamical process of Bar formation nor the evolution of the related gas flows to form the molecular ring, they provide two interesting indications.
1. Only when introducing the effects of the Bar, the model is able to reproduce the radial gas profile properly. The only possible exception resides in a particular combination of an Oort-type SF law with a radial inflow pattern whose velocity decreases inward (model O10RFe; this combination may lead to a peak of the gas distribution in the inner Galactic regions, closed to the observed molecular ring. But this particular, fortunate case does not diminish the general conclusions about the role of the Galactic Bar.
2. In any case (A or B above), overall radial inflows in the disc are indispensable to replenish the inner regions with enough gas that the observed molecular ring can form under the influence of the Bar. This seems to favour disc models with radial inflows, unless one assumes that the gas in the ring has some different origin (gas swept from the Bulge, or accreted later).
To investigate these issues any further, detailed gas–dynamical models are obviously required. Unfortunately, most studies on Bar–induced gas dynamics (see references in §6) concentrate on the observed features of the very inner regions, such as the nuclear ring, the 3 kpc expanding arm, and so forth. Little discussion can be found about the effects of the Bar on more external regions, and on the formation of the molecular ring in particular: whether it is due to gas depletion inside CR as in our case A, or due to gas accumulation at some resonance (e.g. Binney et al. 1991, Fux 1999) as in our case B, or whether it just consists of two or more tightly wound spiral arms (e.g. Englmaier & Gerhard 1999). Further gas–dynamical studies suggesting detailed scenarios, time-scales, and velocity profiles for the formation of the molecular ring would be welcome, for the sake of including the effects of the Bar in chemical evolution models more consistently.
Further investigation of gas–dynamical models on the influence of the Bar on even larger scales (namely, outside its OLR) should be pursued as well, since this is a clue issue related to a claimed discrepancy between the characteristics of the Galactic Bar and the observed metallicity gradient. It is well known that barred galaxies display systematically shallower gradients than ordinary spirals (e.g. Alloin et al. 1981, Vila–Costas & Edmunds 1992, Martin & Roy 1994). This is likely a consequence of the radial mixing induced by bars; in fact, Martin & Roy (1994) found a correlation for external galaxies between the strength of a bar and the metallicity gradient. Taking this empirical relation at face value, the Galactic Bar with an axial ratio of $`0.5`$ should induce a metallicity gradient of –0.03 dex/kpc, much shallower than the observed one of –0.07 dex/kpc, which is typical of a normal Sbc galaxy. To overcome such a puzzle, it has been suggested that the Galactic Bar must be very young ($`<`$1 Gyr), so that there was not enough time yet to flatten the gradient (Gummersbach et al. 1998); but this is in conflict with other estimates of the Bar’s age (e.g. Sevenster 1997, 1999). Alternatively, we suggest that the discrepancy might be only apparent, since the Galactic Bar is quite small, and the Milky Way cannot be properly considered a barred spiral. It might be unlikely that the Bar can influence the metallicity gradient all over the Disc, as in really barred galaxies: Bar–induced radial drifts and corresponding chemical mixing are expected to occur from CR toward the ILR (inflows) and to the OLR (outflows; e.g. Schwarz 1981, 1984; Friedli et al. 1994). Present understanding of the Galactic Bar sets its OLR between 4.5 and 6 kpc (§6 and references therein), so in our models we presumed that the Bar induces negligible mixing beyond these radii, regardless of its age (see also Gerhard 1999). If the situation is as in the models we presented here, in fact, the metallicity gradient in the outer regions is unperturbed and just related to intrinsic Disc properties and/or large–scale viscous flows.
However, gas–dynamical simulations dedicated to the effects of the Galactic Bar over the whole Disc would be necessary, so as to investigate the relation between the Bar, radial mixing and the metallicity gradient more consistently. More in general, including the effects of bar–induced radial flows in the picture of the chemical evolution of spiral galaxies might turn out to be of wide interest, since it is likely that all spirals develop at some point, or have developed in the past, some bar–like structure (Binney 1995). Infrared observations indeed reveal that a large fraction of spirals host a barred structure (e.g. Eskridge et al. 1999), and recent numerical simulations suggest that even weak bars or oval distortions may be able to induce radial drifts to form multiple gaseous rings at the corresponding Lindblad resonances (Jungwiert & Palouš 1996). Bars could even drive secular evolution of spiral discs from late to early type (e.g. Dutil & Roy 1999). Bar–driven radial gas flows might therefore play a fundamental role in the chemical evolution of spiral discs.
###### Acknowledgements.
We thank Joachim Köppen for his advice on numerical modelling of radial gas flows, Yuen K. Ng and Antonella Vallenari for useful discussions about Galactic structure, and our referee, Mike Edmunds, whose suggestions much improved the presentation of our paper. L.P. acknowledges kind hospitality from the Nordita Institute in Copenhagen, from the Observatory of Helsinki and from Sissa/Isas in Trieste. This study has been financed by the Italian MURST through a PhD grant and the contract “Formazione ed evoluzione delle galassie”, n. 9802192401.
## Appendix A The explicit expression of Eq. (18)
Here we sort out the explicit expression (20) of the solution (18) for $`G_i(r_k,t)`$, by calculating the matrix $`e^{t𝒜}`$. Let us first introduce the simplified notation:
$$G_k(t)=G_i(r_k,t),\overline{W_k}=\overline{W}_i(r_k)$$
(21)
For the sake of example, we present here the detailed solution in the case of three shells ($`N=3`$). The relevant system of equations (17) becomes:
$$\{\begin{array}{c}\frac{d}{dt}G_1(t)=\vartheta _1G_1(t)+\gamma _1G_2(t)+W_1\hfill \\ \\ \frac{d}{dt}G_2(t)=\alpha _2G_1(t)+\vartheta _2G_2(t)+\gamma _2G_3(t)+W_2\hfill \\ \\ \frac{d}{dt}G_3(t)=\alpha _3G_2(t)+\vartheta _3G_3(t)+W_3+\overline{\omega _i}\hfill \end{array}$$
(22)
where all the coefficients are considered as constants (see §3.2), though we have omitted the bar over $`W_k`$ and $`\vartheta _k`$ for simplicity. A system like (22) holds for each chemical species $`i`$, but the characteristic matrix $`𝒜`$ is the same for any $`i`$ (see §3.2); in this case:
$$𝒜=\left(\begin{array}{ccc}\vartheta _1& \gamma _1& 0\\ \alpha _2& \vartheta _2& \gamma _2\\ 0& \alpha _3& \vartheta _3\end{array}\right)$$
Let’s first notice, from the definition (9), that $`\gamma _k`$ and $`\alpha _{k+1}`$ can never be both positive: at least one of them must be zero since they are “activated” in the opposite cases of inflow or outflow at $`r_{k+\frac{1}{2}}`$, respectively; if there is no flow at all through $`r_{k+\frac{1}{2}}`$, they both reduce to zero. Therefore, in our calculations we are always entitled to use the condition:
$$\gamma _k\alpha _{k+1}=0k$$
(23)
The eigenvalues of the matrix $`𝒜`$ are
$$\lambda _{1,2,3}=\vartheta _{1,2,3}$$
and the associated eigenvectors are:
$$𝐮_1=\left(\begin{array}{c}1\\ \\ \frac{\alpha _2}{\vartheta _1\vartheta _2}\\ \\ \frac{\alpha _3}{\vartheta _1\vartheta _3}\frac{\alpha _2}{\vartheta _1\vartheta _2}\end{array}\right)𝐮_2=\left(\begin{array}{c}\frac{\gamma _1}{\vartheta _2\vartheta _1}\\ \\ 1\\ \\ \frac{\alpha _3}{\vartheta _2\vartheta _3}\end{array}\right)𝐮_3=\left(\begin{array}{c}\frac{\gamma _1}{\vartheta _3\vartheta _1}\frac{\gamma _2}{\vartheta _3\vartheta _2}\\ \\ \frac{\gamma _2}{\vartheta _3\vartheta _2}\\ \\ 1\end{array}\right)$$
From (19) we get:
$$e^{t𝒜}=\left(\begin{array}{ccc}e^{\vartheta _1t}& \gamma _1f_{21}\left(t\right)& \frac{\gamma _1\gamma _2}{\vartheta _3\vartheta _2}\left[f_{31}\left(t\right)f_{21}\left(t\right)\right]\\ & \\ \alpha _2f_{21}\left(t\right)& e^{\vartheta _2t}& \gamma _2f_{32}\left(t\right)\\ & \\ \frac{\alpha _2\alpha _3}{\vartheta _2\vartheta _1}\left[f_{32}\left(t\right)f_{31}\left(t\right)\right]& \alpha _3f_{32}\left(t\right)& e^{\vartheta _3t}\end{array}\right)$$
where we have indicated with:
$$f_{kl}(t)=f_{lk}(t)\frac{e^{\vartheta _kt}e^{\vartheta _lt}}{\vartheta _k\vartheta _l},k,l=1,2,3$$
Defining:
$$g_k(\mathrm{\Delta }t)_{t_0}^{t_1}e^{\vartheta _k(t_1t)}𝑑t=\frac{e^{\vartheta _k\mathrm{\Delta }t}1}{\vartheta _k}$$
the solution (18) in the case $`N=3`$ becomes:
$$\{\begin{array}{cc}G_1(t_1)=\hfill & e^{\vartheta _1\mathrm{\Delta }t}G_1(t_0)+\hfill \\ & +\gamma _1\frac{e^{\vartheta _2\mathrm{\Delta }t}e^{\vartheta _1\mathrm{\Delta }t}}{\vartheta _2\vartheta _1}G_2(t_0)+\hfill \\ & +\gamma _1\gamma _2\frac{f_{31}(\mathrm{\Delta }t)f_{21}(\mathrm{\Delta }t)}{\vartheta _3\vartheta _2}G_3(t_0)+\hfill \\ & +W_1\frac{e^{\vartheta _1\mathrm{\Delta }t}1}{\vartheta _1}+\hfill \\ & +\gamma _1\frac{W_2}{\vartheta _2\vartheta _1}\left(\frac{e^{\vartheta _2\mathrm{\Delta }t}1}{\vartheta _2}\frac{e^{\vartheta _1\mathrm{\Delta }t}1}{\vartheta _1}\right)+\hfill \\ & +\gamma _1\gamma _2\frac{W_3}{\vartheta _3\vartheta _2}\left[\frac{g_3(\mathrm{\Delta }t)g_1(\mathrm{\Delta }t)}{\vartheta _3\vartheta _1}\frac{g_2(\mathrm{\Delta }t)g_1(\mathrm{\Delta }t)}{\vartheta _2\vartheta _1}\right]\hfill \\ & \\ G_2(t_1)=\hfill & \alpha _2\frac{e^{\vartheta _2\mathrm{\Delta }t}e^{\vartheta _1\mathrm{\Delta }t}}{\vartheta _2\vartheta _1}G_1(t_0)+\hfill \\ & +e^{\vartheta _2\mathrm{\Delta }t}G_2(t_0)+\hfill \\ & +\gamma _2\frac{e^{\vartheta _3\mathrm{\Delta }t}e^{\vartheta _2\mathrm{\Delta }t}}{\vartheta _3\vartheta _2}G_3(t_0)+\hfill \\ & +\alpha _2\frac{W_1}{\vartheta _2\vartheta _1}\left(\frac{e^{\vartheta _2\mathrm{\Delta }t}1}{\vartheta _2}\frac{e^{\vartheta _1\mathrm{\Delta }t}1}{\vartheta _1}\right)+\hfill \\ & +W_2\frac{e^{\vartheta _2\mathrm{\Delta }t}1}{\vartheta _2}+\hfill \\ & +\gamma _2\frac{W_3}{\vartheta _3\vartheta _2}\left(\frac{e^{\vartheta _3\mathrm{\Delta }t}1}{\vartheta _3}\frac{e^{\vartheta _2\mathrm{\Delta }t}1}{\vartheta _2}\right)\hfill \\ & \\ G_3(t_1)=\hfill & \alpha _2\alpha _3\frac{f_{32}(\mathrm{\Delta }t)f_{31}(\mathrm{\Delta }t)}{(\vartheta _2\vartheta _1)}G_1(t_0)+\hfill \\ & +\alpha _3\frac{e^{\vartheta _3\mathrm{\Delta }t}e^{\vartheta _2\mathrm{\Delta }t}}{\vartheta _3\vartheta _2}G_2(t_0)+\hfill \\ & +e^{\vartheta _3\mathrm{\Delta }t}G_3(t_0)+\hfill \\ & +\alpha _2\alpha _3\frac{W_1}{\vartheta _2\vartheta _1}\left[\frac{g_3(\mathrm{\Delta }t)g_2(\mathrm{\Delta }t)}{\vartheta _3\vartheta _2}\frac{g_3(\mathrm{\Delta }t)g_1(\mathrm{\Delta }t)}{\vartheta _3\vartheta _1}\right]+\hfill \\ & +\alpha _3\frac{W_2}{\vartheta _3\vartheta _2}\left(\frac{e^{\vartheta _3\mathrm{\Delta }t}1}{\vartheta _3}\frac{e^{\vartheta _2\mathrm{\Delta }t}1}{\vartheta _2}\right)+\hfill \\ & +(W_3+\overline{\omega _i})\frac{e^{\vartheta _3\mathrm{\Delta }t}1}{\vartheta _3}\hfill \end{array}$$
(24)
With zero flow velocity ($`\alpha _k=\beta _k=\gamma _k=0`$), (24) reduces to the solving formula of the original static model (see PCB98). Notice that the solution $`G_1(t_1)`$ for the $`1^{st}`$ shell includes not only the contribution of the contiguous $`2^{nd}`$ shell, but also a contribution from the $`3^{rd}`$ shell “scaled” by its passage through the $`2^{nd}`$ shell. Similarly, the $`3^{rd}`$ shell is affected not only by the $`2^{nd}`$, but also by the $`1^{st}`$ shell though they are not contiguous.
With analogous procedure, for an arbitrary number $`N`$ of shells the solution is of the kind:
$$G_k(t)=\underset{l=1}{\overset{N}{}}F_{kl}(\mathrm{\Delta }t)G_l(t_0)+\underset{m=0}{\overset{N}{}}H_{km}(\mathrm{\Delta }t)W_m$$
(25)
where:
$$\begin{array}{ccc}F_{kk}(\mathrm{\Delta }t)=\hfill & e^{\vartheta _k\mathrm{\Delta }t}\hfill & \\ F_{kl}(\mathrm{\Delta }t)=\hfill & \gamma _k\gamma _{k+1}\mathrm{}\mathrm{}.\gamma _{l1}_{k(k+1)\mathrm{}.l}(\mathrm{\Delta }t)\hfill & \hfill l>k\\ F_{kl}(\mathrm{\Delta }t)=\hfill & \alpha _{l+1}\mathrm{}\mathrm{}.\alpha _{k1}\alpha _k_{l(l+1)\mathrm{}.k}(\mathrm{\Delta }t)\hfill & \hfill l<k\\ H_{kk}(\mathrm{\Delta }t)=\hfill & g_k(\mathrm{\Delta }t)\hfill & \\ H_{kl}(\mathrm{\Delta }t)=\hfill & \gamma _k\gamma _{k+1}\mathrm{}\mathrm{}.\gamma _{l1}_{k(k+1)\mathrm{}.l}(\mathrm{\Delta }t)\hfill & \hfill l>k\\ H_{kl}(\mathrm{\Delta }t)=\hfill & \alpha _{l+1}\mathrm{}\mathrm{}.\alpha _{k1}\alpha _k_{k(k1)\mathrm{}.l}(\mathrm{\Delta }t)\hfill & \hfill l<k\end{array}$$
and the quantities $``$ and $``$ are constructed by means of recursive formulæ:
$$\begin{array}{cc}_{ki}(\mathrm{\Delta }t)=f_{ki}(\mathrm{\Delta }t)\hfill & _{ki}(\mathrm{\Delta }t)=\frac{g_k\mathrm{\Delta }tg_i\mathrm{\Delta }t}{\vartheta _k\vartheta _i}\hfill \\ & \\ _{kij}=\frac{_{ki}_{kj}}{\vartheta _i\vartheta _j}\hfill & _{kij}=\frac{_{ki}_{kj}}{\vartheta _i\vartheta _j}\hfill \\ & \\ _{kijm}=\frac{_{kij}_{kim}}{\vartheta _j\vartheta _m}\hfill & _{kijm}=\frac{_{kij}_{kim}}{\vartheta _j\vartheta _m}\hfill \\ & \\ _{kijmn}=\frac{_{kijm}_{kijn}}{\vartheta _m\vartheta _n}\hfill & _{kijmn}=\frac{_{kijm}_{kijn}}{\vartheta _m\vartheta _n}\hfill \\ & \\ \mathrm{}\hfill & \mathrm{}\hfill \end{array}$$
The coefficients $`F_{kl}`$ and $`H_{kl}`$ describe the contribution of the generic shell $`l`$ to the chemical evolution of $`k`$. Notice that a shell external to $`k`$ ($`l>k`$) can influence $`k`$ only if all the inflow coefficients $`\gamma `$ in between $`l`$ and $`k`$ are non-zero, namely if there is a continuous inflow from $`l`$ to $`k`$, as expected. The same holds for inner shells $`l<k`$, whose contribution $`F_{kl}`$, $`H_{kl}`$ is non-zero only if none of the intermediate outflow coefficients $`\alpha `$ is zero. The solution (25) gets more and more complicated the larger the number $`N`$ of shells considered, since each shell formally feels the contribution of all the other shells (as already noticed in the above case $`N=3`$). This occurs because (25) would be the exact analytical solution in the case of a differential system with constant coefficients, namely if the $`𝒜`$ matrix in (17) were constant. Then, (25) would describe the complete evolution of any shell $`k`$, which over a galaxy’s lifetime would indeed process and exchange gas drifting from or to rather distant shells. But $`𝒜(t)`$ is not constant even when the flow pattern $`\alpha `$, $`\beta `$, $`\gamma `$ is constant, because $`\eta _k(t)`$ and therefore $`\vartheta _k(t)`$ evolve in time due to SF; in fact, we apply (25) only upon short timesteps $`\mathrm{\Delta }t`$, within which $`𝒜(t)`$ can be considered approximately constant. If $`\mathrm{\Delta }t`$ is short enough with respect to the radial flow velocities — as guaranteed by the Courant condition $`\mathrm{\Delta }t<v\mathrm{\Delta }r`$, see §3.2 —, we can assume that within $`\mathrm{\Delta }t`$ the $`k`$-th shell is affected just by the flows from the contiguous shells $`k`$+1 and $`k`$–1, and not from more distant shells, although all of them formally contribute to the solution. In this approximation we neglect all higher order terms in $`\alpha `$ and $`\gamma `$, namely all the terms $`𝒪(\alpha _i\alpha _j)`$ and $`𝒪(\gamma _i\gamma _j)`$, keeping only the “linear” terms of the kind $`\alpha _kf_{k(k1)}`$ and $`\gamma _kf_{k(k+1)}`$; so the general solution (25) reduces in fact to (20).
## Appendix B Testing the numerical model
Since discretized numerical solutions for partial differential equations containing an “advection term” tend to be affected by instability problems (e.g. Press et al. 1986), we tested our numerical code with gas flows against suitable analytical cases. We report here two representative tests involving pure gas flows (no SF) which allow for exact analytical solutions; to these we compare the predictions of our numerical code with a SF efficiency dropped virtually to zero.
Selecting the timestep. The first reference analytical case is that of a purely gaseous disc of infinite radial extent and flat profile, where the gas is:
1. accreting uniformly with a time-scale $`\tau `$;
2. flowing inward with a constant (in time) and uniform (in space) velocity $`v`$, starting from the “inflow onset time” $`T_{rf}`$.
Such a system is governed by the same differential equation (13) describing the adopted boundary condition at the outer disc edge in the chemical model (§3.1.2). Eq. (13) is a linear, first order, partial differential equation, equivalent to the system:
$$\{\begin{array}{cc}\frac{dr}{dt}=v\hfill & \hfill (a)\\ & \\ \frac{d\sigma }{dt}=Ae^{\frac{t}{\tau }}\frac{v}{r}\sigma \hfill & \hfill (b)\end{array}$$
(26)
Eq. (26a) is solved as:
$$r=v(tt_0)+r_0$$
and substitution into Eq. (26b) yields:
$$\frac{d\sigma }{dt}+\frac{1}{t+p}\sigma =Ae^{\frac{t}{\tau }}p\frac{r_0}{v}t_0$$
(27)
This linear, first order, ordinary differential equation is solved as:
$$\begin{array}{cc}\sigma (t)\hfill & =\sigma (t_0)e^{_{t_0}^t\frac{1}{\xi +p}𝑑\xi }+_{t_0}^te^{_\xi ^t\frac{1}{\theta +p}𝑑\theta }Ae^{\frac{\xi }{\tau }}𝑑\xi \hfill \\ & \\ & =\sigma (t_0)\frac{t_0+p}{t+p}+\frac{A}{t+p}\tau \times \hfill \\ & \times \left[(t_0+p)e^{\frac{t_0}{\tau }}(t+p)e^{\frac{t}{\tau }}+\tau \left(e^{\frac{t_0}{\tau }}e^{\frac{t}{\tau }}\right)\right]\hfill \end{array}$$
Finally, replacing back:
$$r_0=rv(tt_0),t+p=\frac{r}{v}$$
we get:
$$\begin{array}{cc}\sigma (r,t)=\hfill & [1\frac{v}{r}(tt_0)]\sigma (rv(tt_0),t_0)+A\tau \times \hfill \\ & \\ \multicolumn{2}{c}{\times \left[\left(1\frac{v}{r}(tt_0)\right)e^{\frac{t_0}{\tau }}e^{\frac{t}{\tau }}+\frac{v}{r}\tau \left(e^{\frac{t_0}{\tau }}e^{\frac{t}{\tau }}\right)\right]}\end{array}$$
(28)
Let’s now set the initial conditions at $`t_0`$. If radial flows “activate” at a time $`t_0=T_{rf}0`$ the surface density distribution for $`tT_{rf}`$ is determined just by the accretion profile:
$$\sigma (r,T_{rf})=A\tau \left(1e^{\frac{T_{rf}}{\tau }}\right)$$
(see Eq. 2) and Eq. (28) becomes in fact Eq. (14). Here in our test case, Eq. (14) is the exact analytical description of the surface density profile over the whole disc (a part from the centre $`r=0`$, which is a singular point).
As a representative test, let’s consider the case $`\tau =3`$ Gyr, $`T_{rf}=0`$ and $`v=1`$ km sec<sup>-1</sup>. The relevant analytical solution is plot in Fig. 17 for $`t=t_G=15`$ Gyr (thick solid line). The numerical models used for comparison cover the radial range 2 to 20 kpc and adopt a flat accretion profile $`A(r)A`$. Their outer edge does match exactly with the analytical counterpart, since the boundary condition at $`r=20`$ kpc is given by the analytical expression (14) itself. But the predictions of numerical models at inner radii tend to deviate from the reference density profile, and the mismatch is larger:
1. the larger the typical timestep of the model;
2. for a fixed timestep, the thinner the shells (compare upper to lower panel).
For a typical shell width of 1 kpc, for instance, a good match is obtained with model timesteps of $`2\times 10^2`$ Gyr, while model shells of 0.5 kpc an acceptable profile is obtained only with timesteps of $`2\times 10^3`$ Gyr. Therefore, a reliable representation of radial gas flows is obtained only with a suitably small timestep; how small, is related to the width of the shells, namely to the resolution of the grid spacing.
Selecting the grid spacing. Since our models are to simulate a disc with an exponential profile, as a second test we consider a gaseous disc with uniform and constant infall time-scale and inflow velocity, analogous to the previous case, but with an exponential profile. Namely, in the representative differential equation (12) the accretion profile $`A(r)`$ declines exponentially outward:
$$A(r)=A(r_{})e^{\frac{rr_{}}{r_d}}$$
Eq. (12) can then be written:
$$\frac{\sigma }{t}+v\frac{\sigma }{r}=A(r_{})e^{\frac{r_{}}{r_d}}e^{\frac{r}{r_d}}e^{\frac{t}{\tau }}\frac{v}{r}\sigma $$
This is another linear, first order, partial differential equation of the same kind as (13), and can be solved with analogous procedure into:
$$\begin{array}{cc}\sigma (r,t)=\hfill & \left[1\frac{v}{r}(tt_0)\right]\sigma (rv(tt_0),t_0)+\hfill \\ & +\frac{A(r_{})e^{\frac{rr_{}}{r_d}}}{\frac{1}{\tau }+\frac{v}{r_d}}[(1\frac{v}{r}(tt_0))e^{\frac{t_0}{\tau }+\frac{v}{r_d}(tt_0)}+\hfill \\ & e^{\frac{t}{\tau }}+\frac{v}{r}\frac{e^{\frac{t_0}{\tau }+\frac{v}{r_d}(tt_0)}e^{\frac{t}{\tau }}}{\frac{1}{\tau }+\frac{v}{r_d}}]\hfill \end{array}$$
(29)
Notice that Eq. (28) for a flat profile is recovered from (29) for $`r_d\mathrm{}`$. If radial inflows set in at time $`t_0=T_{rf}0`$, from
$$\sigma (r,T_{rf})=A(r_{})e^{\frac{rr_{}}{r_d}}\tau (1e^{\frac{T_{rf}}{\tau }})$$
we get:
$$\begin{array}{cc}\sigma (r,t)\hfill & =A(r_{})e^{\frac{rr_{}}{r_d}}\{\tau (1e^{\frac{T_{rf}}{\tau }})e^{\frac{v}{r_d}(tT_{rf})}+\hfill \\ & +\frac{e^{\frac{v}{r_d}t\left(\frac{1}{\tau }+\frac{v}{r_d}\right)T_{rf}}e^{\frac{t}{\tau }}}{\frac{1}{\tau }+\frac{v}{r_d}}+\hfill \\ & \\ & +\frac{v}{r}[\frac{e^{\frac{v}{r_d}t\left(\frac{1}{\tau }+\frac{v}{r_d}\right)T_{rf}}e^{\frac{t}{\tau }}}{\left(\frac{1}{\tau }+\frac{v}{r_d}\right)^2}(tT_{rf})\times \hfill \\ \multicolumn{2}{c}{\times (\tau (1e^{\frac{T_{rf}}{\tau }})e^{\frac{v}{r_d}(tT_{rf})}+\frac{e^{\frac{v}{r_d}t\left(\frac{1}{\tau }+\frac{v}{r_d}\right)T_{rf}}}{\frac{1}{\tau }+\frac{v}{r_d}})]\}}\end{array}$$
(30)
We compared this second analytical case to our numerical model, where an exponential accretion profile was adopted between 2 and 20 kpc, and at the outer edge the boundary condition (14) was replaced by (30). Our tests showed that in this case the analytical profile is better reproduced the larger the number of shells, i.e. the finer the grid spacing. A good match is obtained especially when model shells are equally spaced in the logarithmic, rather than linear, scale; namely when the shells are chosen, in this case of an exponential accretion profile, so as to contain roughly the same mass, rather than cover the same radial width. Fig. 18 shows in fact the analytical solution (30) (for $`r_d=4`$ kpc, $`\tau =3`$, $`T_{rf}=0`$ and $`v=1`$) together with a corresponding numerical model with 40 shells logarithmically spaced (from 0.1 kpc wide for the inner ones to $``$1 kpc wide for the outer ones).
With such a grid spacing, our tests with the reference flat profile of Fig. 17 indicate $`10^4`$ Gyr as a suitable timestep to obtain stable solutions for velocities up to the order of 1 km sec<sup>-1</sup>.
In the light of all these tests, in our chemical models we adopted a grid spacing of 35 shells from 2.5 to 20 kpc, equally spaced in logarithmic scale and a typical timestep of $`10^4`$ Gyr (see §3.2). Of course, the suitable timestep depends also on the velocity field: flows with higher velocities require shorter integration timesteps. Whenever we need to consider much larger speeds than 1 km sec<sup>-1</sup>, as might be the case for the strong flows induced by the Bar, we reduce the timestep in proportion. |
warning/0002/hep-ph0002210.html | ar5iv | text | # CP VIOLATION IN HYPERON DECAYS
## 1 Introduction
The CPT theorem was proved in 1955 and soon thereafter L$`\ddot{\mathrm{u}}`$ders and Zumino deduced from it the equality of masses and lifetimes between particles and anti-particles. In 1958 Okubo observed that CP violation allows hyperons and antihyperons to have different branching ratios into conjugate channels even though their total rates must be equal by CPT. Somewhat later, this paper inspired Sakharov to his famous work on cosmological baryon-antibaryon asymmetry. In fact, he called this the “Okubo effect”, perhaps a better phrase than the current dull use of “direct CP violation.” Pais extended Okubo’s proposal to asymmetry parameters in $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ decays. The subject was revived in the ’80s and a number of calculations were made. Only now, over 40 years after Okubo’s paper, are these proposals about to be tested in the laboratory.
The reason for the current interest is the need to find CP violation in places other than just $`K_LK_S`$ complex. Only a number of different observations of CP violation in different channels will help us pin down the source and nature of CP violation in or beyond the standard model (SM). From this point of view, hyperon decay is one more weapon in our arsenal in addition to the K system, the B system, the D system, etc.
## 2 Phenomenology of Hyperon Decays
I summarize here the salient features of the phenomenology of non-leptonic hyperon decays . Leaving out $`\mathrm{\Omega }^{}`$ decays, there are seven decay modes $`\mathrm{\Lambda }N\pi ,\mathrm{\Sigma }^\pm N\pi `$ and $`\mathrm{\Xi }\mathrm{\Lambda }\pi `$. The effective matrix element can be written as
$$i\overline{u}_{\overline{p}}(a+b\gamma _5)u_\mathrm{\Lambda }\varphi $$
(1)
for the mode $`\mathrm{\Lambda }p+\pi ^{}`$, where a and b are complex in general. The corresponding element for $`\overline{\mathrm{\Lambda }}\overline{p}+\pi ^+`$ is then:
$$i\overline{v}_{\overline{p}}(a^{}+b^{}\gamma _5)v_{\overline{\mathrm{\Lambda }}}\varphi ^+$$
(2)
It is convenient to express the observables in terms of S and P and write the matrix element as
$$S+P\sigma .\widehat{𝐪}$$
(3)
where q is the proton momentum in the $`\mathrm{\Lambda }`$ rest frame and S and P are:
$`S`$ $`=`$ $`a\sqrt{{\displaystyle \frac{\left\{(m_\mathrm{\Lambda }+m_p)^2m_\pi ^2\right\}}{16\pi m_\mathrm{\Lambda }^2}}}`$
$`P`$ $`=`$ $`b\sqrt{{\displaystyle \frac{\left\{(m_\mathrm{\Lambda }m_p)^2m_\pi ^2\right\}}{16\pi m_\mathrm{\Lambda }^2}}}`$ (4)
In the $`\mathrm{\Lambda }`$ rest-frame, the decay distribution is given by:
$`{\displaystyle \frac{d\mathrm{\Gamma }}{d\mathrm{\Omega }}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }}{8\pi }}\{[1+\alpha <\sigma _\mathrm{\Lambda }>.\widehat{\sigma }]`$ (5)
$`+`$ $`<\sigma _p>.[(\alpha +<\sigma _\mathrm{\Lambda }>.\widehat{𝐪})\widehat{𝐪}+\beta <\sigma _\mathrm{\Lambda }>\times \widehat{𝐪}`$
$`+`$ $`\gamma (\widehat{𝐪}\times (<\sigma _\mathrm{\Lambda }>\times \widehat{𝐪})]\}`$
where $`\mathrm{\Gamma }`$ is the decay rate and is given by:
$$\mathrm{\Gamma }=2𝐪\left\{S^2+P^2\right\}$$
(6)
$`\alpha ,\beta `$ and $`\gamma `$ are given by
$`\alpha `$ $`=`$ $`{\displaystyle \frac{2Re(S^{}P)}{\{S^2+P^2\},}}`$
$`\beta `$ $`=`$ $`{\displaystyle \frac{2Im(SP^{})}{\left\{S^2+P^2\right\}}}`$
$`\gamma `$ $`=`$ $`{\displaystyle \frac{\left\{S^2P^2\right\}}{\left\{S^2+P^2\right\}}}`$ (7)
For a polarized $`\mathrm{\Lambda }`$, the up-down asymmetry of the final proton is given by $`\alpha (\alpha `$ is also the longitudinal polarization of the proton for an unpolarized $`\mathrm{\Lambda })`$. $`\beta `$ and $`\gamma `$ are components of the transverse polarization of proton .
The observed properties of the hyperon decays can be summarised as: (i) the $`\mathrm{\Delta }I=1/2`$ dominance i.e. the $`\mathrm{\Delta }I=3/2`$ amplitudes are about 5% of the $`\mathrm{\Delta }I=1/2`$ amplitudes; (ii) the asymmetry parameter $`\alpha `$ is large for $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ decays, $`\mathrm{\Xi }`$ decays and $`\mathrm{\Sigma }^+p\pi ^0`$ and is near zero for $`\mathrm{\Sigma }^\pm n\pi ^\pm `$; and (iii) the Sugawara-Lee triangle sum rule $`\sqrt{3}A(\mathrm{\Sigma }^+p\pi ^0)A(\mathrm{\Lambda }p\pi ^{})=2A(\mathrm{\Xi }\mathrm{\Lambda }\pi ^{})`$ is satisfied to a level of 5% in both $`s`$ and $`p`$ wave amplitudes.
## 3 CP Violating Observables
Let a particle P decay into several final states $`f_1,f_2`$ etc. The amplitude for P $`f_1`$ is in general:
$$A=A_1e^{i\delta 1}+A_2e^{i\delta 2}$$
(8)
where 1 and 2 are strong interaction eigenstates and $`\delta _i`$ are corresponding final state phases. Then the amplitude for $`\overline{P}\overline{f}_1`$ is
$$\overline{A}=A_1^{}e^{i\delta 1}+A_2^{}e^{i\delta 2}$$
(9)
If $`A_1>>A_2`$, then the rate asymmetry $`\mathrm{\Delta }(=(\mathrm{\Gamma }\overline{\mathrm{\Gamma }})/(\mathrm{\Gamma }+\overline{\mathrm{\Gamma }}))`$ is given by:
$$\mathrm{\Delta }2A_2/A_1sin(\varphi _1\varphi _2)sin(\delta _1\delta _2)$$
(10)
where $`A_i=A_ie^{i\varphi _i}`$. Hence, to get a non-zero rate asymmetry, one must have (i) at least two channels in the final state, (ii) CPV weak phases must be different in the two channels, and (iii) unequal final state scattering phase shifts in the two channels. A similar calculation of the asymmetry of $`\alpha `$ shows that for a single isospin channel dominance,
$$A=(\alpha +\overline{\alpha })/(\alpha \overline{\alpha })=2tan(\delta _s\delta _p)tan(\varphi _s\varphi _p)$$
(11)
In this case the two channels are orbital angular momenta $`0`$ and $`1`$; and even a single isospin mode such as $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$ can exhibit a non-zero A. In B decays an example of a single isospin mode exhibiting CP violating rate asymmetry is $`B\pi \pi `$, i.e. In this case the two eigen-channels with different weak CP phases and different final state phases are $`BD\overline{D}\pi \pi \pi `$ and $`B\pi \pi \pi \pi `$.
To define the complete set of CP violating observables, consider the example of the decay modes $`\mathrm{\Lambda }p\pi ^{}`$ and $`\overline{\mathrm{\Lambda }}\overline{p}\pi ^+`$. The amplitudes are:
$`S`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}S_1e^{i(\delta _1+\varphi _1^s)}+{\displaystyle \frac{1}{\sqrt{3}}}S_3e^{i(\delta _3+\varphi _3^s)}`$
$`P`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}P_1e^{i(\delta _{11}+\varphi _1^p)}+{\displaystyle \frac{1}{\sqrt{3}}}P_3e^{i(\delta _3+\varphi _3^p)}`$ (12)
where $`S_i,P_i`$ are real, $`i`$ refers to the final state isospin (i=2I) and $`\varphi _i`$ are the CPV phases. With the knowledge that $`S_3/S_1`$, $`P_3/P_1<<`$ 1 ; one can write
$`\mathrm{\Delta }_\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{(\mathrm{\Gamma }\overline{\mathrm{\Gamma }})}{(\mathrm{\Gamma }+\overline{\mathrm{\Gamma }})}}\sqrt{2}(S_3/S_1)sin(\delta _3\delta _1)sin(\varphi _3^s\varphi _1^s)`$
$`A_\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{(\alpha +\overline{\alpha })}{(\alpha \overline{\alpha })}}tan(\delta _{11}\delta _1)tan(\varphi _1^p\varphi _1^s)`$
$`B_\mathrm{\Lambda }`$ $`=`$ $`{\displaystyle \frac{(\beta +\overline{\beta })}{(\beta \overline{\beta })}}cot(\delta _{11}\delta _1)tan(\varphi _1^p\varphi _1^s)`$ (13)
The last one $`B_\mathrm{\Lambda }`$ has the peculiar feature that it blows up as the phase shift difference vanishes. The reason is that in the limit of CP conservation $`\beta +\overline{\beta }=0`$ but in the limit of no final state phase difference $`\beta \overline{\beta }=0`$. For $`\pi `$N final states, the phase shifts at $`E_{c.m.}=m_\mathrm{\Lambda }`$ are known and are: $`\delta _1=6^0,\delta _3=3.8^0,\delta _{11}=1.1^0`$ and $`\delta _{31}=0.7^0`$ from the 1965 analysis of Roper et al. with errors estimated at 10%. The CPV phases $`\varphi _i`$ have to be provided by theory.
Similar expressions can be written for other hyperon decays. For example, for $`\mathrm{\Lambda }n\pi ^0`$, $`\mathrm{\Delta }`$ is $`2\mathrm{\Delta }_\mathrm{\Lambda }`$ and $`A`$ and $`B`$ are identical to $`A_\mathrm{\Lambda }`$ and $`B_\mathrm{\Lambda }`$. For $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$ ( and $`\mathrm{\Xi }^0\mathrm{\Lambda }\pi ^0)`$ the asymmetries are:
$`\mathrm{\Delta }_\mathrm{\Xi }`$ $`=`$ $`0`$
$`A_\mathrm{\Xi }`$ $`=`$ $`tan(\delta _{21}\delta _2)tan(\varphi ^p\varphi ^s)`$
$`B_\mathrm{\Xi }`$ $`=`$ $`cot(\delta _{21}\delta _2)tan(\varphi ^p\varphi ^s)`$ (14)
where $`\delta _{21}`$ and $`\delta _2`$ are the $`p`$ and $`s`$-wave $`\mathrm{\Lambda }\pi `$ phase shifts at $`m_\mathrm{\Xi }`$ respectively. Somewhat more complicated expressions can be and have been written for $`\mathrm{\Sigma }`$ decays.
## 4 Calculating CP Phases
In standard model description of the non-leptonic hyperon decays, the effective $`\mathrm{\Delta }S=1`$ Hamiltonian is
$$H_{eff}=\frac{G_F}{\sqrt{2}}U_{ud}^{}U_{us}\underset{i=1}{\overset{12}{}}c_i(\mu )O_i(\mu )$$
(15)
after the short distance QCD corrections (LLO + NLLO) where $`c_i=z_i+y_i\tau (\tau =U_{td}U_{ts}^{}/U_{ud}U_{us})`$, and $`\mu 0(1`$ GeV). For CP violation, the most important operator is:
$$O_6=\overline{d}\lambda _i\gamma _\mu (1+\gamma _5)s\overline{q}\lambda _i\gamma _\mu (1\gamma _5)q$$
(16)
and $`y_60.1`$ at $`\mu 1GeV`$. To estimate the CP phases in Eq. (12), one adopts the following procedure. The real parts(in the approximation that the imaginary parts are very small) are known from the data on rates and asymmetries. The real parts of the amplitudes have also been evaluated in SM with reasonable success with some use of chiral perturbation theory(current algebra and soft pion theorems) and a variety of choices for the baryonic wave functions. The MIT bag model wave function is one such choice which gives conservative results. The same procedure is adopted for calculating the imaginary parts using $`O_6`$. The major uncertainty is in the hadronic matrix elements and the fact that the simultaneous fit of $`s`$ and $`p`$ waves leaves a factor of 2 ambiguity . In the SM, with the Kobayashi-Maskawa phase convention there is no CPV in $`\mathrm{\Delta }I=3/2`$ amplitudes; and for $`\mathrm{\Lambda }`$ and $`\mathrm{\Xi }`$ decays $`\varphi _3=0`$. There is a small electroweak penguin contribution to $`\varphi _3`$ which can be safely neglected. The rate asymmetry is dominated by the s wave amplitudes and the asymmetry $`A_\mathrm{\Lambda }`$ is dominated by the $`\mathrm{\Delta }I=1/2`$ amplitudes. Evaluating the matrix elements in the standard way and with the current knowledge of the K-M matrix elements one finds for the decays $`\mathrm{\Lambda }p\pi ^{}`$ and $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$:
$`\varphi _\mathrm{\Lambda }^s\varphi _\mathrm{\Lambda }^p\mathrm{3.5.10}^4`$
$`\varphi _\mathrm{\Xi }^s\varphi _\mathrm{\Xi }^p\mathrm{1.4.10}^4`$ (17)
With the $`N\pi `$ phase shifts known to be
$$\delta _s\delta _p7^0$$
(18)
one finds for the asymmetry $`A_\mathrm{\Lambda }`$ in the standard model a value of about $`4.10^5`$. For the $`\mathrm{\Xi }\mathrm{\Lambda }\pi ^{}`$ decay mode the phase shifts are not known experimentally and have to be determined theoretically. There are calculations from 1965 which gave large values for $`\delta _s\delta _p20^0`$; however, all recent calculations based on chiral perturbation theory, heavy baryon approximation etc. agree that $`\delta _s\delta _p`$ lies between $`1^0`$ and $`3^0`$ . These techniques have been tested in $`\pi `$-N scattering where they reproduce the known phase shifts within a factor of two. In this case the asymmetry $`A_\mathrm{\Xi }`$ is expected to be $`(0.2`$ to $`0.7)10^5`$. In the Table 1, the SM results for the expected asymmetries in SM are given. Using very crude back of the envelope estimates, similar results are obtained. What is needed is some attention to these matrix elements from the Lattice community.
An experimental measurement of the phase shifts $`\delta _s\delta _p`$ in $`\mathrm{\Lambda }\pi `$ system will put the predictions for $`A_\mathrm{\Xi }`$ on a firmer basis. There is an old proposal due to Pais and Treiman to measure $`\mathrm{\Lambda }\pi `$ phase shifts in $`\mathrm{\Xi }\mathrm{\Lambda }\pi e\nu `$, but this does not seem practical in the near future. Another technique, more feasible, it to measure $`\beta `$ and $`\alpha `$ to high precision in $`\mathrm{\Xi }`$ and $`\overline{\mathrm{\Xi }}`$ decays. Then the combination.
$$(\beta \overline{\beta })/(\alpha \overline{\alpha })=tan(\delta _s\delta _p)$$
(19)
can be used to extract $`\delta _s\delta _p`$. To the extend CP phases are negligible one can also use the approximate relation:
$$\beta /\alpha tan(\delta _s\delta _p)$$
(20)
In $`\mathrm{\Sigma }`$ decays, some asymmetries are quite large but in difficult to measure channels e.g. $`B_\mathrm{\Sigma }`$. In $`\mathrm{\Omega }^{}\mathrm{\Xi }\pi `$ decays the rate asymmetry is larger due to the larger $`\mathrm{\Delta }I=3/2`$ amplitudes. There are no experimental proposals to measure CP asymmetries in $`\mathrm{\Sigma }`$ or $`\mathrm{\Omega }^{}`$ decays at this time.
## 5 Beyond Standard Model
Can new physics scenarios in which the source of CP violation is not K-M matrix yield large enhancements of these asymmetries? We consider some classes of models where these asymmetries can be estimated more or less reliably . It should be kept in mind that any estimates with new physics are at least as uncertain as SM and usually much more prone to uncertainty for obvious reasons.
First there is the class of models which are effectively super-weak . Examples are models in which the K-M matrix is real and the observed CP violation is due to exchange of heavier particles; heavy scalars with FCNC, heavy quarks etc. In all such models direct CP violation is negligible and unobservable and so all asymmetries in hyperon decays are essentially zero. Furthermore, they need to be modified to accommodate the fact that direct CP violation (“Okubo effect”) has now been seen in the kaon decays( the fact that $`ϵ^{}/ϵ`$ is not zero). In the three Higgs doublet model with flavor conservation imposed, the charged Higgs exchange tends to give large effects in direct CP violation as well as large neutron electric dipole moment .
There are two generic classes of left-right symmetric models: (i) Manifest Left- Right symmetric model without $`W_LW_R`$ mixing and (ii) with $`W_LW_R`$ mixing . In (i) $`U_{KM}^L=`$ real and $`U_{KM}^R`$ complex with arbitrary phases but angles given by $`U_{KM}^L`$. Then one gets the “isoconjugate” version in which
$$H_{eff}=\frac{G_FU_{us}}{\sqrt{2}}\left[J_{\mu L}^{}J_{\mu L}+\eta e^{i\beta }J_{\mu R}^{}J_{\mu R}\right]$$
(21)
where $`\eta =m_{WL}^2/m_{WR}^2`$ and $`\beta `$ is the relevant CPV phase. Then $`H_{p.c.}`$ and $`H_{p.v.}`$ have overall phases $`(1+i\eta \beta )`$ and $`(1i\eta \beta )`$ respectively. To account for the observed CPV in K-decay, $`\eta \beta `$ has to be of order $`\mathrm{4.5.10}^5`$. In this model, $`ϵ^{}/ϵ=0`$ and there are no rate asymmetries in hyperon decays but the asymmetries A and B are not zero and e.g. A goes as $`2\eta \beta \mathrm{sin}(\delta _s\delta _p)`$. In the class of models where $`W_LW_R`$ mixing is allowed, the hyperon asymmetries can be enhanced, and also $`ϵ^{}/ϵ`$ is not zero in general (see Table 1).
In MSSM (Minimum Supersymmetric Standard Model) there are new CP violating phases and potentially new contributions to many observables. Until recently the conventional thinking was that the most relevant phase was the one in the squark LL mass terms:
$$m_{12}^2\overline{\stackrel{~}{d}}_L\stackrel{~}{s}_L$$
(22)
and well constrained by $`ϵ`$ so that the contribution to $`ϵ^{}/ϵ`$ would be less than $`2.10^4`$ (similarly for hyperon decays). The new wisdom, painfully learnt after the new results on $`ϵ^{}/ϵ,`$ is that this is not the whole story. There are several ways in which supersymmetric contributions can arise for $`K`$ and hyperon decays.
One example is the lack of degeneracy of $`\stackrel{~}{d}_R`$ and $`\stackrel{~}{u}_R`$ masses. This gives rise to I-spin breaking and in turn can enhance $`ImA_2`$ and contribute to $`ϵ^{}/ϵ`$ at a level of $`10^3`$. For hyperon decays this would lead to a mild enhancement of rate asymmetries but would have no effect on the asymmetries A being probed by E871.
Another possibility is the existence of phases in the L-R squark mass terms. The effect of these on the s-d gluon dipole operator can be parameterised as:
$$\left\{a_{LR}\overline{d}_L\lambda ^a\sigma _{\mu \nu }s_R+a_{RL}\overline{d}_R\lambda ^a\sigma _{\mu \nu }s_L\right\}G_{\mu \nu }^a+h.c.$$
(23)
In terms of $`a_{LR}`$ and $`a_{RL}`$, $`ϵ^{}/ϵ`$ and $`A(\mathrm{\Lambda })`$ can be written as:
$`ϵ^{}/ϵ\alpha `$ $`Im(a_{LR}a_{RL})`$ (24)
$`A(\mathrm{\Lambda })\alpha `$ $`Im(0.2a_{LR}+2.6a_{RL}).`$
The figure shows the range of $`A(\mathrm{\Lambda })`$ for various allowed values of $`a_{LR}`$ and $`a_{RL}`$. Note that $`a_{RL}`$ can yield values for $`A(\mathrm{\Lambda })`$ as large as $`10^3`$ easily probed by E871. This operator is also enhanced in models where CP violation arises thru the exchange of charged scalars such as the Weinberg model.
## 6 Experiments
There have been several proposals to measure hyperon decay asymmetries in $`\overline{p}p\overline{\mathrm{\Lambda }}\mathrm{\Lambda },\overline{p}p\overline{\mathrm{\Xi }}\mathrm{\Xi }`$ and in $`e^+e^{}J/\psi \mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$ but none of these were approved . The only approved and on-going experiment is E871 at Fermilab. In this experiment $`\mathrm{\Xi }^{}`$ and $`\overline{\mathrm{\Xi }}^+`$ are produced and the angular distribution of $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}p\pi ^{}\pi ^{}`$ and $`\overline{\mathrm{\Xi }}^+`$ compared. This experiment effectively measures $`A_\mathrm{\Lambda }+A_\mathrm{\Xi }`$ and will be described in detail by Kam-Biu Luk . To summarize the implications for the measurement of $`A_\mathrm{\Lambda }+A_\mathrm{\Xi }`$ by E871: the SM expectation is about $`4.10^5`$ with a factor of two uncertainty; if new physics should contribute it could be as large as $`10^3`$. A measurement by E871 at the $`10^4`$ level, therefore, will already be a strong discriminant. Eventually, it will be important to know $`A_\mathrm{\Lambda }`$ and $`A_\mathrm{\Xi }`$ separately and the old proposals should be revived.
## 7 $`ϵ^{}/ϵ`$ and Hyperon Decay Asymmetries
It might seem that now that $`ϵ^{}/ϵ`$ has been measured and direct CP violation in $`\mathrm{\Delta }S=1`$ channel been observed, a study of CP violation in hyperon decays is unnecessary and no new information will be obtained. Why is it worthwhile measuring another $`\mathrm{\Delta }S=1`$ process like hyperon decay? The point is that there are important differences and the two are not at all identical. First, there are important differences in the matrix elements. Hyperon matrix elements do not have the kind of large cancellations that plague the kaon matrix elements. The hadronic uncertainties are present for both, but are different. Next, a very important difference is the fact that the K $`\pi \pi `$ decay (and hence $`ϵ^{}`$) is only sensitive to CP violation in the parity violating amplitude and cannot yield any information on parity conserving amplitudes. Hyperon decays, by contrast, are sensitive to both. Thus, $`ϵ^{}/ϵ`$ and hyperon decay CP asymmetries are different and complimentary. The hyperon decay measurements are as important and significant as $`ϵ^{}/ϵ`$.
## Conclusion
The searches for direct CPV are being pursued in many channels: $`\mathrm{\Lambda }N\pi `$, B decays and D decays. Any observation of a signal would be the first outside of $`K^0\overline{K}^0`$ system and would be complimentary to the measurement of $`ϵ^{}/ϵ`$. This will constitute one more step in our bid to confirm or demolish the Standard Kobayashi-Maskawa description of CP violation.
Hyperon decays offer a rich variety of CP violating observables, each with different sensitivity to various sources of CP violation. For example, $`\mathrm{\Delta }_\mathrm{\Lambda }`$ is mostly sensitive to parity violating amplitudes, $`\mathrm{\Delta }_{\mathrm{\Sigma }+}`$ is sensitive only to parity conserving amplitudes, $`A`$ is sensitive to both etc. The size of expected signals vary inversely with the ease of making measurements, i.e. $`\mathrm{\Delta }<A<B`$. Probably because of that, the number of experimental proposals is rather small so far. The one on-going experiment Fermilab E871 can probe $`A`$ to a level of $`10^4`$ which is already in an interesting range. In addition to more experiments, this subject sorely needs more attention devoted to calculating the matrix elements more reliably.
## Acknowledgment
I am grateful to my collaborators Alakabha Datta, Xiao-Gang He, Hitoshi Murayama, Pat O’Donnell, German Valencia and John Donoghue and to members of the E871 collaboration for many discussions. The hospitality of Hai-Yang Cheng, George Hou and their colleagues and staff was memorable and the atmosphere of the conference was most stimulating. This work is supported in part by USDOE under Grant #DE- FG-03-94ER40833.
## References |
warning/0002/nucl-ex0002012.html | ar5iv | text | # REFERENCES
Isotopic Composition of Fragments in Nuclear Multifragmentation
P.M. Milazzo<sup>1</sup>, A.S. Botvina<sup>2,3</sup>, G. Vannini<sup>2</sup>, N. Colonna<sup>4</sup>, F. Gramegna<sup>5</sup>, G. V. Margagliotti<sup>1</sup>, P.F. Mastinu<sup>5</sup>, A. Moroni<sup>6</sup>, R. Rui<sup>1</sup>
<sup>1</sup>Dipartimento di Fisica and INFN, 34127 Trieste, Italy
<sup>2</sup>Dipartimento di Fisica and INFN, 40126 Bologna, Italy
<sup>3</sup>Institute for Nuclear Research, Russian Academy of Science, 117312 Moscow, Russia
<sup>4</sup>INFN, 70126 Bari, Italy
<sup>5</sup>INFN, 35020 Laboratori Nazionali di Legnaro, Italy
<sup>6</sup>Dipartimento di Fisica and INFN, 20133 Milano, Italy
## Abstract
The isotope yields of fragments, produced in the decay of the quasiprojectile in Au+Au peripheral collisions at 35 MeV/nucleon and those coming from the disassembly of the unique source formed in Xe+Cu central reactions at 30 MeV/nucleon, were measured. We show that the relative yields of neutron-rich isotopes increase with the excitation energy in multifragmentation reaction. In the framework of the statistical multifragmentation model which fairly well reproduces the experimental observables, this behaviour can be explained by increasing $`N/Z`$ ratio of hot primary fragments, that corresponds to the statistical evolution of the decay mechanism with the excitation energy: from a compound-like decay to complete multifragmentation.
PACS numbers: 25.70Pq, 25.70-z, 24.60-k
Nuclear fragmentation and its connection to the behaviour of nuclear matter at high excitation energy is the subject of intensive theoretical and experimental investigations . Some general properties of this process are already established: at relatively small excitation energies ($`E^{}`$2–3 A MeV) there is a formation and decay of a long–lived compound-like nucleus system. This process can be described by evaporation/fission–like models. At higher excitation energies (close to the binding energy) there is a complete fast disintegration of the system into fragments. In this case statistical models based on the hypothesis of a nuclear phase transition (simultaneous decay) happen to be very successful .
The parameters of the nuclear system at the break-up have been studied by many methods (see e.g. ). Information about the density in the freeze-out volume is usually extracted from the analysis of velocity correlation functions and kinetic energies of fragments. The temperature of the system can be deduced from different observables: 1) relative populations of unstable nuclear levels, 2) fragments’ kinetic energies, or 3) relative production of isotopes. The last method, based on the statistical properties of double isotope ratios , seems to be the most reliable one, and in fact for the first time it made possible to obtain a nuclear caloric curve as an experimental evidence of a nuclear liquid-gas type phase transition .
Usually these methods need a correction for secondary decay of hot fragments produced in the freeze-out volume . It has to be stressed that the information about chemical composition of hot fragments is of primary importance: depending on their $`N/Z`$ ratio, the decay can in fact proceed differently. Moreover the knowledge of the chemical composition of the primary fragments would allow to more precisely establish the thermodynamical conditions at the freeze-out, e.g., to provide a way to apply the energy balance , as well as it allows to obtain unambiguous data to extract excitation energies of fragments via correlation analysis . Furthermore it can hint for behaviour of the symmetry energy term of the nuclear Equation of State at subnuclear densities and make difference between dynamical and statistical mechanisms of fragmentation .
In this paper we present recent data on the isotope production in heavy-ion collisions at intermediate energies with the aim of studying the isotope content of fragments for different sources and excitation energies during the transition from the low energy decay to the multifragmentation. We’ll show that the behaviour of the experimental isotope yields as a function of source size, isospin and excitation energy can be connected to the corresponding evolution of the $`N/Z`$ ratio of the hot fragments, leading to more insight in the freeze-out condition.
We investigated the Xe+Cu at 30 MeV/nucleon and Au+Au at 35 MeV/nucleon reactions. The experiments were performed at the National Superconducting K1200 Cyclotron Laboratory of the Michigan State University. The angular range 3$`{}_{}{}^{}<\theta _{lab}<`$23 was covered by the $`MULTICS`$ array . The identification thresholds in the $`MULTICS`$ array were about 1.5 MeV/nucleon for charge identification and about 10 MeV/nucleon for mass identification. The MULTICS array consisted of 48 telescopes, each of which was composed of an Ionization Chamber (IC), a Silicon position-sensitive detector (Si) and a CsI crystal. Typical energy resolutions were 2%, 1% and 5% for IC, Si and CsI, respectively. Light charged particles and fragments with charge up to Z=20 were detected at 23$`{}_{}{}^{}<\theta _{lab}<`$160 by the phoswich detectors of the MSU $`Miniball`$ hodoscope . The charge identification thresholds were about 2, 3, 4 MeV/nucleon in the $`Miniball`$ for Z= 3, 10, 18, respectively. The geometric acceptance of the combined array was greater than 87% of 4$`\pi `$.
The multiplicity of detected charged particles (Nc) was used for the reduced impact parameter $`\widehat{b}`$ reconstruction:
$$\widehat{b}=b/b_{max}=\left(_{Nc}^+\mathrm{}P(N^{}c)𝑑N^{}c\right)^{1/2}.$$
Here P(Nc) is the charged particle probability distribution and $`\pi b_{max}^2`$ is the measured reaction cross section for Nc$``$3.
The decay products coming from the decay of the quasiprojectile in peripheral Au+Au 35 MeV/nucleon reaction and those coming from the disassembly of the unique source formed in Xe+Cu 30 MeV/nucleon central collisions have been identified through a careful data selection taking into account for the experimental efficiency distortions on energy and angular distributions . In particular it has been verified that all the detected decay products are emitted nearly isotropically from the same source and that their energy distributions have Maxwellian shapes, i.e. that angular and energy distributions are compatible with a statistical emission, providing an experimental indication that thermalization has been reached . The reconstruction of excitation energies of the sources was carried out analyzing kinematic characteristics of the produced fragments (calorimetric evaluation ).
We studied fragments coming from the excited quasiprojectile Au-like sources produced in peripheral Au+Au collisions with impact parameters 0.95$`>\widehat{b}>`$0.5, corresponding approximately to excitation energies from 3 to 6 MeV/nucleon. Also we considered central ($`\widehat{b}<`$0.2) Xe+Cu collisions which provide sources with excitation around 5.5 MeV/nucleon with approximately the same number of nucleons as Au but with larger charge (the thermal source can be considered as a system after total fusion of colliding nuclei in the center of mass system). In ref. these data were used to obtain a caloric curve. The double isotope ratio method allows to eliminate the need of the knowledge of the initial $`N/Z`$ value of the source, and provides estimates for the temperature. To get information about the composition of hot fragments, new observables should be studied, such as the isotope yields for fixed elements and their evolution with the excitation energy and other parameters of the emitting source. We think that an analysis of the relative isotope production can provide a more reliable information about statistical picture of the process than an analysis of the isobars. The neighbouring isobars (with $`\mathrm{\Delta }`$Z=1) might be produced at different Coulomb barriers, the difference being up to 10 MeV for the Au source. The uncertainty in accounting the real Coulomb energy of the isobars in the freeze-out may essentially exceed the difference in their binding energy ($``$1 MeV) that prevents to conclude unambiguously about thermodynamical parameters.
To avoid a possible problem of preequilibrium in the emission of light charged particles we concentrate mainly on IMFs in this study. In fig. 1 we show the relative isotope yields versus excitation energy obtained from the experimental data for Au sources (each isotope yield is normalized to the total yield for fixed $`Z`$ value). The relative yields change a little as function of the energy, however, one can see a very interesting feature: the relative yields of isotopes with big $`N/Z`$ ratios become larger with increasing excitation energy. In fig. 2 we present the ratios of yields of measured isotopes with the biggest and smallest number of neutrons at fixed $`Z`$ values versus the excitation energy. For all analysed IMFs, the ratio increases considerably in the energy range of $`E^{}`$=3–6 MeV/nucleon. Within such range, the lowest energy corresponds to the onset of multifragmentation with mean IMF multiplicity around one plus a heavy residual, while at the highest energy the fragmentation into many IMFs dominates. Looking at the sources with different $`N/Z`$ ratios we found that the abundances of neutron-rich isotopes are larger for the peripheral quasiprojectile Au than for the central Xe+Cu unique source, at the same excitation energy of about 5.5 MeV/nucleon. The relative yields at this excitation energy are shown in fig. 3. This trend has a natural explanation as the $`N/Z`$ ratio of the Au source is larger than the Xe+Cu one.
If we assume that the yields of the observed isotopes are mainly affected by the secondary decay of hot primary fragments having a slightly larger size, one can consider the presented data for the Au source as an experimental evidence that the $`N/Z`$ ratio of intermediate mass primary fragments increases with the excitation energy of the source. However the physical process behind of this evolution needs a clarification.
It was shown that the statistical multifragmentation model (SMM) very well reproduces the observed charge yields and the $`He`$-$`Li`$ isotope temperatures , as well as the mean fragment kinetic energies. In the following we use the set of SMM parameters which gives the best description of multifragmentation of the sources produced in peripheral collisions . The freeze-out density was taken 1/3$`\rho _0`$ ($`\rho _0`$ is the normal nuclear density). The crucial condition for the present isotope analysis is the requirement of a full description of the charge yields for each considered charge distribution. Under this condition the same parameterization for the central Xe+Cu collisions was used. A possible slight decrease of the source size, as a result of preequilibrium emission, does not affect the conclusions because it hardly changes the $`N/Z`$ ratio of the source ; likewise a small dynamical expansion effect has minor importance. We have checked that the calculated isotope trends (see below) remain stable with respect to reasonable variations of the SMM parameters in the ranges where the charge yields and other observables are reproduced.
We performed a detailed analysis of the isotope production in the framework of this statistical model. Comparison of the SMM predictions with the data is shown in fig. 1, 2 and 3. The qualitative agreement is evident (even quantitative for some important isotopes) and the general trends, especially, the increase with excitation energy of the neutron-rich isotopes with respect to the neutron-deficit ones, are correctly reproduced. However there are few discrepancies in the results which require some justification. In the SMM the Fermi-break-up model is used to describe the secondary decay of fragments with $`A`$16 . It takes into account all ground and nucleon-stable excited states of light fragments and calculates the probabilities of population of these levels microcanonically (according to the available phase space). It does not include matrix elements of the transitions between these states that can be important at small excitation energies. Also the model does not take into account for possible shifts of the nuclear states caused by Coulomb interaction of the excited fragments with the surrounding nuclear matter: these shifts should be calculated in consistent quantum theories. However, in our cases we have a rather high excitation energy of hot primary IMFs: from 2 to 3 A MeV and higher, that is considerably larger than thresholds of the main break-up channels, and the above mentioned problems do not affect the calculated trends. Obviously one can better see the $`N/Z`$ effect in yields of nuclides far from $`\beta `$-stability line, which are less influenced by the shell structure.
According to the SMM predictions at the beginning of multifragmentation the number of primary nucleons in the freeze-out is very small and nearly all available protons and neutrons are bound in hot primary fragments. Their $`N/Z`$ ratio depends on the fragment size. If there are light and heavy fragments in the freeze-out, the light fragments have typically smaller $`N/Z`$ ratio than the heavy ones: these channels are more energetically favorable because of an interplay of the symmetry and Coulomb energies. In fig. 4 we show how the mean $`N/Z`$ ratios for the light ($`Z=810`$) and heavy ($`Z=6870`$) primary fragments evolve with the excitation energy. It is interesting to note that, if a big (residue-like) hot fragment is present in the freeze-out, its $`N/Z`$ ratio can be even larger than the corresponding source ratio because the other light fragments have a considerably lower ratio. At excitations around the multifragmentation threshold ($`E_{thr}^{}=`$3-4 MeV/nucleon) big fragments start to disappear and nearly all available neutrons are combined into hot IMFs giving rise to their neutron content. Finally at very high excitations ($`E^{}`$8 MeV/nucleon) the $`N/Z`$ ratio of the hot IMFs starts to decrease because the number of primary free neutrons increases fastly. This behaviour is responsible for the corresponding trends of the cold fragments produced after deexcitation of these primary IMF (see in Fig.4, e.g., Z=5). Mainly as a consequence of this evolution, we observe an increase of the $`N/Z`$ ratio of the cold fragments in the energy range $`E^{}`$=3–6 MeV/nucleon, and a sizeable change in the relative yields of neutron-rich and neutron-deficit isotopes (see figs. 1, 2). At higher energies this ratio should drop similarly to the hot fragment one.
We should point out that different mechanisms are responsible for the fragment production in the SMM. At small excitation energy light IMFs can be emitted from a compound-like nucleus system. It favours the production of nearly symmetric isotopes with large binding energy (close to $`\beta `$-stability line). Therefore their $`N/Z`$ ratios are usually smaller than the ratio of the source. The probability of the evaporation of IMF is small, however, it contributes to the yield at $`E^{}E_{thr}^{}`$. At excitation energies higher than the threshold the multifragmentation sets in: from a fast break-up into two hot fragments it evolves towards the break-up into three or more fragments with the increase of the source excitation energy . At the multifragmentation the secondary decay of hot primary fragments is the main process defining a relative abundance of particular isotopes. In the SMM calculations the contribution of the evaporated isotopes favours the low ratios presented in fig. 2 at small $`E^{}`$. However, assuming that only the evaporation mechanism exists at high excitations (independent from the reproduction of the charge yield) the observed increase of the ratio can not be explained. That supports the suggested evolution of the decay mechanisms and isotope composition of hot fragments.
It was shown in Ref. that an increase of the $`N/Z`$ source ratio leads to increasing the relative yields of neutron-rich isotopes, in agreement with present results. In their analysis they extract also information about neutron-to-proton ratio ($`n/p`$) at the freeze-out. In the present calculations with SMM we found that for central (Xe+Cu) and peripheral (Au) sources with the excitation of 5.5 A MeV this ratio increases more than increasing the $`N/Z`$ ratio of the corresponding sources, in agreement with . A possible production of neutron-rich hot primary fragments is also predicted by other theoretical models. Dynamical stochastic mean field calculations for Au source predict hot fragments with the same $`N/Z`$ ratio as the SMM. Also an analysis of the correlation functions performed in comes to the conclusion that the larger neutron content of the hot primary fragments corresponds to the experimental data.
In summary, we presented new data on yields of isotopes produced after decay of the Au and Xe+Cu sources at excitation energy range of 3–6 MeV per nucleon, that is around and slightly above the multifragmentation threshold. We found that the experimental relative yields of neutron–rich isotopes increase with excitation energy for the Au sources. The SMM calculations reproduce the whole set of data well enough to support the statistical picture realized in this model. In this approach the energy dependence of the isotopic composition of the produced fragments can be explained in terms of a transition from an evaporation-like emission of few fragments to the total multifragmentation break-up which leads to the increase of neutron content of hot primary intermediate mass fragments.
A.S. Botvina thanks the Istituto Nazionale di Fisica Nucleare (Bologna section, Italy) for hospitality and support.
Figure captions
Fig.1: Relative yields of isotopes of different elements versus excitation energy of Au source. Symbols are experimental data, lines are SMM calculations: <sup>6</sup>Li, <sup>7</sup>Be, <sup>10</sup>B, <sup>11</sup>C (open circles, dashed line); <sup>7</sup>Li, <sup>9</sup>Be, <sup>11</sup>B, <sup>12</sup>C (full circles, solid line); <sup>8</sup>Li, <sup>10</sup>Be, <sup>12</sup>B, <sup>13</sup>C (open squares, dot-dashed line); <sup>14</sup>C (full squares, dotted line). The experimental uncertainties on the excitation energy are the same as shown in fig. 2; error bars on relative yields are smaller than symbols size.
Fig.2: Ratio of relative yields of neutron-rich to neutron-deficit isotopes of Li, Be, B and C fragments versus excitation energy of Au source. Solid circles are the experimental data, while lines refer to SMM calculations; solid squares refer to central Xe+Cu experimental data.
Fig.3: Relative yields of different isotopes for fragments with charges from $`Z`$=1 to $`Z`$=6. Circles are experimental data: the solid ones are for the Au system, the open ones are for the Xe+Cu system at the excitation energy of 5.5 MeV/nucleon. Solid and dashed lines are the corresponding SMM calculations.
Fig.4: The SMM calculations of the mean neutron–to–proton ($`N/Z`$) ratio of hot primary fragments produced at the freeze-out (full lines) and the cold fragments produced after the secondary decay (dot-dashed line) versus excitation energy for the Au source. |
warning/0002/cond-mat0002033.html | ar5iv | text | # A microscopic approach to phase transitions in quantum systems
## Abstract
We present a new theoretical approach for the study of the phase diagram of interacting quantum particles: bosons, fermions or spins. In the neighborhood of a phase transition, the expected renormalization group structure is recovered both near the upper and lower critical dimension. Information on the microscopic hamiltonian is also retained and no mapping to effective field theories is needed. A simple approximation to our formally exact equations is studied for the spin-$`S`$ Heisenberg model in three dimensions where explicit results for critical exponents, critical temperature and coexistence curve are obtained.
Several physical systems, ranging from magnets to superfluids and superconductors, display rich phase diagrams in a temperature regime where quantum effects cannot be neglected. Different scenarios, characterized by competing order parameters and zero temperature phase transitions have been recently advocated also in the framework of high temperature superconductivity where antiferromagnetic order, Cooper pairing and, possibly, phase separation are at play in the same region of the phase diagram . A satisfactory understanding of phase transitions in quantum models has been attained years ago through the seminal work by Hertz who showed that, at low energy and long wavelengths, quantum models may be described by a suitable classical action. However, a quantitative theory of the thermodynamic behavior is still lacking and we mostly rely on mean field approaches or weak coupling renormalization group (RG) calculations , applied to quantum systems via the mapping to the appropriate effective field theory. In particular, the interplay between thermal and quantum fluctuations is expected to give rise to crossover phenomena whose extent strongly depends on the microscopic features of the system. Even for the most extensively studied models, like the Heisenberg antiferromagnet, our knowledge of the phase diagram is in fact limited, and the first precise finite temperature simulation attempting to fill this gap has become available only recently . By contrast, in classical models, numerical simulations are quite efficient even in the neighborhood of critical points and, from the analytical side, microscopic approaches especially devised for the quantitative description of the phase diagram of classical fluids and magnets are available. For instance, the hierarchical reference theory of fluids (HRT) has proven quite accurate in locating the phase transition lines both in lattice and in continuous models.
In this Letter we sketch the derivation of the quantum hierarchical reference theory of fluids (QHRT) which we then apply to the Heisenberg antiferromagnet. We will demonstrate that the known renormalization group equations near four and near two dimensions are naturally recovered within our approach, which therefore unifies two complimentary techniques. On approaching the critical point, the spin velocity vanishes according to the expected dynamical critical exponent for an antiferromagnet. Finally, the phase diagram of this model in three dimensions is computed by numerical integration of a simple approximation to the the QHRT equations, providing a concrete application of our general approach.
The starting point is a microscopic, many body hamiltonian $`H`$ written as the sum of a reference part $`H_0`$ and an interaction term $`V`$. The interaction is assumed to be bilinear in some operator $`\rho (r)`$, which is assumed either linear in bosonic operators or quadratic in fermionic ones:
$$V=\frac{1}{2}𝑑x𝑑y\rho (x)w(xy)\rho (y)$$
(1)
with a non singular (i.e. Fourier transformable) two body potential $`w`$. The properties of the reference system under the action of an external field $`h`$ coupled to the order parameter $`\rho (r)`$ are supposed known. No specific assumption on the reference system is made: in particular we do not need that $`H_0`$ corresponds to a non interacting system, where Wick theorem applies (such a feature is crucial in setting up the QHRT equations). These requirements are indeed rather general and include several models of current interest in many body physics: quantum magnets (where $`\rho (r)`$ represents the local spin variable), fermionic systems, like the Hubbard model, or even the Holstein model for the electron-phonon problem.
The first task is to build up a formal perturbative expansion of the partition function of the model: $`Z=\mathrm{Tr}\mathrm{exp}(\beta H)`$. Following a standard procedure , $`Z/Z_0`$ can be written as the average over the reference distribution function of an imaginary time evolution operator $`U(\beta )`$. When this operator is written as a power series of the interaction $`w(r)`$, we formally recover a perturbative expansion identical to that of a classical partition function for a $`(d+1)`$ dimensional model. The additional “temporal” dimension is limited to the interval $`(0,\beta )`$ and $`w(r)\delta (t)/\beta `$ plays the role of classical two body interaction $`w_c(r,t)`$. The reference system of the associated classical model is implicitly defined by requiring that its correlation functions coincide with those of the quantum reference hamiltonian $`H_0`$. Approximate mappings between a quantum model and an effective classical system have been proposed and studied in the literature in order to clarify the role of thermal and quantum fluctuations. The novelty of our approach is that $`i)`$ it is exact and $`ii)`$ it applies to all temperature regimes, including the $`T0`$ limit. Having reduced the quantum problem to a classical one, we can directly apply the techniques developed in that framework. In particular, the already mentioned HRT is an implementation of the momentum space renormalization method which preserves information on the details of the microscopic hamiltonian. In HRT, different Fourier components of the two body interaction are included gradually, starting from the shortest wavelength: physically this corresponds to a smooth turning on of fluctuations over larger and larger lengthscales. This procedure can be carried out exactly by defining a sequence of auxiliary systems interacting via a potential whose Fourier components coincide with $`w(k)`$ for $`k>Q`$ and vanish elsewhere. As a result, we obtain a set of coupled differential equations expressing the change in the free energy and in the correlation functions of the model when a given Fourier component $`Q`$ of the potential is included. The initial condition represents the system in which fluctuations are frozen, and in fact coincides with the known mean field result. As the cut-off wavevector $`Q0`$, the fully interacting system is recovered. Mean field approximation therefore corresponds to neglecting the change in the properties of the model as described by our differential equations. Details can be found in Ref. .
As an example, we study the the spin-$`S`$ antiferromagnetic Heisenberg model on a hypercubic lattice in $`d`$ dimension:
$$H=H_0+V=h\underset{𝐑}{}e^{i𝐠𝐑}S_𝐑^z+J\underset{<𝐑,𝐑^{}>}{}𝐒_𝐑𝐒_𝐑^{}$$
(2)
where the sum is restricted to nearest neighbors and $`𝐠`$ is the antiferromagnetic wavevector of components $`g_i=\pi `$. In applying the HRT approach to a quantum model, we decided to impose the cut-off $`Q`$ only to the spatial Fourier components of $`w_c(r,t)`$. Different physical models might require other choices of the cut-off, whose only role is to continuously connect the reference system to the fully interacting one, which is recovered in the $`Q0`$ limit. The exact evolution equation for the Helmholtz free energy density $`a`$ of the system describes how $`a`$ is modified due to a change in the cut-off $`Q`$:
$`{\displaystyle \frac{da^Q}{dQ}}={\displaystyle \frac{1}{2\beta }}{\displaystyle _{𝐩\mathrm{\Sigma }_Q}}{\displaystyle \underset{\omega }{}}`$ $`\{`$ $`2\mathrm{ln}\left[(1F_{xx}^Q(𝐩,\omega )w(𝐩))(1+F_{xx}^Q(𝐩^{},\omega )w(𝐩))F_{xy}^Q(𝐩,\omega )F_{xy}^Q(𝐩^{},\omega )w(𝐩)^2\right]`$ (3)
$`+`$ $`\mathrm{ln}\left[(1F_{zz}^Q(𝐩,\omega )w(𝐩))(1+F_{zz}^Q(𝐩^{},\omega )w(𝐩))\right]\}`$ (4)
Here $`𝐩^{}=𝐩+𝐠`$, $`w(𝐩)=2J\gamma (𝐩)=2J_i\mathrm{cos}(p_i)`$ is the Fourier transform of the interaction, the summation is over the Matsubara frequencies $`\omega _n=2\pi n/\beta `$, the $`(d1)`$ dimensional integral is restricted to the surface $`\mathrm{\Sigma }_Q`$ defined by $`\gamma (𝐩)=\sqrt{d^2Q^2}`$ and the functions $`F_{ij}^Q(𝐩,\omega )=<S^i(𝐩,\omega )S^j(𝐩,\omega )>`$ are the Fourier transforms of the spin-spin dynamical correlation functions (in imaginary time) for a system where only fluctuations of wavevector $`𝐩`$ such that $`|\gamma (𝐩)|<\sqrt{d^2Q^2}`$ are included. The isotropy of the model implies $`F_{yy}=F_{xx}`$ and $`F_{yx}=F_{xy}`$. Analogous equations can be derived for the many spin correlation functions of the model. The (infinite) set of differential equations forms the QHRT hierarchy. When $`Q=d`$, fluctuations are neglected and the exact initial condition for the first QHRT equation (4) coincides with the mean field free energy density. The magnetic structure factors at $`Q=d`$ can be explicitly written as:
$`F_{xx}^Q(𝐩,\omega )`$ $`=`$ $`{\displaystyle \frac{\mu _{}w(𝐩)}{m^2\omega ^2+\mu _{}^2w(𝐩)^2}};`$ (5)
$`F_{xy}^Q(𝐩,\omega )`$ $`=`$ $`{\displaystyle \frac{m^1\omega }{m^2\omega ^2+\mu _{}^2w(𝐩)^2}};`$ (6)
$`F_{zz}^Q(𝐩,\omega )`$ $`=`$ $`{\displaystyle \frac{\delta _{\omega ,0}}{\mu _{||}+w(𝐩)}}`$ (7)
where $`m`$ is the staggered magnetization and $`\mu _{},\mu _{||}`$ are known functions of $`m`$. For $`S=1/2`$: $`\mu _{}=2Tm^1\mathrm{tanh}^1(2m)`$ and $`\mu _{||}=4T(14m^2)^1`$. Note that the dependence of the transverse magnetic structure factors in (7) on frequency and momentum is consistent with the first order spin wave result at zero temperature and reproduces the known single mode approximation which well represents antiferromagnetic correlations at low temperatures . The longitudinal correlations in equation (7) are instead purely classical and, as such, satisfy the relationship $`T\chi _{||}(k)=S_{||}(k)`$ between longitudinal susceptibility $`\chi _{||}(k)`$ and the corresponding static structure factor $`S_{||}(k)`$.
Equation (4), although formally exact, is not closed because the evolution of the free energy depends on the unknown magnetic structure factors of the model $`F_{ij}^Q(𝐩,\omega )`$. Therefore, we have to introduce some approximate parametrization of the structure factors in terms of the free energy. The simple approximation we have studied is to retain the form of $`F_{ij}^Q(𝐩,\omega )`$ as given in Eq. (7) but imposing thermodynamic consistency in order to determine the two scalar parameters $`\mu _{}`$ and $`\mu _{||}`$. More precisely, for every $`Q`$ we related the transverse and longitudinal staggered susceptibilities to the free energy via the exact sum rules:
$`(\mu _{}2dJ)^1`$ $`=`$ $`F_{xx}^Q(𝐠,\omega =0)=m(a^Q/m)^1`$ (8)
$`(\mu _{||}2dJ)^1`$ $`=`$ $`F_{zz}^Q(𝐠,\omega =0)=(^2a^Q/m^2)^1`$ (9)
From the adopted structure of the dynamical correlation functions, we also obtain the relationship between the parameters entering $`F_{ij}^Q`$ and the zero temperature non linear sigma model coupling constants: uniform transverse susceptibility $`\chi _0=1/(4d)`$, spin wave velocity $`c=\sqrt{4d}m`$ and spin stiffness $`\rho _s=m^2`$. The hydrodynamic relation $`\chi _0c^2=\rho _s`$ is automatically satisfied by our ansatz for arbitrary spontaneous magnetization $`m`$. It is interesting to note that from our parametrization, the scaling of the spin wave velocity $`c`$ on approaching the critical temperature ($`t=(T_cT)/T_c0`$) gives $`cm`$ that is $`ct^\beta `$ along the coexistence curve. The dynamic scaling hypothesis instead predicts $`ct^{\nu (z1)}`$ where $`z`$ is the dynamical critical exponent. By use of scaling laws and recalling that in our approximation the correlation critical exponent vanishes ($`\eta =0`$), we get $`z=d/2`$ which is the expected result for an antiferromagnet (i.e. model G) . Equation (4), together with (7) and (9) give rise to a partial differential equation for the free energy density of the Heisenberg model $`a^Q(m)`$ as a function of the cut-off $`Q`$ and of the magnetization $`m`$. The frequency sum can be carried out analytically giving the final equation:
$$\frac{da^Q}{dQ}=\frac{1}{2\beta }_{𝐩\mathrm{\Sigma }_Q}\left\{4\mathrm{ln}\left[\frac{\mathrm{sinh}\left(\frac{1}{2}\beta m\mu _{}\right)}{\mathrm{sinh}\left(\frac{1}{2}\beta m\sqrt{\mu _{}^2w(𝐩)^2}\right)}\right]+\mathrm{ln}\left[\frac{\mu _{||}^2}{\mu _{||}^2w(𝐩)^2}\right]\right\}$$
(10)
We numerically solved this partial differential equation for several values of the spin $`S`$ and different temperatures in order to study the phase diagram of this system. Note that $`S`$ just enters the theory through the initial condition $`a^Q(m)`$ at $`Q=d`$, while the form of the differential equation is unaffected by $`S`$. Before showing the numerical results, however, it is useful to discuss the behavior of Eq. (10) near a phase transition. In particular, we studied the neighborhood of the critical point (in $`d>2`$) and the low temperature region. Both at the critical point and along the coexistence curve the susceptibilities diverge due to the presence of critical fluctuations and Goldstone bosons, respectively. From Eq. (9) we conclude that at long wavelengths (i.e. $`Q0`$) and near a phase transition we have $`\mu _{}\mu ||2dJ`$. In this region equation (10) simplifies and, by rescaling the free energy as $`a^QQ^d`$ and the magnetization as $`mQ^{(2d)/2}`$, it reduces to the RG equation obtained by Stanley et al. for a $`O(3)`$ symmetric $`\varphi ^4`$ hamiltonian. Such an equation has been analyzed near four dimension and proved to give the correct critical exponents to first order in the $`ϵ=4d`$ expansion. In three dimensions, the numerical solution of the universal fixed point equation gives for the correlation length critical exponent the result $`\nu =0.826`$, to be compared with the accepted value $`\nu =0.71`$. The other critical exponents follow from the scaling laws, noting that our analytical form of the two point functions (7) forces the anomalous dimension exponent to vanish $`\eta =0`$. We therefore find non classical exponents in three dimensions. Special care must be paid when dealing with the $`T0`$ limit of our equation. In this case, the asymptotic form of the equation changes and it can be shown to give rise, near a hypothetical quantum critical point, to critical exponents appropriate for a $`O(3)`$ model in $`d+1`$ dimensions as expected. A separate analysis should be carried out in the low temperature phase. If symmetry is spontaneously broken, following Chakravarty et al. , we may ask how quantum and thermal fluctuations modify the zero field magnetization. In order to answer this question we perform a Legendre transform on our equation (10): we first derive it with respect to the magnetization $`m`$ obtaining an evolution equation for the magnetic field $`h^Q(m)`$ at fixed $`m`$. Then we find the equation governing the evolution of the spontaneous magnetization $`m^Q`$ implicitly defined by the requirement $`h^Q(m^Q)=0`$ at every $`Q`$. This procedure gives rise to a differential equation for $`m^Q`$. In the $`Q0`$ limit, taking into account that the longitudinal susceptibility diverges more slowly than $`Q^2`$, QHRT reduces to a simple ordinary differential equation:
$$\frac{dm^Q}{dQ}=K_d\left(\frac{Q}{\sqrt{d}}\right)^{d2}\left[\mathrm{tanh}\left(Q\beta m^Q\right)\right]^1$$
(11)
where $`K_d`$ is a geometrical factor (ratio between the solid angle and volume of the Brillouin zone). By introducing the rescaled variable $`g=\sqrt{4d}(Q/\sqrt{d})^{d1}/m^Q`$, equation (11) becomes identical, to order $`g^2`$, to the known weak coupling RG equations for the non linear sigma model applied by Chakravarty et al. to the analysis of the antiferromagnetic Heisenberg model at long wavelengths and low temperatures. As an example, we plot in Fig. 1 the RG flux of $`g^Q`$ obtained by the integration of the full QHRT equation (10). We clearly see the effect of the unstable zero temperature weak coupling fixed point while, for the nearest neighbor Heisenberg model, the other fixed point ($`g_c`$), governing the quantum critical regime, has no effect on the RG trajectories. This analysis shows that the single mode approximation to QHRT reproduces the correct long wavelength structure both near four and two dimensions. Furthermore, we expect QHRT to be superior to the weak coupling renormalization group equations because our non perturbative approach also describes the critical region and the high temperature regime where $`m^Q0`$ as $`Q0`$, corresponding to $`g\mathrm{}`$ i.e. to the strong coupling phase of the non linear sigma model.
Finally, we present few results of the numerical integration of Eq. (10) in three dimensions. As already pointed out in the classical case , the HRT approach is able to correctly implement Maxwell construction at first order phase transitions and in fact the free energy density $`a^Q(m)`$ at the end of the integration, i.e. in the $`Q0`$ limit, becomes rigorously flat in a finite region of the magnetization axis for $`T<T_c`$. Therefore it is easy to extract from the numerical output the critical temperature and the coexistence curve, shown in Fig. 2 for several values of the spin $`S`$. The zero temperature limits of this curve agree within about $`2\%`$ with the accepted estimates based on spin wave theory, Monte Carlo simulations or series expansions. Regrettably, for the spontaneous magnetization at finite temperature there are just few available results going beyond mean field approaches. Simulation data for the classical $`S\mathrm{}`$ case and recent series expansion for the $`S=1/2`$ model seem to give somewhat larger coexistence regions. However, we believe that a more systematic analysis of these models by accurate numerical techniques is necessary before reaching a definite conclusion on the accuracy in the determination of the coexistence curve. The critical temperature for the classical model is known by several methods to be $`T_c=1.443J`$ while for the $`S=1/2`$ case it has been recently estimated as $`T_c=0.946J`$ by use of a newly developed Quantum Monte Carlo method and $`T_c=0.93J`$ by series expansions. Our results are a few percent lower, being $`T_c=1.419J`$ for $`S=\mathrm{}`$ and $`T_c=0.90J`$ for $`S=1/2`$. From the solution of the QHRT equation we also obtain other important information on the model, for instance, the equation of state, the specific heat and also the temperature dependent dynamical structure factors, via analytic continuation of the adopted expressions (7). In order to improve the QHRT results we have just discussed, other approximate expressions for the magnetic structure factors should be examined, possibly keeping the same form (7) but allowing for a non trivial renormalization factor for the uniform susceptibility. This method can be applied in a straightforward way to other models of interest in quantum many body physics, like the Hubbard model, and may help to determine the location of the magnetic phase transitions and the possible occurrence of phase separation in a purely repulsive electron system. |
warning/0002/hep-th0002149.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Anomaly and tadpole vanishing conditions have been recognized for long as intriguingly related constraints in the construction of consistent vacuum configurations with open and unoriented strings. Indeed infinity and anomaly cancellations go hand by hand in the $`SO(32)`$ type I theory . Although considerable experimental evidence has been accumulated in favour of a one-to-one correspondence between the two kinds of consistency conditions in models with $`𝒩=(1,0)`$ supersymmetry in $`D=6`$, already in $`D=4`$ there are exceptions to this rule and very little is known so far in the non-supersymmetric case. At first sight it is not obvious why the vanishing of RR tadpoles, a transverse-channel constraint on the closed-string coupling to boundaries and crosscaps, should in general be related to the vanishing of irreducible anomalies, a constraint on the spectrum that is naturally encoded in the direct-channel amplitudes.
The issue has been first addressed in where the RR tadpole in $`D=10`$ was shown to induce a violation of BRST invariance on the string worldsheet. In tadpole cancellation was proposed as the open-string analogue of closed-string modular invariance, that cuts off the UV region responsible of anomalies. Additional evidence for a tadpole-anomaly correspondence was given in , where some $`𝒩=(1,0)`$ supersymmetric models in $`D=6`$ and some non-supersymmetric models in $`D=10`$ were shown to be free of irreducible anomalies thanks to the vanishing of RR tadpoles. The cancellation of the left-over reducible anomalies relied on the presence of extra RR antisymmetric tensors that participate in a generalized Green-Schwarz mechanism . The general belief on an IR-UV correspondence between RR tadpoles and anomalies has been reinforced by the advent of D-branes . RR charge neutrality of vacuum configurations with closed unoriented strings requires the introduction of D-branes and their open-string excitations. The new trend motivated the study of , where the correspondence was analyzed for supersymmetric D-brane configurations at orbifold singularities in non-compact spaces. RR tadpoles from twisted sectors were identified with irreducible gauge anomalies. More recently, an important contribution has been given in , where a careful comparison has been drawn in the context of geometric (supersymmetric) orbifold compactifications to $`D=4,6`$. In all cases studied in , RR tadpoles associated to string amplitudes in sectors where the orbifold group acts without fixed tori have been put in a one-to-one correspondence with irreducible anomalies. The remaining tadpole conditions have been left as seemingly unrelated to the conditions for anomaly cancellation. The intrisic reasons for such a tadpole-anomaly correspondence and the range of its validity inside the more general web of open string compactifications remains unclear.
The increasing interest in the phenomenological perspectives of open string vacuum configurations with various patterns of supersymmetry breaking deserves a thorough analysis of this long-standing problem. In the present paper we address the problem in full generality. Our analysis relies on the topological properties of the internal SCFT (super-conformal field theory) and encompasses not only unoriented descendants of geometric orbifolds and Gepner models, but also descendants of asymmetric orbifolds and free fermionic constructions that result in left-right symmetric parent theories. We find that RR-tadpole conditions always correspond to the cancellation of some sort of anomaly in the effective theory. The contribution to the anomalies from a given sector of the internal SCFT is measured by its “chiral” Witten index, $`=tr_R()^F`$ . This index is defined by a chiral vacuum amplitude in the odd spin structure, where the worldsheet supercurrent is periodic, and effectively counts the chiral asymmetry of each sector <sup>1</sup><sup>1</sup>1For definitions and applications of the index in the open string context see . Tadpoles associated to sectors with non-vanishing Witten index will be precisely identified with irreducible terms in the canonical anomaly polynomial. This extends the notion of sectors without fixed tori that appear in the orbifold constructions discussed in to the present more general context. In addition an attentive study of the tadpoles that arise from sectors with vanishing Witten index suggests their connection to anomalous amplitudes involving a larger number of external insertions.
We would like to stress that nowhere in the paper we will assume that the theory enjoys target space supersymmetry. Our discussion is completely general and applies to non-supersymmetric models that can arise for instance from superstring compactifications with or without brane supersymmetry.
In order to keep the notation as compact as possible, the language we use is the one of rational superconformal field theories but with obvious modifications one can accomodate in this setting all “irrational” models studied so far such as tori and orbifolds thereof. The only crucial ingredient in our considerations is the worldsheet consistency between direct (loop) and transverse (tree) channel for the relevant amplitudes. This imposes highly non-trivial constraints in the construction of a sensible open string model. Throughout the paper we will assume that a solution to these fundamental constraints (not to be confused with the tadpole-anomaly conditions under consideration) has been found. Another important ingredient is the “modular invariance” of string amplitudes in the odd-spin structure which allows us to relate irreducible anomalies in the direct channel to RR tadpoles in the transverse channel.
After showing the equivalence between RR tadpoles and irreducible terms in the anomaly polynomial, we will pass to the study of the reducible ones. We parametrize the anomaly polynomials in terms of reflection coefficients for the massless closed string states in front of crosscaps and boundaries. We will find that, barring some minor ambiguities in $`D=8,10`$ if the Chan-Paton group is non semi-simple, the anomaly polynomial admits a unique factorization in terms of a sum of as many products as sectors that flow in the transverse channel and have non-vanishing Witten index. This allow us to factorize the reducible anomaly in such a way as to immediately expose the R-R fields that participate in the generalized GSS mechanism and to extract the corresponding terms in the effective lagrangian. In particular one can immediately identify the “anomalous” $`U(1)`$’s.
The plan of the paper is as follows. In Section 2 we briefly review the construction of open string descendants. In Section 3 we establish the precise relation between RR tadpoles for sectors with non-vanishing Witten index and irreducible anomalies. In Section 4 we study the reducible part of the anomaly polynomials in any even dimensions and extract the WZ couplings in the effective theories. In Section 5 we discuss RR tadpoles arising from sectors with vanishing Witten index and elaborate on their connection to anomalous amplitudes with higher number of external legs. In Section 6 we discuss our results and add some concluding remarks.
## 2 Open string descendants
In this section we review some relevant features in the construction of open string descendants of generic left-right symmetric “compactifications” of type II and type 0 superstrings. By superstring compactification we mean any vacuum configuration whose worldsheet dynamics is governed by a superconformal field theory (SCFT) irrespective of the presence of unbroken target-space supersymmetry.
A closed string compactification to $`D`$ dimensions is defined by tensoring an internal SCFT, with left and right moving central charges $`(c,\overline{c})=(\frac{3}{2}(10D),\frac{3}{2}(10D))`$, with a $`(c,\overline{c})=(\frac{3}{2}D,\frac{3}{2}D)`$ spacetime theory realized in terms of free worldsheet bosons and fermions $`\{X^\mu ,\psi ^\mu ,\stackrel{~}{\psi }^\mu \}`$. The Hilbert space of string states can be decomposed into a (generically infinite) sum $`=_{i,j}_i\overline{}_j`$. The index $`i`$ labels the primary fields $`\mathrm{\Phi }^i`$ of some chiral algebra $`𝒢`$, that includes an $`𝒩=1`$ superconformal algebra at least. The left-moving subspace $`_i`$ thus consists in the tower of descendants under $`𝒢`$ of the groundstate $`|h_i=\mathrm{\Phi }^i|0`$. The spectrum of the left-moving Hamiltonian $`L_0`$ in the sector $`_i`$ is conveniently assembled in the holomorphic character
$$𝒳_i(\tau )=\mathrm{Tr}__i^{}q^{L_0\frac{c}{24}},$$
(1)
with $`q=e^{2\pi i\tau }`$. A prime in (1) indicates the omission of the contribution of the bosonic zero-modes to the trace. Moreover, it is always understood that states in the NS (sub)sector, where the worldsheet supercurrent is anti-periodic, enter with a plus sign and states in the R (sub)sector, where the worldsheet supercurrent is periodic, enter with a minus sign. This follows from modular invariance at two loops and implements the correct relation between spin and statistics. For the right-moving excitations on the string worldsheet we assume $`\overline{𝒢}𝒢`$ and define anti-holomorphic characters $`\overline{𝒳}_j(\overline{\tau })`$ similarly.
The Lorentz transformation properties of states in $`_i`$ are encoded in the spacetime part of the character $`𝒳_i`$, that can be decomposed into “characters” of the $`SO(D2)`$ little group
$`\chi _O+\chi _V`$ $`=`$ $`\left({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{0}{0}\right]}{\eta ^3}}\right)^{\frac{D2}{2}}\chi _O\chi _V=\left({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{0}{\frac{1}{2}}\right]}{\eta ^3}}\right)^{\frac{D2}{2}}`$
$`\chi _S+\chi _C`$ $`=`$ $`\left({\displaystyle \frac{\vartheta \left[\genfrac{}{}{0pt}{}{\frac{1}{2}}{0}\right]}{\eta ^3}}\right)^{\frac{D2}{2}}\chi _S\chi _C=\left({\displaystyle \frac{i\vartheta \left[\genfrac{}{}{0pt}{}{\frac{1}{2}}{\frac{1}{2}}\right]}{\eta ^3}}\right)^{\frac{D2}{2}}.`$ (2)
The labels $`O,V,S,C`$ denotes the scalar, vector, spinor left (L) and spinor right (R) representation of the affine transverse Lorentz algebra at level $`\kappa =1`$. Once the contributions of the spacetime bosonic and fermionic coordinates are included, modular invariance ensures that any holomorphic character is given by a sum over spin structures:
$$𝒳_i=\underset{\alpha ,\beta =0,\frac{1}{2}}{}\left(\frac{\vartheta \left[\genfrac{}{}{0pt}{}{\alpha }{\beta }\right]}{\eta ^3}\right)^{\frac{D2}{2}}_i\left[\genfrac{}{}{0pt}{}{\alpha }{\beta }\right]=_i\mathrm{\Theta }_D+\mathrm{}.$$
(3)
and similarly for the antiholomorphic part. In the second equality we have isolated the contribution of the odd spin structure on which our analys will be focussed. In this spin-structure worldsheet bosons and fermions share the same boundary conditions. The worldsheet supercurrent is periodic and the contributions of the massive excitations to the sum cancel against each other. The contribution of the spacetime super-coordinates,
$$\mathrm{\Theta }_D=\left(\frac{i\vartheta \left[\genfrac{}{}{0pt}{}{\frac{1}{2}}{\frac{1}{2}}\right]}{\eta ^3}\right)^{\frac{D2}{2}},$$
(4)
though formally zero, is a convenient book-keeping of the chiral asymmetry. The result of the trace in each sector of the internal SCFT theory, $`_i_i\left[\genfrac{}{}{0pt}{}{\frac{1}{2}}{\frac{1}{2}}\right]`$, is an integer (the Witten index ) that counts the difference between the number of bosonic and fermionic ground-states. From the right-hand side of (3) we notice that $`_i`$ effectively counts the difference between the number of left- ($`L`$) and right- ($`R`$) “handed” spacetime fermions.
The spectrum of perturbative closed string states is then packaged in the one-loop (torus) partition function. Neglecting overall volume factors and the modular integration ($`_{}\frac{d^2\tau }{\tau _2^2}`$) one finds
$$𝒯=\tau _2^{\frac{D2}{2}}\underset{ij}{}𝒯_{ij}𝒳_i\overline{𝒳}_j.$$
(5)
$`𝒯_{ij}`$ are positive integer coefficients with $`𝒯_{00}=1`$ for the uniqueness of the identity sector, i.e. of the graviton. The powers of $`\tau _2`$ arise from the momentum integrations in the non-compact directions. The condition for left-right symmetry on the worldsheet translates into the constraints $`𝒯_{ij}=𝒯_{ji}`$. The coefficients $`𝒯_{ij}`$ are also highly restricted by one-loop modular invariance. The characters are known to provide a unitary representation of the modular group $`SL(2,\text{Z}\text{Z})`$ generated by the transformations $`T`$ and $`S`$ under which
$`T:𝒳_i(\tau +1)`$ $`=`$ $`e^{2\pi i(h_ic/24)}𝒳_i(\tau )`$
$`S:𝒳_i({\displaystyle \frac{1}{\tau }})`$ $`=`$ $`{\displaystyle \underset{j}{}}(i\tau )^{\frac{D2}{2}}S_{ij}𝒳_j(\tau ).`$ (6)
After the resolution of fixed-point ambiguities, related to the presence of different sectors with the same character, the transformation $`S`$ is represented by a symmetric matrix that satisfies $`(ST)^3=S^2=C`$, with $`C`$ the charge conjugation matrix.
In most of our subsequent discussion, we will restrict our attention to the case in which the boundary conditions associated to the introduction of open and unoriented strings preserve the diagonal combination of the worldsheet symmetries of the bulk theory together with their target-space byproducts. This generalizes the notion of a BPS D-brane. With a slight change of notation, that amounts to decomposing the characters of the parent theory into their irreducible components with respect to some lower (super-)symmetry, one can easily accomodate models with brane (super-)symmetry breaking. This generalizes the notion of a non-BPS D-brane. Moreover we will display formulas which are more akin to the rational context, where the number of characters is finite. With some additional care, the results can be adapted to irrational contexts such as toroidal or orbifold compactifications at irrational values of the moduli (radii, shapes and Wilson lines).
At generic points of the moduli space of toroidal compactifications the index $`i`$ may be thought to run over an infinite number of primary fields, one for each independent choice of internal momenta and windings. Similarly in orbifold compactifications, only states that are invariant under the action of the orbifold group will enter the trace (1). Chiral string excitations are then organized according to their eigenvalues under the action of the orbifold group. Taking $`\text{Z}\text{Z}_N`$ for simplicity, one is lead to define the characters
$$𝒳_{gh}=\frac{1}{N}\underset{k=0}{\overset{N1}{}}\rho _{gk}\omega _h^k$$
(7)
with $`g,h=0,1,\mathrm{}N1`$ labelling all possible twists in the $`\sigma `$ and $`\tau `$ directions and $`\omega _h=e^{\frac{2\pi ih}{N}}`$. In the unwisted sectors and in sectors with fixed tori, there are infinitely many characters depending on the choices of allowed momenta and windings. In twisted sectors without fixed tori, a finite number of twisted characters lives at each fixed point. The chiral amplitudes $`\rho _{gh}`$ are defined by the traces
$$\rho _{gh}=\mathrm{Tr}_ghq^{L_0c/24}.$$
(8)
in the $`g`$-twisted sector. Explicit expressions for these traces in $`D=4,6`$ can be found in the appendix A of .
By suitably tuning the parameters, any toroidal or orbifold compactification can be described in terms of a rational SCFT. Moreover in both the rational and the irrational contexts the number of massless characters, i.e. sectors with massless ground-states, is finite. Since only these states will be relevant to our analysis the two cases can be treated in parallel. Focussing on the odd spin-structure, where anomalies potentially reside, one is effectively dealing with the “topological” part of the theory that is largely independent of the moduli of the SCFT. For instance, only massless sectors with non-vanishing Witten index enter the computation of the Euler number of the “compactification manifold” $`M`$ ,
$$\chi (M)=\underset{i,j}{}𝒯_{ij}_i_j,$$
(9)
and the other curvature invariants that allow one to identify the topological class the vacuum configuration belongs to. On the other hand, massive string states do not contribute to the Witten index since they always come in Bose-Fermi degenerate pairs. Ground-states that can become massless at special points (e.g. decompactification limit) of the moduli space, though not contributing to the naïve Witten index, enter the computation of anomalous amplitudes in a subtle way.
Let us now consider the unoriented descendant of a left-right symmetric closed-string theory, the “parent” theory, defined by (5). For definiteness we will consider orientifold reductions by the worldsheet parity operator $`\mathrm{\Omega }`$. More general unoriented descendants that combine $`\mathrm{\Omega }`$ with other internal symmetries can be similarly studied. The unoriented closed string spectrum is determined by halving (5) and adding to it the Klein-bottle projection. Up to overall volume factors and the modular integration ($`\frac{dt}{t}`$) one has
$$𝒦=\frac{1}{2}\mathrm{Tr}_{closed}(\mathrm{\Omega }q^{L_0\frac{c}{24}})=\frac{1}{2t^{\frac{D}{2}}}\underset{i}{}K^i𝒳_i(2it)=\frac{1}{2t^{\frac{D}{2}}}\underset{i}{}K^i_i\mathrm{\Theta }_D(2it)+\mathrm{}.$$
(10)
In the right hand side, for later convenience, only the contribution of the odd spin-structure $`\mathrm{\Theta }_D`$ has been displayed. The sum in (10) is restricted to the diagonal terms ($`𝒯_{ii}0`$) in (5). For permutation modular invariants ($`𝒯_{ij}=0,1`$) the Klein-bottle coefficients $`K^i`$, implementing the action of $`\mathrm{\Omega }`$ on $`_i\overline{}_i`$, are related to the ones in the torus partition function through $`K^i=\pm 𝒯_{ii}`$.
Whenever the action of $`\mathrm{\Omega }`$ is not free the consistency of the theory requires the inclusion of boundaries and the corresponding open string sectors ending on them. The open string partition function consists in the Annulus and Möbius-strip amplitudes. Neglecting the overall volume factors and the modular integration ($`\frac{dt}{t}`$) one has
$`𝒜`$ $`=`$ $`{\displaystyle \frac{1}{2t^{\frac{D}{2}}}}{\displaystyle \underset{i,a,b}{}}A_{ab}^in^an^b𝒳_i\left({\displaystyle \frac{it}{2}}\right)={\displaystyle \frac{1}{2t^{\frac{D}{2}}}}{\displaystyle \underset{i,a,b}{}}A_{ab}^in^an^b_i\mathrm{\Theta }_D\left({\displaystyle \frac{it}{2}}\right)+\mathrm{}.`$
$``$ $`=`$ $`{\displaystyle \frac{1}{2t^{\frac{D}{2}}}}{\displaystyle \underset{i,a}{}}M_a^in^a\widehat{𝒳}_i\left({\displaystyle \frac{it}{2}}+{\displaystyle \frac{1}{2}}\right)={\displaystyle \frac{1}{2t^{\frac{D}{2}}}}{\displaystyle \underset{i,a}{}}_iM_a^in^a\widehat{\mathrm{\Theta }}_D\left({\displaystyle \frac{it}{2}}+{\displaystyle \frac{1}{2}}\right)`$ (11)
A proper basis of real hatted characters $`\widehat{𝒳}_i=T^{1/2}𝒳_i`$ has been introduced in $``$ . The indices $`a,b`$ run over the number of boundaries, i.e. independent Chan-Paton charges, with integer multiplicities $`n^a`$. Barring some exceptions to be discussed later, the number of independent charges one should introduce equals the number of characters $`𝒳_i`$ that are paired with their charge conjugates $`\overline{𝒳}_{i^C}`$ in the parent theory (5). The integers $`A_{ab}^i`$ and $`M_a^i`$ count the number of times the sector $`i`$ runs in an Annulus and Möbius-strip loop, respectively, of open strings with ends on $`(a,b)`$ and $`(a,a)`$ respectively. Here and in the following, we will not distinguish between “complex” (unitary) and “real” (orthogonal or symplectic) Chan-Paton multiplicities $`n^a`$. The sum over $`a`$ will include two contributions for the former and only one for the latter, thus reproducing not only the correct dimensions of the representations but also the correct orientation of the boundary.
The quantities appearing in (10), (11) are highly constrained by two consistency requirements.
The first is the requirement of a consistent interpretation of the transverse channel of these amplitudes as the tree-level exchange of closed string states between boundaries and/or crosscaps. This implies that under $`t1/t`$ the coefficients of the above three amplitudes should reconstruct perfect squares
$`\stackrel{~}{𝒦}`$ $`=`$ $`{\displaystyle \frac{2^{\frac{D}{2}}}{2}}{\displaystyle \underset{i}{}}(\mathrm{\Gamma }^i)^2𝒳_i(q)={\displaystyle \frac{2^{\frac{D}{2}}}{2}}{\displaystyle \underset{i}{}}(\mathrm{\Gamma }_i)^2_i\mathrm{\Theta }_D(q)+\mathrm{}`$
$`\stackrel{~}{𝒜}`$ $`=`$ $`{\displaystyle \frac{2^{\frac{D}{2}}}{2}}{\displaystyle \underset{i,a}{}}(B_a^in^a)^2𝒳_i(q)={\displaystyle \frac{2^{\frac{D}{2}}}{2}}{\displaystyle \underset{i,a}{}}(B_a^in^a)^2_i\mathrm{\Theta }_D(q)+\mathrm{}`$
$`\stackrel{~}{}`$ $`=`$ $`{\displaystyle \frac{2}{2}}{\displaystyle \underset{i,a}{}}(\mathrm{\Gamma }^iB_a^in^a)\widehat{𝒳}_i(q)={\displaystyle \frac{2}{2}}{\displaystyle \underset{i,a}{}}(\mathrm{\Gamma }^iB_a^in^a)_i\widehat{\mathrm{\Theta }}_D(q)+\mathrm{}`$ (12)
The relative powers of 2 result from the different rescalings of the modular parameters ($`\tau _𝒦=2it,\tau _𝒜=it/2,\tau _{}=it/2+1/2`$) that naturally enter the definition of the amplitudes in the direct channel. These rescalings are necessary in order for the amplitudes in the transverse channel to be expressed in terms of the common length of the tube $`\mathrm{}=\frac{1}{2\pi }\mathrm{log}q`$. The coefficients $`\mathrm{\Gamma }^i`$ and $`B_a^i`$ should then be interpreted as the reflection coefficient (one-point function) on a crosscap and on a boundary of type $`a`$ respectively. They are related to the integer coefficients in the direct channel by suitable modular transformations:
$`K^i`$ $`=`$ $`{\displaystyle \underset{j}{}}S^i{}_{j}{}^{}\mathrm{\Gamma }_{}^{j}\mathrm{\Gamma }^j`$
$`A_{ab}^i`$ $`=`$ $`{\displaystyle \underset{j}{}}S^i{}_{j}{}^{}B_{a}^{j}B_b^j`$
$`M_a^i`$ $`=`$ $`{\displaystyle \underset{j}{}}P^i{}_{j}{}^{}\mathrm{\Gamma }_{}^{j}B_a^j`$ (13)
with $`S`$ defined in (6) and $`PT^{1/2}ST^2ST^{1/2}`$.
In addition one should require that fluxes of massless RR-fields should not be trapped inside the “compactification” manifold. This leads to the tadpole cancellation conditions:
$$2^{D/2}\mathrm{\Gamma }^i+B_a^in^a=0$$
(14)
where $`i`$ runs over any sector that flow in the transverse channel and contain RR massless states. Understanding how to precisely relate these conditions to another consistency requirement, the absence of irreducible anomalies in the low-energy field theory, will be the main subject of our investigation.
Finding a solution to (13) is in general a highly non-trivial step in the construction of an open descendant. We will always assume that this has been properly done. It is worth stressing that a “canonical” solution always exists if the torus modular invariant is given by the charge conjugation matrix $`𝒯_{ij}=C_{ij}`$ (“geometric compactifications” such as tori, symmetric orbifolds<sup>2</sup><sup>2</sup>2As we will see, in some orbifold models this may require some identifications among the Chan-Paton charges. and Gepner models belong in this category). In this case the number of independent Chan-Paton charges is exactly equal to the number of characters and one can use the same kind of indices, say $`i,j,\mathrm{}`$, to label both. The ansatz amounts to taking
$`K^i`$ $`=`$ $`Y_{00}^i`$
$`A_{ij}^k`$ $`=`$ $`N_{ij}^k`$
$`M_j^i`$ $`=`$ $`Y_{0j}{}_{}{}^{i},`$ (15)
where
$$N_{ij}{}_{}{}^{k}=\underset{l}{}\frac{S_{il}S_{jl}(S^{})^{lk}}{S_{0l}}$$
(16)
are the fusion rule coefficients, that compute how many independent couplings of the primary fields $`i`$ and $`j`$ give the primary field $`k`$, and
$$Y_{ij}{}_{}{}^{k}=\underset{l}{}\frac{S_{il}P_{jl}(P^{})^{lk}}{S_{0l}}$$
(17)
are (not necessarily positive) integers that satisfy $`Y_{ij}{}_{}{}^{k}=N_{ij}{}_{}{}^{k}(mod\mathrm{\hspace{0.25em}2})`$. The boundary and crosscap reflection coefficients then read
$`\mathrm{\Gamma }^i`$ $`=`$ $`{\displaystyle \frac{P^i_0}{\sqrt{S_{0i}}}}`$
$`B^i_j`$ $`=`$ $`{\displaystyle \frac{S^i_j}{\sqrt{S_{0i}}}}.`$ (18)
Notice that quantities in the transverse amplitudes (12) are related to the spectra (10), (11) through the $`S`$ and $`P`$ modular transformations, whose details vary from model to model. Fortunately this is not the case for the contribution of the odd-spin structure in (3). Each sector of the internal SCFT only enters through its Witten index $`_i`$ that is modular invariant and the modular transformations of $`\mathrm{\Theta }_D`$ are simply
$`\mathrm{\Theta }_D(1/it)`$ $`=`$ $`(it)^{\frac{D2}{2}}\mathrm{\Theta }_D(it)`$
$`\mathrm{\Theta }_D(it+1)`$ $`=`$ $`\mathrm{\Theta }_D(it).`$ (19)
This elementary observation will be the basis of all our manipulations in what follows.
## 3 Anomalies in open string descendants
In this section we discuss the relation between RR tadpoles and irreducible anomalies in a generic even-dimensional open descendant. We first review (and slightly adapt) the results of , which allow us to reproduce the anomaly polynomial by an string computation in the unoriented descendant. Our strategy will be to use modular invariance of the odd-spin structure (19) to parametrize the anomaly polynomial in terms of the reflection coefficients in the transverse channel. On the one hand, this allows us to identify the precise combination of RR tadpoles that corresponds to the irreducible gravitational anomaly. On the other hand, using the completeness properties of the boundary reflection coefficients, this will enable us to map RR tadpoles in sectors with non-vanishing Witten index to irreducible gauge anomalies. After imposing the tadpole/anomaly conditions, the anomaly polynomial turns out to be reducible. It thus admits a simple factorization that suggests the participation of various RR antisymmetric tensors of different rank to a generalized mechanism of anomaly cancellation.
Anomalies are encoded in the odd-spin structure part of the one-loop string amplitudes
$$\underset{f=1}{\overset{N}{}}V_F(p_f,\xi _f)\underset{g=1}{\overset{M}{}}V_G(p_g,h_g)$$
(20)
involving $`N`$ gauge field $`V_F`$ and $`M`$ graviton $`V_G`$ vertex operators with polarization and momenta $`\xi _f,p_f`$ and $`h_g,p_g`$, respectively. The number of external legs is such that $`N+M=D/2+1`$, where $`D`$ is the (even) number of non-compact dimensions. One of the vertices has to be taken with a longitudinal polarization. Since the parent theory is assumed to be anomaly free <sup>3</sup><sup>3</sup>3Modular invariance ensures that this is always the case. anomalous contributions to the above amplitudes only arise from the Klein bottle, Annulus and Möbius strip. In the odd spin-structure there is one supermodulus (zero-mode of the spin 3/2 bosonic superghost $`\beta `$) and one conformal Killing spinor (zero-mode of the spin -1/2 bosonic superghost $`\gamma `$). In order to dispose of the former a picture changing operator is to be inserted. In order to dispose of the latter, one of the vertices (let’s say the one with longitudinal polarization) should be taken in the $`(1)`$ picture. As a result the total superghost charge remains zero as expected for surfaces with vanishing Euler characteristic. After simple manipulations (see for details) a generating function for the anomalous amplitudes can be written as
$$A=_0^{\mathrm{}}𝑑t\frac{d}{dt}e^{S_0+S_F+S_R}_{odd}$$
(21)
where $`S_0`$ is the free action, and the exponentiated effective vertex operators $`S_F,S_R`$ are given by
$`S_F`$ $`=`$ $`{\displaystyle 𝑑sF^a\lambda ^a}`$
$`S_G`$ $`=`$ $`{\displaystyle d^2zR_{\mu \nu }[X^\mu (+\overline{})X^\nu +(\psi ^\mu \stackrel{~}{\psi }^\mu )(\psi ^\nu \stackrel{~}{\psi }^\nu )]}`$ (22)
with
$$F^a=\frac{1}{2}F_{\mu \nu }^a(\psi _0^\mu \psi _0^\nu )R_{\mu \nu }=\frac{1}{2}R_{\mu \nu \rho \sigma }(\psi _0^\rho \psi _0^\sigma ).$$
(23)
bilinears in the zero modes of non-compact fermionic coordinates. After integration over the grassmanian variables, $`F^a`$ and $`R_{\mu \nu }`$ behave as two-forms.
The above determinant has been computed in (see also for additional insights). The crucial point is that the determinant, being defined by a trace in the odd-spin structure where worldsheet bosons and fermions share the same boundary conditions, is given by a $`t`$-independent topological invariant. Each sector of the internal SCFT only enters through its Witten index $`_i`$ and the final result can be written as
$`𝒦_{odd}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}_iK^iI_A(R)`$
$`𝒜_{odd}`$ $`=`$ $`={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a,b}{}}_iA_{ab}^ich_{n^a}(F)ch_{n^b}(F)I_{1/2}(R)`$
$`_{odd}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a}{}}_iM_a^in^ach_{n^a}(2F)I_{1/2}(R)`$ (24)
where $`ch(F)`$ is the Chern character and $`I_A(R)`$ and $`I_{1/2}(R)`$ represent the contributions to the gravitational anomaly of a self-dual antisymmetric tensor and a complex spin 1/2 L-fermion respectively. The additional factor of one-half in the Annulus and Möbius-strip amplitudes reflects the fact that they are counting real fermions. The relation between spin and statistics is responsabile for the extra minus sign of the Klein-bottle contribution with respect to the Annulus and Möbius strip. The former can only contribute loops of (anti)self-dual antisymmetric tensors while the latter can only contribute fermionic loops. In terms of the Hirzebruch polynomial $`\widehat{}(R)`$ and A-roof genus $`\widehat{𝒜}(R)`$ one has
$`I_A(R)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\widehat{}(R)={\displaystyle \frac{1}{8}}+{\displaystyle \frac{1}{48}}R^2+{\displaystyle \frac{1}{32}}\left({\displaystyle \frac{7}{45}}R^4{\displaystyle \frac{1}{9}}(R^2)^2\right)`$ (25)
$`+{\displaystyle \frac{1}{128}}\left({\displaystyle \frac{496}{2835}}R^6{\displaystyle \frac{588}{2835}}R^2R^4+{\displaystyle \frac{140}{2835}}(R^2)^3\right)+\mathrm{}`$
$`I_{1/2}(R)`$ $`=`$ $`\widehat{𝒜}(R)=1+{\displaystyle \frac{1}{48}}R^2+{\displaystyle \frac{1}{32}}\left({\displaystyle \frac{1}{180}}R^4+{\displaystyle \frac{1}{72}}(R^2)^2\right)`$ (26)
$`+{\displaystyle \frac{1}{128}}\left({\displaystyle \frac{1}{2835}}R^6+{\displaystyle \frac{1}{1080}}R^2R^4+{\displaystyle \frac{1}{1296}}(R^2)^3\right)+\mathrm{}`$
$`ch_{n^a}(F)`$ $`=`$ $`tr_{n^a}(\mathrm{exp}iF){\displaystyle \underset{k}{}}{\displaystyle \frac{1}{k!}}F_a^k`$ (27)
where by $`R^{2m}`$ we mean $`\mathrm{tr}_VR^{2m}`$ and wedge products are always understood. In the last line we have introduced the shorthand notation $`F_a^k=tr_{n^a}F^k`$ for later convenience.
The absence of anomalies in the parent theory allows one to rewrite the Klein-bottle contribution in terms of the spectrum of massless closed-string states of the unoriented descendant. As in , we find the expected field-theory result
$$𝒦=\frac{1}{2}\underset{i}{}_iK_iI_A(R)=(n_A^Ln_A^R)I_A(R)+(n_{3/2}^Ln_{3/2}^R)I_{3/2}(R)+(n_{1/2}^Sn_{1/2}^C)I_{1/2}(R),$$
(28)
where $`n_A^{L,R}`$, $`n_{1/2}^{L,R}`$ and $`n_{3/2}^{L,R}`$ are respectively the numbers of antisymmetric tensors, spin $`1/2`$ and spin $`3/2`$ massless closed-string states with definite chirality properties and and $`I_{3/2}(R)=(ch_V(R)1)\widehat{𝒜}(R)`$. Similarly the contribution to the anomaly polynomial coming from the Annulus and Möbius-strip amplitudes (24) reproduces the expected field theory result. Piecing everything together yields
$$𝒫(R,F)=\underset{i}{}_i\left[\frac{1}{2}K^iI_A(R)+\frac{1}{4}\left(\underset{a,b}{}A_{ab}^ich_{n^a}(F)ch_{n^b}(F)+\underset{a}{}M_a^ich_{n^a}(2F)\right)I_{\frac{1}{2}}(R)\right].$$
(29)
In (29) the trace in the (anti-)symmetric representation of the gauge group $`a`$ is written in terms of the trace in the $`n_a`$-dimensional fundamental representation by means of
$$ch_{\frac{1}{2}n^a(n^a\pm 1)}(F)=\frac{1}{2}\left[ch_{n^a}(F)^2\pm ch_{n^a}(2F)\right].$$
(30)
Few remarks are in order. First, modular invariance of the parent closed string theory not only implies the vanishing of the irreducible anomalies but also of reducible ones, i.e. no GS mechanism is at work in the parent theory. This is well-known for left-right symmetric compactifications but we suspect it to be true in any perturbative vacuum configuration. Second, the Klein bottle only contributes to the pure gravitational anomaly. We have not bother to turn on the field-strengths of any closed-string vector boson since all massless fermions of a left-right symmetric closed-string theory are neutral with respect to them. It is well-known that RR vector bosons are not minimally couple to perturbative string states and that non-abelian NS-NS vector bosons may be present in left-right symmetric theories only if not supersymmetric. Finally the CPT theorem in $`D=4k`$ dimensions requires $`n_A^L=n_A^R`$, $`n_{3/2}^L=n_{3/2}^R`$ and $`n_{1/2}^L=n_{1/2}^R`$. As a consequence closed string states do not contribute to the anomaly polynomial in these dimensions. The Klein bottle may give a non-trivial contribution only in $`D=4k+2`$ dimensions. By the same token, the open string sector may contribute to the anomaly in $`D=4k`$ dimensions only when the massless fermions are chiral and belong to complex representations of the Chan-Paton group. Since spinorial representations of orthogonal groups cannot appear perturbatively and symplectic groups only admit (pseudo-)real representations, only when unitary groups are present with associated “complex” Chan-Paton charges can the theory be chiral in $`D=4k`$ dimensions. In $`D=4k+2`$ dimensions, on the contrary, barring few exceptions, any unpaired massless fermion contributes to gauge, gravitational and mixed anomalies.
We are now ready to pin down the UV-IR correspondence between anomalies and tadpoles in open descendants. Using the general relations (13) and the fact that the Witten index, being a $`t`$-independent integer number is invariant under modular transformations, i.e.
$`{\displaystyle \underset{i}{}}_iS^i_j`$ $`=`$ $`_j`$
$`{\displaystyle \underset{i}{}}_iP^i_j`$ $`=`$ $`_j,`$ (31)
one can freely transfer the information encoded in the massless sectors with non-vanishing Witten index from one channel to the other. Indeed, saturing (13) with $`_i`$ and using (31) one can establish the much simpler dictionary
$`{\displaystyle \underset{i}{}}_iK^i`$ $`=`$ $`{\displaystyle \underset{i}{}}_i\mathrm{\Gamma }^i\mathrm{\Gamma }^i`$
$`{\displaystyle \underset{i}{}}_iA_{ab}^i`$ $`=`$ $`{\displaystyle \underset{i}{}}_iB_a^iB_b^i`$
$`{\displaystyle \underset{i}{}}_iM_a^i`$ $`=`$ $`{\displaystyle \underset{i}{}}_iB_a^i\mathrm{\Gamma }^i,`$ (32)
that does not explicitly involve any modular transformation. Plugging these relations into the anomaly polynomial (29) yields
$$𝒫(R,F)=\underset{i}{}_i\left[\frac{1}{2}(\mathrm{\Gamma }^i)^2I_A(R)+\frac{1}{4}\left(\underset{a}{}B_a^ich_{n^a}(F)\right)^2I_{\frac{1}{2}}(R)\frac{1}{4}\mathrm{\Gamma }^i\underset{a}{}B_a^ich_{n^a}(2F)I_{\frac{1}{2}}(R)\right].$$
(33)
The consistency of the effective theory requires the cancellation of the irreducible terms in (33). In particular, for the pure gravitational anomaly, using the relation
$$I_A^D(R)|_{irr}=\frac{2^{D/2}(2^{D/2}1)}{2}I_{1/2}^D(R)|_{irr},$$
(34)
the vanishing of the coefficient of $`tr_VR^{\frac{D}{2}+1}`$ in the expansion of (29) requires
$$\underset{i}{}_i\left(2^{D/2}\mathrm{\Gamma }^i+\underset{a}{}B_a^in^a\right)\left((2^{D/2}1)\mathrm{\Gamma }^i+\underset{a}{}B_a^in^a\right)=0.$$
(35)
Notice that this is a specific linear combination of the RR tadpole conditions (14). In particular the cancellation of the tadpoles automatically implies the cancellation of this irreducible anomaly.
Similarly the absence of irreducible gauge anomaly, i.e. the terms $`tr_{n^b}F^{\frac{D}{2}+1}`$ in the expansion of (29), is equivalent to the conditions <sup>4</sup><sup>4</sup>4Notice that $`U(n)`$ field strengths couple to $`n`$ and $`\overline{n}`$ complex charges, and therefore the corresponding gauge anomalies (both irreducible and reducible) will always include the sum of two terms involving traces in the $`n`$ and $`\overline{n}`$ after using (30).
$$\underset{i}{}_i\left(2^{D/2}\mathrm{\Gamma }^i+\underset{a}{}B_a^in^a\right)B_b^i=0.$$
(36)
As apparent one has in principle one tadpole condition for each factor in the Chan-Paton group. Since the number of independent boundary conditions coincides with the number of independent Chan-Paton multiplicities, one may suspect to be on the right track. Indeed one can turn an index $`a`$, that labels the factors in the Chan-Paton group, into an index $`i`$, that labels the characters that flow in the transverse channel, by making use of the completeness relation,
$$\underset{a}{}B_a^iB_a^j=\frac{\mathrm{\Pi }^{ij}}{\sqrt{S_{0i}S_{0j}}},$$
(37)
valid for any permutation modular invariant parent theory. $`\mathrm{\Pi }^{ij}`$ is the projector onto the set of characters that flow in the transverse channel and $`S_{0i}`$ are the elements in the first row/column of the matrix $`S_{ij}`$ representing modular $`S`$-transformations (6). $`S_{0i}`$ are always positive and are sometimes called “quantum dimensions”.
By saturating the Chan-Paton index $`b`$ in (36) with $`B_b^j`$ we are left with
$$_i\left(2^{D/2}\mathrm{\Gamma }^i+B_a^in^a\right)\mathrm{\Pi }^{ij}=0.$$
(38)
This precisely reproduces all tadpole cancellation conditions for massless R-R closed string states flowing in the transverse channel (as the presence of $`\mathrm{\Pi }^{ij}`$ indicates) and belonging to sectors with non-vanishing Witten index $`_i0`$. It is amusing to observe that in a consistent open descendant in $`D=4k+2`$ the absence of irreducible gauge anomalies always implies the absence of irreducible gravitational anomalies. The vanishing of the first factor in (14) is selected as the relevant tadpole condition.
A comment is in order. As the attentive reader might have observed, a rephrasing of the above conclusion is necessary in the context of orbifolds. Indeed, in these constructions the worldsheet consistency (13) between the direct (one-loop) and transverse (tree) channel often restricts the number of allowed CP charges. As a result the index $`a`$ often runs over a smaller subset than the set of characters flowing in the transverse channel. If this is the case, the index $`i`$ in the completeness conditions (37) and tadpoles (35,36) should be understood as running over an smaller subset defined by the linear combination of characters flowing in this channel. In the new basis (usually conveniently written in terms of the chiral amplitudes) additional “sectors” with vanishing Witten index arise and the corresponding tadpole information is consequently lost in (36).
Modulo this subtlety, we conclude that any massless RR tadpole arising from a sector with non-vanishing Witten index can be identified with some irreducible anomaly. In particular string theory defines a notion of irreducible gauge anomalies even for non-semisimple groups as the ones originating from amplitudes in which all the vertex operators are inserted on the same boundary. Let us now concentrate on the left over reducible part of the anomaly polynomial (33).
## 4 Anomaly polynomial and WZ couplings
Once the tadpole conditions are imposed the anomaly polynomial $`𝒫(R,F)`$ is reducible and can be factorized into a sum of products. This factorization allows one to extract the WZ and anomalous couplings of the R-R massless states that participate in a generalized GSS mechanism of anomaly cancellation <sup>5</sup><sup>5</sup>5Explicit examples of WZ lagrangians in supersymmetric $`D=4,6`$ open string compactifications have been worked it out recently in . Since the factorization is somewhat sensitive to the (even) dimension of the non-compact part of the target space it is convenient to decompose the total anomaly polynomial in terms of its $`D`$-dimensional components
$$𝒫(R,F)=\underset{D}{}𝒫_D(R,F),$$
(39)
where $`𝒫_D(R,F)`$, the anomaly polynomial in $`D`$-dimension, is a $`(D+2)`$-form.
After some simple algebra one finds that the factorization is almost always unique. Only in $`D=8`$ and $`D=10`$ we find mild ambiguities. The anomaly polynomials in $`D`$ dimensions explicitly factorize as
$`𝒫_0(R,F)`$ $`=`$ $`0`$
$`𝒫_2(R,F)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a,b}{}}_iB_a^iB_b^i\left[X_2^aX_2^b\right]`$
$`𝒫_4(R,F)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a,b}{}}_iB_a^iB_b^i\left[2X_2^aX_4^b\right]`$
$`𝒫_6(R,F)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a,b}{}}_iB_a^iB_b^i\left[X_4^aX_4^b+2X_2^aX_6^b(1)\right]`$
$`𝒫_8(R,F)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a,b}{}}_iB_a^iB_b^i\left[2X_4^aX_6^b(\xi _b)+2X_2^aX_8^b(\xi _b)\right]`$
$`𝒫_{10}(R,F)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i,a,b}{}}_iB_a^iB_b^i\left[2X_2^aX_{10}^b(\xi _b)+2X_4^aX_8^b(1)+X_6^a(\xi _a)X_6^b(\xi _b)\right],`$ (40)
where, in order to keep the above formulas as compact as possible, we have introduced the definitions
$`X_2^a`$ $`=`$ $`iF_a`$
$`X_4^a`$ $`=`$ $`{\displaystyle \frac{i^2}{2!}}(F_a^2{\displaystyle \frac{n^a}{32}}R^2)`$
$`X_6^a(\xi _a)`$ $`=`$ $`{\displaystyle \frac{i^3}{3!}}(F_a^3{\displaystyle \frac{\xi _a}{16}}F_aR^2)`$
$`X_8^a(\xi _a)`$ $`=`$ $`{\displaystyle \frac{i^4}{4!}}(F_a^4+{\displaystyle \frac{\xi _a2}{8}}F_a^2R^2+{\displaystyle \frac{54\xi _a}{1024}}n^a(R^2)^2+{\displaystyle \frac{1}{256}}n^aR^4)`$
$`X_{10}^a(\xi _a)`$ $`=`$ $`{\displaystyle \frac{i^5}{5!}}(F_a^5+{\displaystyle \frac{5\xi _a10}{24}}F_a^3R^2+{\displaystyle \frac{105\xi _a^2}{768}}F_a(R^2)^2+{\displaystyle \frac{1}{96}}F_aR^4).`$ (41)
As mentioned above, the only ambiguities we find are parametrized by real constants $`\xi _a`$ that enter the definitions of the factors $`X_6^a`$, $`X_8^a`$ and $`X_{10}^a`$ in the anomaly polynomials for $`D=8,10`$. Previous analyses correspond to the choice $`\xi _a=1`$. The parameters $`\xi _a`$ could be genuine free parameters in the effective lagrangian or an artifact of the procedure. In order to settle this issue one should compute additional string amplitudes on the disk. It is remarkable that these ambiguities are absent in the supersymmetric case and whenever $`tr_{n^a}F=0`$ for any $`a`$, since all $`\xi _a`$ would disappear from the above formulas.
The nice feature of (41) is that, one can easily recognize the closed string RR fields participating in the GSS mechanism of anomay cancellation . Indeed the index $`i`$ labels the character<sup>6</sup><sup>6</sup>6More precisely the index $`i`$ refers to the pair $`(i,i^C)`$ of characters that are paired in the torus partition function. to which the corresponding closed string field belongs while $`B_a^i`$ measures the strength of the coupling of a “generalized brane” of type $`a`$ to these massless R-R states and their massive (super-)partners. Each sector includes $`p+1`$-form potentials of various degrees. The generalized GSS mechanism requires that a $`p+1`$-form potential $`C_{p+1}^i`$ in the sector labelled by $`i`$ couple with strength $`B_a^i`$ to a $`Dp1`$-form $`X_{Dp1}^a`$ of type $`a`$ in the factorized anomaly polynomial. We can then extract from (41) all WZ couplings responsabile for this GSS mechanism in the effective $`D`$-dimensional theory
$`L_0`$ $`=`$ $`0`$ (42)
$`L_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,a}{}}B_a^iC_0^iX_2^a`$ (43)
$`L_4`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,a}{}}B_a^i\left(C_0^iX_4^a+C_2^iX_2^a\right)`$ (44)
$`L_6`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,a}{}}B_a^i\left(C_0^iX_6^a(1)+C_2^iX_4^a+C_4^iX_2^a\right)`$ (45)
$`L_8`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,a}{}}B_a^i\left(C_0^iX_8^a(\xi )+C_2^iX_6^a(\xi )+C_4^iX_4^a+C_6^iX_2^a\right)`$ (46)
$`L_{10}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,a}{}}B_a^i\left(C_0^iX_{10}^a(\xi )+C_2^iX_8^a(1)+C_4^iX_6^a(\xi )+C_6^iX_4^a+C_8^iX_2^a\right)`$ (47)
The factorization of the transverse channel amplitudes in the odd spin structure gives rise to the “odd” propagator which turns a $`(p+1)`$-form into its dual $`(Dp3)`$-form. More precisely
$$dC_{p+1}^idC_{Dp3}^j=_i\mathrm{\Pi }^{ij}𝒱_D.$$
(48)
with $`𝒱_D`$ a regularized D-dimensional volume. Notice that one can simultaneously rescale any $`C_{p+1}^i`$ by a factor $`\alpha _{p+1}^i`$ and its dual $`C_{Dp1}^i`$ by a factor $`1/\alpha _{p+1}^i`$ without modifying the factorization of the anomaly polynomial. The couplings (47) will of course absorb these renormalizations.
From the expression of (47) one can easily recognize the anomalous $`U(1)`$’s. They are to be identified as the abelian factors of type $`a`$ entering in $`X_2^a`$ and coupling anomalously to $`C_{D2}^i`$ with a non-vanishing reflection coefficient $`B_a^i`$.
Although we have concentrated our attention on the CP-odd part of the effective lagrangian, in supersymmetric models some interesting CP-even coupling are related to (47) by supersymmetry. WZ lagrangians in supersymmetric $`D=4,6`$ open string models have been worked it out in . The terms displyed in (47) cannot be the whole story. RR fields belonging to sectors with vanishing Witten index are not present. In particular the WZ couplings RR axion belonging to the identity sector in any supersymmetric vacuum configurations in $`D=4`$ are not reproduced by (47). In order to get around this problem one has to find a way to disclose the anomalous content encoded in RR tadpoles of sectors with vanishing Witten index.
## 5 Tadpoles with $`=0`$ and anomalous amplitudes
In section 3 we have shown that RR tadpoles in sectors with non-vanishing Witten index are in one-to-one correspondence with irreducible anomalies. Unfortunately the vanishing of the Witten index $`_i`$ in a given internal sector makes the associated anomaly condition (36) empty and the corresponding tadpole information lost. In this section we argue that a careful look into the anomaly structure of the effective theory may overcome this apparent asymmetry. This will be done at the price of considering anomalous amplitudes involving non-trivial insertions in the internal SCFT. The correct choice is given by the minimal number of insertions that removes the Bose-Fermi degeneracy in a given sector with $`=0`$. As a result one can once again translate the associated tadpole condition into a condition for anomaly cancellation through (36).
Rather than being general let us illustrate this point in the simplest context of even-dimensional ($`d=2\mathrm{}=10D`$) toroidal compactifications of the $`SO(32)`$ type I theory. To start with, we will neither turn on Wilson lines nor expectation values of the $`B_{NSNS}`$-field (equivalent to non-commuting Wilson lines) . We will include their effects in the analysis later on. In this simple situation we have a single massless character $`𝒳_0`$. Its Witten index is clearly zero due to the $`2\mathrm{}`$ zero-modes $`\mathrm{\Psi }_0^I`$ of the internal fermionic coordinates ($`I=1\mathrm{}2\mathrm{}`$). In order to get a non-trivial contribution from this sector one should compute anomalous amplitudes involving insertions in the internal SCFT such that all fermionic zero-modes $`\mathrm{\Psi }_0^I`$ are soaked up. The contribution is now measured by non-vanishing generalized Witten indices $`\widehat{}_0=TrF^{\mathrm{}}()^F`$ <sup>7</sup><sup>7</sup>7The “quasi” topological index $`=TrF()^F`$ was introduced in the $`N=2`$ context in . Irrespective of the dimension, anomalous amplitudes descending from the $`D=10`$ exagon meet the requirements. The string computation is clearly the same as in ten dimensions but the interpretation in terms of the lower dimensional field theory is different. With respect to the canonical anomaly, the anomalous amplitudes under consideration are associated to higher derivative CP-odd terms involving scalars and KK vectors. These receive anomalous contributions not only from loops of massless fermions and but also from their higher KK-modes. The presence of such anomalous terms reflects the presence of a non-trivial coupling of the universal RR 2-form that compensates for the anomalous variation of the fermion determinants only when $`N=32`$ as in $`D=10`$, i.e. only when the RR tadpole of the identity sector with $`_0=0`$, is enforced. If one turns on continuous Wilson lines, thus generically breaking $`SO(32)`$ to $`U(1)^{16}`$, one is effectively trading some anomalous contribution of the massless fermions to the CP-odd thresholds with an additional contribution of the massive BPS states. If one turns on non-commuting Wilson lines, equivalent to a quantized background for the NS-NS antisymmetric tensor $`B_{NSNS}`$, one is reducing both the rank of the Chan-Paton group as well as the number of gaugini and the coupling of the universal RR 2-form. Notice that this has a neat counterpart in the dual heterotic string. In this case, modular invariance replaces the RR tadpole condition in ensuring that any anomalous variation of the one-loop effective action is zero. When a gauge symmetry $`G`$ in the heterotic string is realized in terms of a world-sheet current algebra at level $`\kappa `$ the effective GS coupling is rescaled by that same factor. This is precisely what happens in the type I settings we have just discussed if one identifies the level of the worldsheet current algebra $`\kappa `$ with $`2^{r/2}`$ where $`r`$ is the rank of $`B_{NSNS}`$.
In non-trivial open string configurations involving sectors with vanishing Witten index $`_i=0`$, other than the identity sector, one can apply similar considerations. The vanishing of $`_i`$ signals the presence of a certain number $`d`$ of unsoaked fermionic zero-modes together with their bosonic superpartner under worldsheet supersymmetry. This essentially free piece of the internal SCFT is more like the $`T^d`$-toroidal compactification we have discussed so far. Indeed one can trace the different higher dimensional origins of the relevant anomalous amplitudes by their different scalings with the internal “volumes”. By volume here we do not neccesarily mean a geometrical one but simply the result of the integration over the compact bosonic zero modes $`X_0^I`$ ($`I=1,\mathrm{}d`$). In the context of geometric orbifolds, different amplitudes<sup>8</sup><sup>8</sup>8In our approach, amplitudes are linear combinations of characters. The former follow from the latter by inverting (7). are in general associated to different fixed tori. The different scalings of the anomalous amplitudes allow one to identify the constraints on the different sets of branes and the corresponding RR tadpoles in sectors with $`_i=0`$.
It is instructive to explore the implications of similar anomalous diagrams in different string settings. In vacuum configurations for the heterotic string, moduli fields such as the radii of toroidal compactifications and the blowing up modes of an orbifold specify non-linear $`\sigma `$-models with continuous non-compact symmetries. These continuous symmetries combine geometric transformations such as internal diffeomorphisms with non-geometric symmetries that arise because of the extended nature of the fundamental constituents. The related presence of winding states breaks the continuous symmetries down to discrete ones (“T-duality”). Some of these discrete symmetries of the tree-level effective lagrangian are broken by quantum effects. On the one hand, they act by chiral transformations on the fermions that transform with moduli-dependent phases. The discrete charges (modular weights) are very sensitive to the sector each massless fermion belongs to. On the other hand, massive string states give rise to moduli-dependent thresholds corrections to CP-odd terms, related to the anomalous amplitudes discussed above, in theories with fixed tori, i.e. internal sectors with $`_i=0`$. In $`𝒩=1`$ supersymmetric compactifications, these CP-odd threshold corrections are anomalous in that they violate an integrability condition . This violation is not renormalized beyond one-loop . The combined variation of the contribution of the massless fermions and the CP-odd thresholds under discrete T-duality transformations is non-zero. There is howevere a universal GS counterterm that cancels the left-over discrete anomaly . Worldsheet modular invariance at one-loop thus guarantees the absence of target space modular anomalies.
After inclusion of all stable bound-states of strings and branes “T-duality” is expected to be promoted to “U-duality”. In theories with open and unoriented strings, T-duality is broken at the perturbative level by the presence of the background brane configuration. Still, by including the contribution of branes wrapping around cycles of the internal manifold one expects to recover a discrete T-duality symmetry isomorphic to a T-duality subgroup of the U-duality of the parent theory. Absence of T-duality anomalies would then be a consequence of the vanishing of RR-tadpoles, including sectors with $`=0`$, much in the same way as absence of T-duality anomalies is a consequence of modular invariance in closed-string theories.
Following the same line of reasoning, one may argue that the absence of discrete anomalies in non-supersymmetric vacua requires the absence of all massless RR tadpoles. Although very little is known about the effective lagrangian in this case we expect the anomalous terms including the anomalous CP-odd thresholds to be “topological” and as such to be non-renormalized beyond one loop except possibly by world-sheet instantons or D-istantons.
Let us conclude this section with a side remark on the odd-dimensional case. Although there are no local anomalies for odd dimensional manifolds without boundaries and/or branes, global anomalies cannot be a priori excluded . It is reassuring to observe that these are absent for toroidal compactifications with a Chan-Paton group of even rank as required by RR tadpole cancellation, i.e. $`N=32\times 2^{r/2}`$ or, for open descendants without open strings, $`N=0`$.
## 6 Discussion and concluding remarks
Most of the vacuum configurations with open and unoriented strings are not simply geometric compactifications of ten-dimensional theories with open and unoriented superstrings. In a non-trivial background, genuine lower dimensional open string excitations can be located at new kinds of branes that might not admit a large volume description. Still one can try to extract some insights from the anomaly polynomial we have derived in some generality. In principle, once the topological data of the “compactification” manifold are identified, one can extract some useful information about the relevant brane configuration and/or the vacuum gauge bundle by comparing the WZ couplings found in the genuine string computation with the ones that would result from a naive compactification. All additional terms are to be related to extra bound-states of branes and/or to non-trivial configurations of the background gauge fields. We have parametrize the relevant information in terms of the boundary reflection coefficients $`B_a^i`$, i.e. one-point functions of massless states of type $`(i,i^C)`$ on the disk with boundary conditions of type $`a`$. These $`B_a^i`$ enjoy some interesting completeness properties and behave as vielbeins, i.e. they allow one to transform an index of type $`i`$ labelling a sector of the closed string spectrum that flows in the transverse channel into an index of type $`a`$ that labels the independent Chan-Paton multiplicities. It is conceivable that the $`B_a^i`$ satisfy some interesting evolution equation that would allow one to compute them starting from special “rational” points in the moduli space of the vacuum configuration.
At these points, classifying all possible boundary conditions compatible with preserving half of the bulk symmetries is expected to be feasible. The worldsheet SCFT description pursued in this paper and the topological nature of the anomaly-related terms should be sufficient to show that the $`B_a^i`$ for massless sectors with non-vanishing Witten index should be independent of the moduli. Indeed, for D-branes wrapped around supersymmetric cycles in Calabi-Yau spaces it is known that $`B_a^i`$ for $`i`$ a massless sector are quasi-topological . For A-type boundary conditions the $`B_a^i`$, with $`i`$ running over the (c,c) primary fields and $`a`$ over the middle cohomology cycles, only depends on the complex structure deformations and may be computed at large volumes. For B-type boundary conditions the $`B_a^i`$, with $`i`$ running over the (c,a) primary fields and $`a`$ over the vertical cohomology cycles, only depends on the Kähler structure deformations and can be computed in terms of the “quantum intersection form”, i.e. the “topological intersection form” (an integer) plus worldsheet instanton corrections. For massive sectors this is no longer expected to be the case, but the polynomial equations satisfied by the $`B_a^i`$ may rescue the situation . It would be interesting to find a way to extend the arguments of based on the topological twist of $`𝒩=2`$ SCFT to worldsheet theories that only enjoy $`𝒩=1`$ superconformal symmetry.
We have not discussed the tadpoles in the NS-NS sector that may be disposed of by means of a Fischler-Susskind mechanism. The latter may destabilize the vacuum configuration. A signal could be found in behaviour of the $`B_a^i`$ at points where a given brane configuration could decay into its constituents or tachyon condensation could trigger the formation of a stable non-BPS bound state.
The applications of the anomaly-tadpole correspondence we have established go far beyond of the contexts exploited in this paper. Indeed, our results only rely on the consistency between the open and closed string channel in the presence of boundaries and crosscaps. Relying on the $`𝒩=(1,1)`$ superconformal symmetry on the worldsheet enables one to isolate the odd spin-structure as the source of all potential anomalies. Modular invariance then allows one to identify anomalous amplitudes with RR tadpoles. Similar analyses can be applied to any (BPS or not) brane configuration satisfying the consistency requirements. Tadpoles are associated to closed string absortion processes encoded in the dynamics of boundaries and crosscaps on the worldsheet, Anomalies enter as one-loop CP-odd effects in the effective theory governing the dynamics of the brane configuration. The UV-IR dictionary relating these two effects is expected to be realized in any SYM/SUGRA correspondence described in terms of boundaries on the string worldsheet.
Acknowledgements
We would like to thank A. Sagnotti and Ya.S.Stanev for valuable discussions, M. Serone for collaboration at early stages of this project and G. Pradisi and G.C. Rossi for useful comments. |
warning/0002/hep-th0002230.html | ar5iv | text | # 1 Introduction and summary of the results
## 1 Introduction and summary of the results
Holography states that a $`(d+1)`$-dimensional gravitational theory (referred to as the bulk theory) should have a description in terms of a $`d`$-dimensional field theory (referred to as the boundary theory) with one degree of freedom per Planck area . The arguments leading to the holographic principle use rather generic properties of gravitational physics, indicating that holography should be a feature of any quantum theory of gravity. Nevertheless it has been proved a difficult task to find examples where holography is realized, let alone to develop a precise dictionary between bulk and boundary physics. The AdS/CFT correspondence provides such a realization with a rather precise computational framework . It is, therefore, desirable to sharpen the existing dictionary between bulk/boundary physics as much as possible. In particular, one of the issues one would like to understand is how spacetime is built holographically out of field theory data.
The prescription of gives a concrete proposal for a holographic computation of physical observables. In particular, the partition function of string theory compactified on AdS spaces with prescribed boundary conditions for the bulk fields is equal to the generating functional of conformal field theory correlation functions, the boundary value of fields being now interpreted as sources for operators of the dual conformal field theory (CFT). String theory on anti-de Sitter (AdS) spaces is still incompletely understood. At low energies, however, the theory becomes a gauged supergravity with an AdS ground state coupled to Kaluza-Klein (KK) modes. On the field theory side, this corresponds to the large $`N`$ and strong ’t Hooft coupling regime of the CFT. So in the AdS/CFT context the question is how one can reconstruct the bulk spacetime out of CFT data. One can also pose the converse question: given a bulk spacetime, what properties of the dual CFT can one read off?
The prescription of equates the on-shell value of the supergravity action with the generating functional of connected graphs of composite operators. Both sides of this correspondence, however, suffer from infinities —infrared divergences on the supergravity side and ultraviolet divergences on the CFT side. Thus, the prescription of should more properly be viewed as an equality between bare quantities. Ones needs to renormalize the theory to obtain a correspondence between finite quantities. It is one of the aims of this paper to present a systematic way of performing such renormalization.
The CFT data<sup>4</sup><sup>4</sup>4We assume that the CFT we are discussing has an AdS dual. Our results only depend on the spacetime dimension and apply to all cases where the AdS/CFT duality is applicable, so we shall not specify any particular CFT model. that we will use are: which operators are turned on, and what is their vacuum expectation value. Since the boundary metric (or, more properly, the boundary conformal structure) couples to the boundary stress-energy tensor, the reconstruction of the bulk metric to leading order involves a detailed knowledge of the way the energy-momentum tensor is encoded holographically. There is by now an extended literature on the study of the stress-energy tensor in the context of the AdS/CFT correspondence starting from . We will build on these and other related works . Our starting point will be the calculation of the infrared divergences of the on-shell gravitational action . Minimally subtracting the divergences by adding counterterms leads straightforwardly to the results in . After the subtractions have been made one can remove the (infrared) regulator and obtain a completely explicit formula for the expectation value of the dual stress-energy tensor in terms of the gravitational solution.
We will mostly concentrate on the gravitational sector, i.e. in the reconstruction of the bulk metric, but we will also discuss the coupling to scalars. Our approach will be to build perturbatively an Einstein manifold of constant negative curvature (which we will sometimes refer to as an asymptotically AdS space) as well as a solution to the scalar field equations on this manifold out of CFT data. The CFT data we start from is what sources are turned on. We will include a source for the dual stress-energy tensor as well as sources for scalar composite operators. This means that in the bulk we need to solve the gravitational equations coupled to scalars given a conformal structure at infinity and appropriate Dirichlet boundary conditions for the scalars. It is well-known that if one considers the standard Euclidean AdS (i.e., with isometry $`SO(1,d+1)`$), the scalar field equation with Dirichlet boundary conditions has a unique solution. In the Lorentzian case, because of the existence of normalizable modes, the solution ceases to be unique. Likewise, the Dirichlet boundary condition problem for (Euclidean) gravity has a unique (up to diffeomorphisms) smooth solution in the case the bulk manifold in topologically a ball and the boundary conformal structure sufficiently close to the standard one . However, given a boundary topology there may be more than one Einstein manifold with this boundary. For example, if the boundary has the topology of $`S^1\times S^{d1}`$, there are two possible bulk manifolds : one which is obtained from standard AdS by global identifications and is topologically $`S^1\times R^d`$, and another, the Schwarzschild-AdS black hole, which is topologically $`R^2\times S^{d1}`$.
We will make no assumption on the global structure of the space or on its signature. The CFT should provide additional data in order to retrieve this information. Indeed, we will see that only the information about the sources leaves undetermined the part of the solution which is sensitive on global issues and/or the signature of spacetime. To determine that part one needs new CFT data. To leading order these are the expectation values of the CFT operators.
In particular, in the case of pure gravity, we find that generically a boundary conformal structure is not sufficient in order to obtain the bulk metric. One needs more CFT data. To leading order one needs to specify the expectation value of the boundary stress-energy tensor. Since the gravitational field equation is a second order differential equation, one may expect that these data are sufficient in order to specify the full solution. In general, however, non-local observables such as Wilson loops may be needed in order to recover global properties of the solution and reconstruct the metric in the deep interior region. Furthermore, higher point functions of the stress-energy tensor may be necessary if higher derivatives corrections such as $`R^2`$ terms are included in the action. We emphasize that we make no assumption about the regularity of the solution. Under additional assumptions the metric may be determined by fewer data. For example, as we mentioned above, under certain assumptions on the topology and the boundary conformal structure one obtains a unique smooth solution . Another example is the case when one restricts oneself to conformally flat bulk metrics. Then a conformally flat boundary metric does yield a unique, up to diffeomorphisms and global identifications, bulk metric .
Turning things around, given a specific solution, we present formulae for the expectation values of the dual CFT operators. In particular, in the case the operator is the stress energy tensor, our formulae have a “dual” meaning : both as the expectation value of the stress-energy tensor of the dual CFT and as the quasi-local stress-energy tensor of Brown and York . We provide very explicit formulae for the stress-energy tensor associated with any solution of Einstein’s equations with negative constant curvature.
Let us summarize these results for spacetime dimension up to six. The first step is to rewrite the solution in the Graham-Fefferman coordinate system
$$ds^2=G_{\mu \nu }dx^\mu dx^\nu =\frac{l^2}{r^2}\left(dr^2+g_{ij}(x,r)dx^idx^j\right)$$
(1.1)
where
$$g(x,r)=g_{(0)}+r^2g_{(2)}+\mathrm{}+r^dg_{(d)}+h_{(d)}r^d\mathrm{log}r^2+𝒪(r^{d+1})$$
(1.2)
The logarithmic term appears only in even dimensions and only even powers of $`r`$ appear up to order $`r^{[(d1)]}`$, where $`[a]`$ indicates the integer part of $`a`$. $`l`$ is a parameter of dimension of length related to the cosmological constant as $`\mathrm{\Lambda }=\frac{d(d1)}{2l^2}`$. Any asymptotically AdS metric can be brought in the form (1.1) near the boundary (, see also ). Once this coordinate system has been reached, the stress-energy tensor reads
$$T_{ij}=\frac{dl^{d1}}{16\pi G_\text{N}}g_{(d)ij}+X_{ij}[g_{(n)}].$$
(1.3)
where $`X_{ij}[g_{(n)}]`$ is a function of $`g_{(n)}`$ with $`n<d`$. Its exact form depends on the spacetime dimension and it reflects the conformal anomalies of the boundary conformal field theory. In odd (boundary) dimensions, where there are no gravitational conformal anomalies, $`X_{ij}`$ is equal to zero. The expression for $`X_{ij}[g_{(n)}]`$ for $`d=2,4,6`$ can be read off from (3.10), (3.15) and (3.16), respectively. The universal part of (1.3) (i.e. with $`X_{ij}`$ omitted) was obtained previously in . Actually, to obtain the dual stress-energy tensor it is sufficient to only know $`g_{(0)}`$ and $`g_{(d)}`$ as $`g_{(n)}`$ with $`n<d`$ are uniquely determined from $`g_{(0)}`$, as we will see. The coefficient $`h_{(d)}`$ of the logarithmic term in the case of even $`d`$ is also directly related to the conformal anomaly: it is proportional to the metric variation of the conformal anomaly.
It was pointed out in that this prescription for calculating the boundary stress-energy tensor provides also a novel, free of divergences<sup>5</sup><sup>5</sup>5 We emphasize, however, that one has to subtract the logarithmic divergences in even dimensions in order for the stress-energy tensor to be finite., way of computing the gravitational quasi-local stress-energy tensor of Brown and York . This approach was recently criticized in , and we take this opportunity to address this criticism. Conformal anomalies reflect infrared divergences in the gravitational sector . Because of these divergences one cannot maintain the full group of isometries even asymptotically. In particular, the isometries of AdS that rescale the radial coordinate (these correspond to dilations in the CFT) are broken by infrared divergences. Because of this fact, bulk solutions that are related by diffeomorphisms that yield a conformal transformation in the boundary do not necessarily have the same mass. Assigning zero mass to the spacetime with boundary $`R^d`$, one obtains that, due to the conformal anomaly, the solution with boundary $`R\times S^{d1}`$ has non-zero mass. This parallels exactly the discussion in field theory. In that case, starting from the CFT on $`R^d`$ with vanishing expectation value of the stress-energy tensor, one obtains the Casimir energy of the CFT on $`R\times S^{d1}`$ by a conformal transformation . The agreement between the gravitational ground state energy and the Casimir energy of the CFT is a direct consequence of the fact that the conformal anomaly computed by weakly coupled gauge theory and by supergravity agree . It should be noted that, as emphasized in , agreement between gravity/field theory for the ground state energy is achieved only after all ambiguities are fixed in the same manner on both sides.
A conformal transformation in the boundary theory is realized in the bulk as a special diffeomorphism that preserves the form of the coordinate system (1.1) . Using these diffeomorphisms one can easily study how the (quantum, i.e., with the effects of the conformal anomaly taken into account) stress-energy tensor transforms under conformal transformations. Our results, when restricted to the cases studied in the literature , are in agreement with them. We note that the present determination is considerably easier than the one in .
The discussion is qualitatively the same when one adds matter to the system. We discuss scalar fields but the discussion generalizes straightforwardly to other kinds of matter. We study both the case the gravitational background is fixed and the case gravity is dynamical.
Let us summarize the results for the case of scalar fields in a fixed gravitational background (given by a metric of the form (1.1)). We look for solutions of massive scalar fields with mass $`m^2=(\mathrm{\Delta }d)\mathrm{\Delta }`$ that near the boundary have the form (in the coordinate system (1.1))
$$\mathrm{\Phi }(x,r)=r^{d\mathrm{\Delta }}\left(\varphi _{(0)}+r^2\varphi _{(2)}+\mathrm{}+r^{2\mathrm{\Delta }d}\varphi _{(2\mathrm{\Delta }d)}+r^{2\mathrm{\Delta }d}\mathrm{log}r^2\psi _{(2\mathrm{\Delta }d)}\right)+𝒪(r^{\mathrm{\Delta }+1}).$$
(1.4)
The logarithmic terms appears only when $`2\mathrm{\Delta }d`$ is an integer and we only consider this case in this paper. We find that $`\varphi _{(n)}`$, with $`n<2\mathrm{\Delta }d`$, and $`\psi _{(2\mathrm{\Delta }d)}`$ are uniquely determined from the scalar field equation. This information is sufficient for a complete determination of the infrared divergences of the on-shell bulk action. In particular, the logarithmic term $`\psi _{(2\mathrm{\Delta }d)}`$ in (1.4) is directly related to matter conformal anomalies. These conformal anomalies were shown not to renormalize in . We indeed find exact agreement with the computation in . Adding counterterms to cancel the infrared divergences we obtain the renormalized on-shell action. We stress that even in the case of a free massive scalar field in a fixed AdS background one needs counterterms in order for the on-shell action to be finite (see (5.1)). The coefficient $`\varphi _{(2\mathrm{\Delta }d)}`$ is left undetermined by the field equations. It is determined, however, by the expectation value of the dual operator. Differentiating the renormalized on-shell action one finds (up to terms contributing contact terms in the 2-point function)
$$O(x)=(2\mathrm{\Delta }d)\varphi _{(2\mathrm{\Delta }d)}(x)$$
(1.5)
This relation, with the precise proportionality coefficient, has first been derived in . The value of the proportionality coefficient is crucial in order to obtain the correct normalization of the 2-point function in standard AdS background .
In the case the gravitational background is dynamical we find that, for scalars that correspond to irrelevant operators, our perturbative treatment is consistent only if one considers single insertions of the irrelevant operator, i.e. the source is treated as an infinitesimal parameter, in agreement with the discussion in . For scalars that correspond to marginal and relevant operators one can compute perturbatively the back-reaction of the scalars to the gravitational background. One can then regularize and renormalize as in the discussion of pure gravity or scalars in a fixed background. For illustrative purposes we analyze a simple example.
This paper is organized as follows. In the next section we discuss the Dirichlet problem for AdS gravity and we obtain an asymptotic solution for a given boundary metric (up to six dimensions). In section 3 we use these solutions to obtain the infrared divergences of the on-shell gravitational action. After renormalizing the on-shell action by adding counterterms, we compute the holographic stress-energy tensor. Section 4 is devoted to the study of the conformal transformation properties of the boundary stress-energy tensor. In section 5 we extend the analysis of sections 2 and 3 to include matter. In appendices A and D we give the explicit form of the solutions discussed in section 2 and section 5. Appendix B contains the explicit form of the counterterms discussed in section 3. Finally, in appendix C we present a proof that the coefficient of the logarithmic term in the metric (present in even boundary dimensions) is proportional to the metric variation of the conformal anomaly.
## 2 Dirichlet boundary problem for AdS gravity
The Einstein-Hilbert action for a theory on a manifold $`M`$ with boundary $`M`$ is given by<sup>6</sup><sup>6</sup>6 Our curvature conventions are as follows $`R_{ijk}{}_{}{}^{l}=_i\mathrm{\Gamma }_{jk}{}_{}{}^{l}+\mathrm{\Gamma }_{ip}{}_{}{}^{l}\mathrm{\Gamma }_{jk}^{}{}_{}{}^{p}ij`$ and $`R_{ij}=R_{ikj}^k`$. We these conventions the curvature of AdS comes out positive, but we will still use the terminology “space of constant negative curvature”. Notice also that we take $`\text{d}^{d+1}x=\text{d}^dx_0^{\mathrm{}}\text{d}r`$ and the boundary is at $`r=0`$ (in the coordinate system (1.1)). The minus sign in front of the trace of the second fundamental form is correlated with the choice of having $`r=0`$ in the lower end of the radial integration.
$$S_{\text{gr}}[G]=\frac{1}{16\pi G_\text{N}}[_M\text{d}^{d+1}x\sqrt{G}(R[G]+2\mathrm{\Lambda })_M\text{d}^dx\sqrt{\gamma }\mathrm{\hspace{0.17em}2}K],$$
(2.1)
where $`K`$ is the trace of the second fundamental form and $`\gamma `$ is the induced metric on the boundary. The boundary term is necessary in order to get an action which only depends on first derivatives of the metric , and it guarantees that the variational problem with Dirichlet boundary conditions is well-defined.
According to the prescription of , the conformal field theory effective action is given by evaluating the on-shell action functional. The field specifying the boundary conditions for the metric is regarded as a source for the boundary operator. We therefore need to obtain solutions to Einstein’s equations,
$$R_{\mu \nu }\frac{1}{2}RG_{\mu \nu }=\mathrm{\Lambda }G_{\mu \nu },$$
(2.2)
subject to appropriate Dirichlet boundary conditions.
Metrics $`G_{\mu \nu }`$ that satisfy (2.2) have a second order pole at infinity. Therefore, they do not induce a metric at infinity. They do induce, however, a conformal class, i.e. a metric up to a conformal transformation. This is achieved by introducing a defining function $`r`$, i.e. a positive function in the interior of $`M`$ that has a single zero and non-vanishing derivative at the boundary. Then one obtains a metric at the boundary by $`g_{(0)}=r^2G|_M`$ <sup>7</sup><sup>7</sup>7Throughout this article the metric $`g_{(0)}`$ is assumed to be non-degenerate. For studies of the AdS/CFT correspondence in cases where $`g_{(0)}`$ is degenerate we refer to .. However, any other defining function $`r^{}=r\mathrm{exp}w`$ is as good. Therefore, the metric $`g_{(0)}`$ is only defined up to a conformal transformation.
We are interested in solving (2.2) given a conformal structure at infinity. This can be achieved by working in the coordinate system (1.1) introduced by Feffermam and Graham . The metric in (1.1) contains only even powers of $`r`$ up to the order we are interested in (see also ). For this reason, it is convenient to use the variable $`\rho =r^2`$ , <sup>8</sup><sup>8</sup>8Greek indices, $`\mu ,\nu ,..`$ are used for $`d+1`$-dimensional indices, Latin ones, $`i,j,..`$ for $`d`$-dimensional ones. To distinguish the curvatures of the various metrics introduced in (2) we will often use the notation $`R_{ij}[g]`$ to indicate that this is the Ricci tensor of the metric $`g`$, etc.
$`ds^2=G_{\mu \nu }dx^\mu dx^\nu =l^2\left({\displaystyle \frac{d\rho ^2}{4\rho ^2}}+{\displaystyle \frac{1}{\rho }}g_{ij}(x,\rho )dx^idx^j\right)`$
$`g(x,\rho )=g_{(0)}+\mathrm{}+\rho ^{d/2}g_{(d)}+h_{(d)}\rho ^{d/2}\mathrm{log}\rho +\mathrm{}`$ (2.3)
where the logarithmic piece appears only for even $`d`$. The sub-index in the metric expansion (and in all other expansions that appear in this paper) indicates the number of derivatives involved in that term, i.e. $`g_{(2)}`$ contains two derivatives, $`g_{(4)}`$ four derivatives, etc. It follows that the perturbative expansion in $`\rho `$ is also a low energy expansion. We set $`l=1`$ from now on. One can easily reinstate the factors of $`l`$ by dimensional analysis.
One can check that the curvature of $`G`$ satisfies
$$R_{\kappa \lambda \mu \nu }[G]=(G_{\kappa \mu }G_{\lambda \nu }G_{\kappa \nu }G_{\lambda \mu })+𝒪(\rho )$$
(2.4)
In this sense the metric is asymptotically anti-de Sitter. The Dirichlet problem for Einstein metrics satisfying (2.4) exactly (i.e. not only to leading order in $`\rho `$) was solved in .
In the coordinate system (2), Einstein’s equations read
$`\rho [2g^{\prime \prime }2g^{}g^1g^{}+\mathrm{Tr}(g^1g^{})g^{}]+\mathrm{Ric}(g)(d2)g^{}\mathrm{Tr}(g^1g^{})g`$ $`=`$ $`0`$ (2.5)
$`_i\mathrm{Tr}(g^1g^{})^jg_{ij}^{}`$ $`=`$ $`0`$ (2.6)
$`\mathrm{Tr}(g^1g^{\prime \prime }){\displaystyle \frac{1}{2}}\mathrm{Tr}(g^1g^{}g^1g^{})`$ $`=`$ $`0,`$ (2.7)
where differentiation with respect to $`\rho `$ is denoted with a prime, $`_i`$ is the covariant derivative constructed from the metric $`g`$, and $`\mathrm{Ric}(g)`$ is the Ricci tensor of $`g`$.
These equations are solved order by order in $`\rho `$. This is achieved by differentiating the equations with respect to $`\rho `$ and then setting $`\rho =0`$. For even $`d`$, this process would have broken down at order $`d/2`$ if the logarithm was not introduced in (2). $`h_{(d)}`$ is traceless, $`\mathrm{Tr}g_{(0)}^1h_{(d)}=0`$, and covariantly conserved, $`^ih_{(d)ij}=0`$. We show in appendix C that $`h_{(d)}`$ is proportional to the metric variation of the corresponding conformal anomaly, i.e. it is proportional to the stress-energy tensor of the theory with action the conformal anomaly. In any dimension, only the trace of $`g_{(d)}`$ and its covariant divergence are determined. Here is where extra data from the CFT are needed: as we shall see, the undetermined part is specified by the expectation value of the dual stress-energy tensor.
We collect in appendix A the results for $`g_{(n)}`$, $`h_{(d)}`$ as well as the results for the trace and divergence $`g_{(d)}`$. In dimension $`d`$ the latter are the only constraints that equations (2.7) yield for $`g_{(d)}`$. From this information we can parametrize the indeterminacy by finding the most general $`g_{(d)}`$ that has the determined trace and divergence.
In $`d=2`$ and $`d=4`$ the equation for the coefficient $`g_{(d)}`$ has the form of a conservation law
$$^ig_{(d)ij}=^iA_{(d)ij},d=2,4$$
(2.8)
where $`A_{(d)ij}`$ is a symmetric tensor explicitly constructed from the coefficients $`g_{(n)},n<d`$. The precise form of the tensor $`A_{(d)ij}`$ is given in appendix A (eq.(A.8)). The integration of this equation obviously involves an “integration constant” $`t_{ij}(x)`$, a symmetric covariantly conserved tensor the precise form of which can not be determined from Einstein’s equations.
In two dimensions, we get (see also )
$$g_{(2)ij}=\frac{1}{2}(Rg_{(0)ij}+t_{ij}),$$
(2.9)
where the symmetric tensor $`t_{ij}`$ should satisfy
$$^it_{ij}=0,\mathrm{Tr}t=R.$$
(2.10)
In four dimensions we obtain<sup>9</sup><sup>9</sup>9From now on we will suppress factors of $`g_{(0)}`$. For instance, $`\mathrm{Tr}g_{(2)}g_{(4)}=\mathrm{Tr}[g_{(0)}^1g_{(2)}g_{(0)}^1g_{(4)}]`$. Unless we explicitly mention to the contrary, indices will be raised and lowered with the metric $`g_{(0)}`$, all contractions will be made with this metric.
$$g_{(4)ij}=\frac{1}{8}g_{(0)ij}[(\mathrm{Tr}g_{(2)})^2\mathrm{Tr}g_{(2)}^2]+\frac{1}{2}(g_{(2)}^2)_{ij}\frac{1}{4}g_{(2)ij}\mathrm{Tr}g_{(2)}+t_{ij},$$
(2.11)
The tensor $`t_{ij}`$ satisfies
$$^it_{ij}=0,\mathrm{Tr}t=\frac{1}{4}[(\mathrm{Tr}g_{(2)})^2\mathrm{Tr}g_{(2)}^2].$$
(2.12)
In six dimensions the equation determining the coefficient $`g_{(6)}`$ is more subtle than the one in (2.8). It given by
$$^ig_{(6)ij}=^iA_{(6)ij}+\frac{1}{6}\mathrm{Tr}(g_{(4)}_jg_{(2)})$$
(2.13)
where the tensor $`A_{(6)ij}`$ is given in (A.8). It contains a part which is antisymmetric in the indices $`i`$ and $`j`$. Since $`g_{(6)ij}`$ is by definition a symmetric tensor the integration of equation (2.13) is not straightforward. Moreover, it is not obvious that the last term in (2.13) takes a form of divergence of some local tensor. Nevertheless, this is indeed the case as we now show. Let us define the tensor $`S_{ij}`$,
$`S_{ij}=^2C_{ij}2R_{ij}^{kl}C_{kl}+4(g_{(2)}g_{(4)}g_{(4)}g_{(2)})_{ij}+{\displaystyle \frac{1}{10}}(_i_jBg_{(0)ij}^2B)`$ (2.14)
$`+{\displaystyle \frac{2}{5}}g_{(2)ij}B+g_{(0)ij}({\displaystyle \frac{2}{3}}\mathrm{Tr}g_{(2)}^3{\displaystyle \frac{4}{15}}(\mathrm{Tr}g_{(2)})^3+{\displaystyle \frac{3}{5}}\mathrm{Tr}g_{(2)}\mathrm{Tr}g_{(2)}^2),`$
where
$$C_{ij}=(g_{(4)}\frac{1}{2}g_{(2)}^2+\frac{1}{4}g_{(2)}\mathrm{Tr}g_{(2)})_{ij}+\frac{1}{8}g_{(0)ij}B,B=\mathrm{Tr}g_2^2(\mathrm{Tr}g_2)^2.$$
The tensor $`S_{ij}`$ is a local function of the Riemann tensor. Its divergence and trace read
$$^iS_{ij}=4\mathrm{T}\mathrm{r}(g_{(4)}_jg_{(2)}),\mathrm{Tr}S=8\mathrm{T}\mathrm{r}(g_{(2)}g_{(4)}).$$
(2.15)
With the help of the tensor $`S_{ij}`$ the equation (2.13) can be integrated in a way similar to the $`d=2,4`$ cases. One obtains
$$g_{(6)ij}=A_{(6)ij}\frac{1}{24}S_{ij}+t_{ij}.$$
(2.16)
Notice that tensor $`S_{ij}`$ contains an antisymmetric part which cancels the antisymmetric part of the tensor $`A_{(6)ij}`$ so that $`g_{(6)ij}`$ and $`t_{ij}`$ are symmetric tensors, as they should. The symmetric tensor $`t_{ij}`$ satisfies
$$^it_{ij}=0,\mathrm{Tr}t=\frac{1}{3}[\frac{1}{8}(\mathrm{Tr}g_{(2)})^3\frac{3}{8}\mathrm{Tr}g_{(2)}\mathrm{Tr}g_{(2)}^2+\frac{1}{2}\mathrm{Tr}g_{(2)}^3\mathrm{Tr}g_{(2)}g_{(4)}].$$
(2.17)
Notice that in all three cases, $`d=2,4,6`$, the trace of $`t_{ij}`$ is proportional to the holographic conformal anomaly. As we will see in the next section, the symmetric tensors $`t_{ij}`$ are directly related to the expectation value of the boundary stress-energy tensor.
When $`d`$ is odd the only constraint on the coefficient $`g_{(d)ij}(x)`$ is that it is conserved and traceless
$$^ig_{(d)ij}=0,\mathrm{Tr}g_{(d)}=0.$$
(2.18)
So that we may identify
$$g_{(d)ij}=t_{ij}.$$
(2.19)
## 3 The holographic stress-energy tensor
We have seen in the previous section that given a conformal structure at infinity we can determine an asymptotic expansion of the metric up to order $`\rho ^{d/2}`$. We will now show that this term is determined by the expectation value of the dual stress-energy tensor.
According to the AdS/CFT prescription, the expectation value of the boundary stress-energy tensor is determined by functionally differentiating the on-shell gravitational action with respect to the boundary metric. The on-shell gravitational action, however, diverges. To regulate the theory we restrict the bulk integral to the region $`\rho ϵ`$ and we evaluate the boundary term at $`\rho =ϵ`$. The regulated action is given by
$`S_{\text{gr,reg}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi G_\text{N}}}\left[{\displaystyle _{\rho ϵ}}\text{d}^{d+1}x\sqrt{G}(R[G]+2\mathrm{\Lambda }){\displaystyle _{\rho =ϵ}}\text{d}^dx\sqrt{\gamma }\mathrm{\hspace{0.17em}2}K\right]=`$
$`=`$ $`{\displaystyle \frac{1}{16\pi G_\text{N}}}{\displaystyle \text{d}^dx\left[_ϵ\text{d}\rho \frac{d}{\rho ^{d/2+1}}\sqrt{detg(x,\rho )}+\frac{1}{\rho ^{d/2}}(2d\sqrt{detg(x,\rho )}+4\rho _\rho \sqrt{detg(x,\rho )})|_{\rho =ϵ}\right]}`$
Evaluating (3) for the solution we obtained in the previous section we find that the divergences appears as $`1/ϵ^k`$ poles plus a logarithmic divergence ,
$$S_{\text{gr,reg}}=\frac{l}{16\pi G_\text{N}}\text{d}^dx\sqrt{detg_{(0)}}\left(ϵ^{d/2}a_{(0)}+ϵ^{d/2+1}a_{(2)}+\mathrm{}+ϵ^1a_{(d2)}\mathrm{log}ϵa_{(d)}\right)+𝒪(ϵ^0),$$
(3.2)
where the coefficients $`a_{(n)}`$ are local covariant expressions of the metric $`g_{(0)}`$ and its curvature tensor. We give the explicit expressions, up to the order we are interested in, in appendix B.
We now obtain the renormalized action by subtracting the divergent terms, and then removing the regulator,
$$S_{\text{gr,ren}}[g_{(0)}]=\underset{ϵ0}{lim}\frac{1}{16\pi G_\text{N}}[S_{\text{gr,reg}}\text{d}^dx\sqrt{detg_{(0)}}\left(ϵ^{d/2}a_{(0)}+ϵ^{d/2+1}a_{(2)}+\mathrm{}+ϵ^1a_{(d2)}\mathrm{log}ϵa_{(d)}\right)]$$
(3.3)
The expectation value of the stress-energy tensor of the dual theory is given by
$$T_{ij}=\frac{2}{\sqrt{detg_{(0)}}}\frac{S_{\text{gr,ren}}}{g_{(0)}^{ij}}=\underset{ϵ0}{lim}\frac{2}{\sqrt{detg(x,ϵ)}}\frac{S_{\text{gr,ren}}}{g^{ij}(x,ϵ)}=\underset{ϵ0}{lim}\left(\frac{1}{ϵ^{d/21}}T_{ij}[\gamma ]\right)$$
(3.4)
where $`T_{ij}[\gamma ]`$ is the stress-energy tensor of the theory at $`\rho =ϵ`$ described by the action in (3.3) but before the limit $`ϵ0`$ is taken ($`\gamma _{ij}=1/ϵg_{ij}(x,ϵ)`$ is the induced metric at $`\rho =ϵ`$). Notice that the asymptotic expansion of the metric only allows for the determination of the divergences of the on-shell action. We can still obtain, however, a formula for $`T_{ij}`$ in terms of $`g_{(n)}`$ since, as (3.4) shows, we only need to know the first $`ϵ^{d/21}`$ orders in the expansion of $`T_{ij}[\gamma ]`$.
The stress-energy tensor $`T_{ij}[\gamma ]`$ contains two contributions,
$$T_{ij}[\gamma ]=T_{ij}^{\text{reg}}+T_{ij}^{\text{ct}},$$
(3.5)
$`T_{ij}^{\text{reg}}`$ comes from the regulated action in (3) and $`T_{ij}^{\text{ct}}`$ is due to the counterterms. The first contribution is equal to
$$T_{ij}^{\text{reg}}[\gamma ]=\frac{1}{8\pi G_\text{N}}(K_{ij}K\gamma _{ij})=\frac{1}{8\pi G_\text{N}}(_ϵg_{ij}(x,ϵ)+g_{ij}(x,ϵ)\mathrm{Tr}[g^1(x,ϵ)_ϵg(x,ϵ)]+\frac{1d}{ϵ}g_{ij}(x,ϵ))$$
(3.6)
The contribution due to counterterms can be obtained from the results in appendix B. It is given by
$`T_{ij}^{\text{ct}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G_\text{N}}}((d1)\gamma _{ij}+{\displaystyle \frac{1}{(d2)}}(R_{ij}{\displaystyle \frac{1}{2}}R\gamma _{ij})`$ (3.7)
$`{\displaystyle \frac{1}{(d4)(d2)^2}}[^2R_{ij}+2R_{ikjl}R^{kl}+{\displaystyle \frac{d2}{2(d1)}}_i_jR{\displaystyle \frac{d}{2(d1)}}RR_{ij}`$
$`{\displaystyle \frac{1}{2}}\gamma _{ij}(R_{kl}R^{kl}{\displaystyle \frac{d}{4(d1)}}R^2{\displaystyle \frac{1}{d1}}^2R)]T_{ij}^a\mathrm{log}ϵ)`$
where $`T_{ij}^a`$ is the stress-energy tensor of the action $`\text{d}^dx\sqrt{det\gamma }a_{(d)}`$. As it is shown in Appendix C, $`T_{ij}^a`$ is proportional to the tensor $`h_{(d)ij}`$ appearing in the expansion (2).
The stress tensor $`T_{ij}[g_{(0)}]`$ is covariantly conserved with respect to the metric $`g_{(0)ij}`$. To see this, notice that each of $`T_{ij}^{\text{reg}}`$ and $`T_{ij}^{\text{ct}}`$ is separately covariantly conserved with respect to the induced metric $`\gamma _{ij}`$ at $`\rho =ϵ`$: for $`T_{ij}^{\text{reg}}`$ one can check this by using the second equation in (2.7), for $`T_{ij}^{\text{ct}}`$ this follows from the fact that it was obtained by varying a local covariant counterterm. Since all divergences cancel in (3.4), we obtain that the finite part in (3.4) is conserved with respect to the metric $`g_{(0)ij}`$.
We are now ready to calculate $`T_{ij}`$. By construction (and we will verify this below) the divergent pieces cancel between $`T^{\text{reg}}`$ and $`T^{\text{ct}}`$.
### 3.1 $`d=2`$
In two dimensions we obtain
$$T_{ij}=\frac{l}{16\pi G_\text{N}}t_{ij}$$
(3.8)
where we have used (2.9) and (2.10) and the fact that $`T_{ij}^a=0`$ since $`R`$ is a topological invariant (and reinstated the factor of $`l`$). As promised, $`t_{ij}`$ is directly related to the boundary stress-energy tensor. Taking the trace we obtain
$$T_i^i=\frac{c}{24\pi }R$$
(3.9)
where $`c=3l/2G_\text{N}`$, which is the correct conformal anomaly .
Using our results, one can immediately obtain the stress-energy tensor of the boundary theory associated with a given solution $`G`$ of the three dimensional Einstein equations: one needs to write the metric in the coordinate system (2) and then use the formula
$$T_{ij}=\frac{2l}{16\pi G_\text{N}}(g_{(2)ij}g_{(0)ij}\mathrm{Tr}g_{(2)}).$$
(3.10)
From the gravitational point of view this is the quasi-local stress energy tensor associated with the solution $`G`$.
### 3.2 $`d=4`$
To obtain $`T_{ij}`$ we first need to rewrite the expressions in $`T^{\text{ct}}`$ in terms of $`g_{(0)}`$. This can be done with the help of the relation
$$R_{ij}[\gamma ]=R_{ij}[g_{(0)}]+\frac{1}{4}ϵ(2R_{ik}R^k{}_{j}{}^{}2R_{ikjl}R^{kl}\frac{1}{3}_i_jR+^2R_{ij}\frac{1}{6}^2Rg_{(0)ij})+𝒪(ϵ^2).$$
(3.11)
After some algebra one obtains,
$`T_{ij}[g_{(0)}]={\displaystyle \frac{1}{8\pi G_\text{N}}}\underset{ϵ0}{lim}`$ $`[{\displaystyle \frac{1}{ϵ}}(g_{(2)ij}+g_{(0)ij}\mathrm{Tr}g_{(2)}+{\displaystyle \frac{1}{2}}R_{ij}{\displaystyle \frac{1}{4}}g_{(0)ij}R)`$ (3.12)
$`+\mathrm{log}ϵ(2h_{(4)ij}T_{ij}^a)`$
$`2g_{(4)ij}h_{(4)ij}g_{(2)ij}\mathrm{Tr}g_{(2)}{\displaystyle \frac{1}{2}}g_{(0)ij}\mathrm{Tr}g_{(2)}^2`$
$`{\displaystyle \frac{1}{8}}(R_{ik}R^k{}_{j}{}^{}2R_{ikjl}R^{kl}{\displaystyle \frac{1}{3}}_i_jR+^2R_{ij}{\displaystyle \frac{1}{6}}^2Rg_{(0)ij})`$
$`{\displaystyle \frac{1}{4}}g_{(2)ij}R+{\displaystyle \frac{1}{8}}g_{(0)ij}(R_{kl}R^{kl}{\displaystyle \frac{1}{6}}R^2)].`$
Using the explicit expression for $`g_{(2)}`$ and $`h_{(4)}`$ given in (A.1) and (A.10) one finds that both the $`1/ϵ`$ pole and the logarithmic divergence cancel. Notice that had we not subtracted the logarithmic divergence from the action, the resulting stress-energy tensor would have been singular in the limit $`ϵ0`$.
Using (2.11) and (2.12) and after some algebra we obtain
$$T_{ij}=\frac{1}{8\pi G_\text{N}}[2t_{ij}3h_{(4)}].$$
(3.13)
Taking the trace we get
$$T_i^i=\frac{1}{16\pi G_\text{N}}(2a_{(4)}),$$
(3.14)
which is the correct conformal anomaly .
Notice that since $`h_{(4)ij}=\frac{1}{2}T_{ij}^a`$ the contribution in the boundary stress energy tensor proportional to $`h_{(4)ij}`$ is scheme dependent. Adding a local finite counterterm proportional to the trace anomaly will change the coefficient of this term. One may remove this contribution from the boundary stress energy tensor by a choice of scheme.
Finally, one can obtain the energy-momentum tensor of the boundary theory for a given solution $`G`$ of the five dimensional Einstein equations with negative cosmological constant. It is given by
$$T_{ij}=\frac{4}{16\pi G_\text{N}}[g_{(4)ij}\frac{1}{8}g_{(0)ij}[(\mathrm{Tr}g_{(2)})^2\mathrm{Tr}g_{(2)}^2]\frac{1}{2}(g_{(2)}^2)_{ij}+\frac{1}{4}g_{(2)ij}\mathrm{Tr}g_{(2)}],$$
(3.15)
where we have omitted the scheme dependent $`h_{(4)}`$ terms. From the gravitational point of view this is the quasi-local stress energy tensor associated with the solution $`G`$.
### 3.3 $`d=6`$
The calculation of the boundary stress tensor in $`d=6`$ case goes along the same lines as in $`d=2`$ and $`d=4`$ cases although it is technically involved. Up to local traceless covariantly conserved term (proportional to $`h_{(6)}`$) the results is
$$T_{ij}=\frac{3}{8\pi G_\text{N}}(g_{(6)ij}A_{(6)ij}+\frac{1}{24}S_{ij}).$$
(3.16)
where $`A_{(6)ij}`$ is given in (A.8) and $`S_{ij}`$ in (2.14). It is covariantly conserved and has the correct trace
$$T_i^i=\frac{1}{8\pi G_\text{N}}(a_{(6)}),$$
(3.17)
reproducing correctly the conformal anomaly in six dimensions .
Given an asymptotically AdS solution in six dimensions equation (3.16) yields the quasi-local stress energy tensor associated with it.
### 3.4 $`d=2k+1`$
In this case one can check that the counterterms only cancel infinities. Evaluating the finite part we get
$$T_{ij}=\frac{d}{16\pi G_\text{N}}g_{(d)ij}.$$
(3.18)
where $`g_{(d)ij}`$ can be identified with a traceless covariantly conserved tensor $`t_{ij}`$. In odd boundary dimensions there are no gravitational conformal anomalies, and indeed (3.18) is traceless. As in all previous cases, one can also read (3.18) as giving the quasi-local stress energy tensor associated with a given solution of Einstein’s equations.
### 3.5 Conformally flat bulk metrics
In this subsection we discuss a special case where the bulk metric can be determined to all orders given only a boundary metric. It was shown in that, given a conformally flat boundary metric, equations (2.7) can be integrated to all orders if the bulk Weyl tensor vanishes<sup>10</sup><sup>10</sup>10 In it was proven that if the bulk metric satisfies Einstein’s equations and it has a vanishing Weyl tensor, then the corresponding boundary metric has to be conformally flat. The converse is not necessarily true: one can have Einstein metrics with non-vanishing Weyl tensor which induce a conformally flat metric in the boundary.. We show that the extra condition in the bulk metric singles out a specific vacuum of the CFT.
The solution obtained in is given by
$$g(x,\rho )=g_{(0)}(x)+g_{(2)}(x)\rho +g_{(4)}(x)\rho ^2,g_{(4)}=\frac{1}{4}(g_{(2)})^2$$
(3.19)
where $`g_{(2)}`$ is given in (A.1) (we consider $`d>2`$), and all other coefficients $`g_{(n)}`$, $`n>4`$ vanish. Since $`g_{(4)}`$ and $`g_{(6)}`$ are now known, one can obtain a local formula for the dual stress energy tensor in terms of the curvature by using (2.11) and (2.16).
In $`d=4`$, using (2.11) and $`g_{(4)}=\frac{1}{4}(g_{(2)})^2`$, one obtains
$$t_{ij}=t_{ij}^{\text{cf}}\frac{1}{4}(g_{(2)})_{ij}^2+\frac{1}{4}g_{(2)ij}\mathrm{Tr}g_{(2)}\frac{1}{8}g_{(0)ij}[(\mathrm{Tr}g_{(2)})^2\mathrm{Tr}g_{(2)}^2].$$
(3.20)
It is easy to check that trace of $`t_{ij}^{\text{cf}}`$ reproduces (2.12). Furthermore, by virtue of Bianchi’s, one can show that $`t_{ij}^{\text{cf}}`$ is covariantly conserved. It is well-known that the stress-energy tensor of a quantum field theory on a conformally flat spacetime is a local function of the curvature tensor (see for example ). Our equation (3.20) reproduces the corresponding expression given in .
In $`d=6`$, using (2.16) and $`g_{(6)}=0`$ we find
$`t_{ij}=t_{ij}^{\text{cf}}[{\displaystyle \frac{1}{4}}g_{(2)}^3{\displaystyle \frac{1}{4}}g_{(2)}^2\mathrm{Tr}g_{(2)}+{\displaystyle \frac{1}{8}}g_{(2)}(\mathrm{Tr}g_{(2)})^2{\displaystyle \frac{1}{8}}g_{(2)}\mathrm{Tr}g_{(2)}`$
$`+g_{(0)}({\displaystyle \frac{1}{8}}\mathrm{Tr}g_{(2)}\mathrm{Tr}g_{(2)}^2{\displaystyle \frac{1}{12}}\mathrm{Tr}g_{(2)}^3{\displaystyle \frac{1}{24}}(\mathrm{Tr}g_{(2)})^3)]_{ij}.`$ (3.21)
One can verify that the trace of $`t_{ij}^{\text{cf}}`$ reproduces (2.17) (taking into account that $`g_{(4)}=\frac{1}{4}g_{(2)}^2`$) and that $`t_{ij}^{\text{cf}}`$ is covariantly conserved (by virtue of Bianchi’s).
Following the analysis in the previous subsections we obtain
$$T_{ij}=\frac{d}{16\pi G_\text{N}}t_{ij}^{\text{cf}}.$$
(3.22)
So, we explicitly see that the global condition we imposed on the bulk metric implies that we have picked a particular vacuum in the conformal field theory.
Note that the tensors $`t_{ij}^{\text{cf}}`$ in (3.20), (3.5) are local polynomial functions of the Ricci scalar and the Ricci tensor (but not of the Riemann tensor) of the metric $`g_{(0)ij}`$. It is perhaps an expected but still a surprising result that in conformally flat backgrounds the anomalous stress tensor is a local function of the curvature.
## 4 Conformal transformation properties of the stress-energy tensor
In this section we discuss the conformal transformation properties of the stress-energy tensor. These can be obtained by noting that conformal transformations in the boundary originate from specific diffeomorphisms that preserve the form of the metric (2). Under these diffeomorphisms $`g_{ij}(x,\rho )`$ transforms infinitesimally as
$$\delta g_{ij}(x,\rho )=2\sigma (1\rho _\rho )g_{ij}(x,\rho )+_ia_j(x,\rho )+_ja_i(x,\rho ),$$
(4.1)
where $`a_j(x,\rho )`$ is obtained from the equation
$$a^i(x,\rho )=\frac{1}{2}_0^\rho \text{d}\rho ^{}g^{ij}(x,\rho ^{})_j\sigma (x).$$
(4.2)
This can be integrated perturbatively in $`\rho `$,
$$a^i(x,\rho )=\underset{k=1}{}a_{(k)}^i\rho ^k.$$
(4.3)
We will need the first two terms in this expansion,
$$a_{(1)}^i=\frac{1}{2}^i\sigma ,a_{(2)}^i=\frac{1}{4}g_{(2)}^{ij}_j\sigma .$$
(4.4)
We can now obtain the way the $`g_{(n)}`$’s transform under conformal transformations
$`\delta g_{(0)ij}=2\sigma g_{(0)ij},`$
$`\delta g_{(2)ij}=_ia_{(1)j}+_ja_{(1)i}`$
$`\delta g_{(3)ij}=\sigma g_{(3)ij},`$
$`\delta g_{(4)ij}=2\sigma (g_{(4)}+h_{(4)})+a_{(1)}^k_kg_{(2)ij}+_ia_{(2)j}+_ja_{(2)i}+g_{(2)ik}_ja_{(1)}^k+g_{(2)jk}_ia_{(1)}^k`$
$`\delta g_{(5)ij}=3\sigma g_{(3)ij},`$ (4.5)
where the term $`h_{(4)}`$ in $`g_{(4)}`$ is only present when $`d=4`$. One can check from the explicit expressions for $`g_{(2)}`$ and $`g_{(4)}`$ in (A.1) that they indeed transform as (4). An alternative way to derive the transformation rules above is to start from (A.1) and perform a conformal variation. In the variations (4) were integrated leading to (A.1) up to conformally invariant terms.
Equipped with these results and the explicit form of the energy-momentum tensors, we can now easily calculate how the quantum stress-energy tensor transforms under conformal transformations. We use the term “quantum stress-energy tensor” because it incorporates the conformal anomaly. In the literature such transformation rules were obtained by first integrating the conformal anomaly to an effective action. This effective action is a functional of the initial metric $`g`$ and of the conformal factor $`\sigma `$. It can be shown that the difference between the stress-energy tensor of the theory on the manifold with metric $`ge^{2\sigma }`$ and the one on the manifold with metric $`g`$ is given by the stress-energy tensor derived by varying the effective action with respect to $`g`$.
In any dimension the stress-energy tensor transforms classically under conformal transformations as
$$\delta T_{\mu \nu }=(d2)\sigma T_{\mu \nu }$$
(4.6)
This transformation law is modified by the quantum conformal anomaly. In odd dimensions, where there is no conformal anomaly, the classical transformation rule (4.6) holds also at the quantum level. Indeed, for odd $`d`$, and by using (3.18) and (4), one easily verifies that the holographic stress-energy tensor transforms correctly.
In even dimensions, the transformation (4.6) is modified. In $`d=2`$, it is well-known that one gets an extra contribution proportional to the central charge. Indeed, using (3.10) and the formulae above we obtain
$$\delta T_{ij}=\frac{l}{8\pi G_\text{N}}(_i_j\sigma g_{(0)ij}^2\sigma )=\frac{c}{12}(_i_j\sigma g_{(0)ij}^2\sigma ),$$
(4.7)
which is the correct transformation rule.
In $`d=4`$ we obtain,
$`\delta T_{ij}`$ $`=`$ $`2\sigma T_{ij}+{\displaystyle \frac{1}{4\pi G_\text{N}}}(2\sigma h_{(4)}+{\displaystyle \frac{1}{4}}^k\sigma [_kR_{ij}{\displaystyle \frac{1}{2}}(_iR_{jk}+_jR_{ik}){\displaystyle \frac{1}{6}}_kRg_{(0)ij}]`$ (4.8)
$`+{\displaystyle \frac{1}{48}}(_i\sigma _jR+_i\sigma _jR)+{\displaystyle \frac{1}{12}}R(_i_j\sigma g_{(0)ij}^2\sigma )`$
$`+{\displaystyle \frac{1}{8}}[R_{ij}^2\sigma (R_{ik}^k_j\sigma +R_{jk}^k_i\sigma )+g_{(0)ij}R_{kl}^k^l\sigma ]).`$
The only other result known to us is the result in , where they computed the finite conformal transformation of the stress-energy tensor but for a conformally flat metric $`g_{(0)}`$. For conformally flat backgrounds, $`h_{(4)}`$ vanishes because it is the metric variation of a topological invariant. The terms proportional to a single derivative of $`\sigma `$ vanish by virtue of Bianchi identities and the fact that the Weyl tensor vanishes for conformally flat metrics. The remaining terms, which only contain second derivatives of $`\sigma `$, can be shown to coincide with the infinitesimal version of (4.23) in .
One can obtain the conformal transformation of the stress energy tensor in $`d=6`$ in a similar fashion but we shall not present this result here.
## 5 Matter
In the previous sections we examined how spacetime is reconstructed (to leading order) holographically out of CFT data. In this section we wish to examine how field theory describing matter on this spacetime is encoded in the CFT. We will discuss scalar fields but the techniques are readily applicable to other kinds of matter.
The method we will use is the same as in the case of pure gravity, i.e. we will start by specifying the sources that are turned on, find how far we can go with only this information and then input more CFT data. We will find the same pattern: knowledge of the sources allows only for determination of the divergent part of the action. The leading finite part (which depends on global issues and/or the signature of spacetime) is determined by the expectation value of the dual operator. We would like to stress that in the approach we follow, i.e. regularize, subtract all infinities by adding counterterms and finally remove the regulator to obtain the renormalized action, all normalizations of the physical correlation functions are fixed and are consistent with Ward identities.
Other papers that discuss similar issues include .
### 5.1 Dirichlet boundary problem for scalar fields in a fixed gravitational background
In this section we consider scalars on a fixed gravitational background. This is taken to be of the generic form (2). In most of the literature the fixed metric was taken to be that of standard AdS, but with not much more effort one can consider the general case.
The action for massive scalar is given by
$$S_\text{M}=\frac{1}{2}\text{d}^{d+1}x\sqrt{G}\left(G^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+m^2\mathrm{\Phi }^2\right)$$
(5.1)
where $`G_{\mu \nu }`$ has an expansion of the form (2).
We take the scalar field $`\mathrm{\Phi }`$ to have an expansion of the form
$$\mathrm{\Phi }(x,\rho )=\rho ^{(d\mathrm{\Delta })/2}\varphi (x,\rho ),\varphi (x,\rho )=\varphi _{(0)}+\varphi _{(2)}\rho +\mathrm{},$$
(5.2)
where $`\mathrm{\Delta }`$ is the conformal dimension of the dual operator. We take the dimension $`\mathrm{\Delta }`$ to be quantized as $`\mathrm{\Delta }=\frac{d}{2}+k,k=0,1,..`$. This is often the case for operators of protected dimension. For the case of scalars that correspond to operators of dimensions $`\frac{d}{2}1\mathrm{\Delta }<\frac{d}{2}`$ we refer to . Inserting (5.2) in the field equation,
$$(\text{ }\text{ }\text{ }\text{ }_G+m^2)\mathrm{\Phi }=0,$$
(5.3)
where $`\text{ }\text{ }\text{ }\text{ }_G\mathrm{\Phi }=\frac{1}{\sqrt{G}}_\mu (\sqrt{G}G^{\mu \nu }_\nu \mathrm{\Phi })`$, we obtain that the mass $`m^2`$ and the conformal dimension $`\mathrm{\Delta }`$ are related as $`m^2=(\mathrm{\Delta }d)\mathrm{\Delta }`$, and that $`\varphi `$ satisfies
$$[(d\mathrm{\Delta })_\rho \mathrm{log}g\varphi +2(2\mathrm{\Delta }d2)_\rho \varphi \text{ }\text{ }\text{ }\text{ }_g\varphi ]+\rho [2_\rho \mathrm{log}g_\rho \varphi 4_\rho ^2\varphi ]=0.$$
(5.4)
Given $`\varphi _{(0)}`$ one can determine recursively $`\varphi _{(n)},n>0`$. This is achieved by differentiating (5.4) and setting $`\rho `$ equal to zero. We give the result for the first couple of orders in appendix D. This process breaks down at order $`\mathrm{\Delta }d/2`$ (provided this is an integer, which we assume throughout this section) because the coefficient of $`\varphi _{(2\mathrm{\Delta }d)}`$ (the field to be determined) becomes zero. This is exactly analogous to the situation encountered for even $`d`$ in the gravitational sector. Exactly the same way as there, we introduce at this order a logarithmic term, i.e. the expansion of $`\mathrm{\Phi }`$ now reads,
$$\mathrm{\Phi }=\rho ^{(d\mathrm{\Delta })/2}(\varphi _{(0)}+\rho \varphi _{(2)}+\mathrm{})+\rho ^{\mathrm{\Delta }/2}(\varphi _{(2\mathrm{\Delta }d)}+\mathrm{log}\rho \psi _{(2\mathrm{\Delta }d)}+\mathrm{}).$$
(5.5)
The equation (5.4) now determines all terms up to $`\varphi _{(2\mathrm{\Delta }d2)}`$, the coefficient of the logarithmic term $`\psi _{(2\mathrm{\Delta }d)}`$, but leaves undetermined $`\varphi _{(2\mathrm{\Delta }d)}`$. This is analogous to the situation discussed in section 2 where the term $`g_{(d)}`$ was undetermined. It is well known that precisely at order $`\rho ^{\mathrm{\Delta }/2}`$ one finds the expectation value of the dual operator. We will review this argument below, and also derive the exact proportionality coefficient. Our result is in agreement with .
We proceed to regularize and then renormalize the theory. We regulate by integrating in the bulk from $`\rho ϵ`$,<sup>11</sup><sup>11</sup>11 This regularization for scalar fields in a fixed AdS background was considered in . In these papers the divergences were computed in momentum space, but no counterterms were added to cancel them. Addition of boundary counterterms to cancel infinities for scalar fields was considered in , and more recently in .
$`S_{\text{M}\text{,reg}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\rho ϵ}}\text{d}^{d+1}x\sqrt{G}\left(G^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+m^2\mathrm{\Phi }^2\right)`$
$`=`$ $`{\displaystyle _{\rho =ϵ}}\text{d}^dx\sqrt{g(x,ϵ)}ϵ^{\mathrm{\Delta }+d/2}[{\displaystyle \frac{1}{2}}(d\mathrm{\Delta })\varphi ^2(x,ϵ)+ϵ\varphi (x,ϵ)_ϵ\varphi (x,ϵ)]`$
$`=`$ $`{\displaystyle \text{d}^dx\sqrt{g_{(0)}}[ϵ^{\mathrm{\Delta }+d/2}a_{(0)}^\text{M}+ϵ^{\mathrm{\Delta }+d/2+1}a_{(2)}^\text{M}+\mathrm{}+ϵa_{(2\mathrm{\Delta }d+2)}^\text{M}\mathrm{log}ϵa_{(2\mathrm{\Delta }d)}]}+𝒪(ϵ^0)`$
Clearly, with $`\mathrm{\Delta }d/2`$ a positive integer there are finite number of divergent terms. The logarithmic divergence appears exactly when $`\mathrm{\Delta }=d/2+k,k=0,1,..`$, in agreement with the analysis in , and is directly related to the logarithmic term in (5.5). The first few of the power law divergences read
$$a_{(0)}^\text{M}=\frac{1}{2}(d\mathrm{\Delta })\varphi _{(0)}^2,a_{(2)}^\text{M}=\frac{1}{4}\mathrm{Tr}g_{(2)}\varphi _{(0)}^2+(d\mathrm{\Delta }+1)\varphi _{(0)}\varphi _{(2)}.$$
(5.7)
Given a field of specific dimension it is straightforward to compute all divergent terms.
We now proceed to obtain the renormalized action by adding counterterms to cancel the infinities,
$$S_{\text{M}\text{,ren}}=\underset{ϵ0}{lim}[S_{\text{M}\text{,reg}}\text{d}^dx\sqrt{g_{(0)}}[ϵ^{\mathrm{\Delta }+d/2}a_{(0)}^\text{M}+ϵ^{\mathrm{\Delta }+d/2+1}a_{(2)}^\text{M}+\mathrm{}+ϵa_{(2\mathrm{\Delta }d+2)}^\text{M}\mathrm{log}ϵa_{(2\mathrm{\Delta }d)}]$$
(5.8)
Exactly as in the case of pure gravity, and since the regulated theory lives at $`\rho =ϵ`$, one needs to rewrite the counterterms in terms of the field living at $`\rho =ϵ`$, i.e. the induced metric $`\gamma _{ij}(x,ϵ)`$ and the field $`\mathrm{\Phi }(x,ϵ)`$, or equivalently $`g_{ij}(x,ϵ)`$ and $`\varphi (x,ϵ)`$. This is straightforward but somewhat tedious: one needs to invert the relation between $`\varphi `$ and $`\varphi _{(0)}`$ and between $`g_{ij}`$ and $`g_{(0)ij}`$ to sufficiently high order. This then allows to express all $`\varphi _{(n)}`$, and therefore all $`a_{(n)}^\text{M}`$, in terms of $`\varphi (x,ϵ)`$ and $`g_{ij}(x,ϵ)`$ (the $`\varphi _{(n)}`$’s are determined in terms of $`\varphi _{(0)}`$ and $`g_{(0)}`$ by solving (5.4) iteratively). Explicitly, the first two orders read,
$`S_{\text{M}\text{,ren}}`$ $`=`$ $`\underset{ϵ0}{lim}[{\displaystyle \frac{1}{2}}{\displaystyle _{\rho ϵ}}\text{d}^{d+1}x\sqrt{G}(G^{\mu \nu }_\mu \mathrm{\Phi }_\nu \mathrm{\Phi }+m^2\mathrm{\Phi }^2)`$
$`+{\displaystyle _{\rho =ϵ}}\sqrt{\gamma }[{\displaystyle \frac{(d\mathrm{\Delta })}{2}}\mathrm{\Phi }^2(x,ϵ)+{\displaystyle \frac{1}{2(2\mathrm{\Delta }d2)}}(\mathrm{\Phi }(x,ϵ)\text{ }\text{ }\text{ }_\gamma \mathrm{\Phi }(x,ϵ)+{\displaystyle \frac{d\mathrm{\Delta }}{2(d1)}}R[\gamma ]\mathrm{\Phi }^2(x,ϵ))+\mathrm{}]]`$
The addition of the first counterterm was discussed in . The action (5.1) with only the counterterms written explicitly is finite for fields of $`\mathrm{\Delta }<d/2+2`$. As remarked above, it is straightforward to obtain all counterterms needed in order to make the action finite for any field of any mass. These counterterms contain also logarithmic subtractions that lead to the conformal anomalies discussed in . For instance, if $`\mathrm{\Delta }=\frac{1}{2}d+1`$, the coefficient $`[2(2\mathrm{\Delta }d2)]^1`$ in (5.1) is replaced by $`\frac{1}{4}\mathrm{log}ϵ`$. An alternative way to derive the counterterms is to demand that the expectation value $`O`$ is finite. This holds in the case of pure gravity too, i.e. the counterterms can also be derived by requiring finiteness of $`T_{\mu \nu }`$ .
The expectation value of the dual operator is given by
$$O(x)=\frac{1}{\sqrt{detg_{(0)}}}\frac{\delta S_{\text{M}\text{,ren}}}{\delta \varphi _{(0)}}=\underset{ϵ0}{lim}\frac{1}{\sqrt{detg(x,ϵ)}}\frac{\delta S_{\text{M}\text{,ren}}}{\delta \varphi (x,ϵ)}.$$
(5.10)
Exactly as in the case of pure gravity, the expectation value receives a contribution both from the regulated part and from the counterterms. We obtain,
$$O(x)=(2\mathrm{\Delta }d)\varphi _{(2\mathrm{\Delta }d)}+F(\varphi _{(n)},\psi _{(2\mathrm{\Delta }d)},g_{(m)}),n<2\mathrm{\Delta }d$$
(5.11)
where we used that $`\varphi _{(2\mathrm{\Delta }d)}`$ is linear in $`\varphi _{(0)}`$ (notice that the action (5.1) does not include interactions). $`F(\varphi _{(n)},\psi _{(2\mathrm{\Delta }d)},g_{(m)})`$ is a local function of $`\varphi _{(n)}`$ with $`n<2\mathrm{\Delta }d`$, $`\psi _{(2\mathrm{\Delta }d)}`$ and $`g_{(m)}`$. These terms are related to contact terms in correlation functions of $`O`$ with itself and with the stress-energy tensor. Its exact form is straightforward but somewhat tedious to obtain (just use (5.1) and (5.10)).
As we have promised, we have shown that the coefficient $`\varphi _{(2\mathrm{\Delta }d)}`$ is related with the expectation value of the dual CFT operator. In the case that the background geometry is the standard Euclidean AdS one can readily obtain $`\varphi _{(2\mathrm{\Delta }d)}`$ from the unique solution of the scalar field equation with given Dirichlet boundary conditions. One finds that $`\varphi _{(2\mathrm{\Delta }d)}`$ is proportional to (an integral involving) $`\varphi _{(0)}`$. Therefore, $`\varphi _{(2\mathrm{\Delta }d)}`$ carries information about the 2-point function. The factor $`(\mathrm{\Delta }d/2)`$ is crucial in order for the 2-point function to be normalized correctly . We refer to for a detailed discussion of this point.
We finish this section by calculating the conformal anomaly associated with the scalar fields and in the case the background is (locally) standard AdS (i.e. $`g_{(n)}=0`$, for $`0<n<d`$). Equation (5.4) simplifies and can be easily solved. One gets
$`\varphi _{(2n)}`$ $`=`$ $`{\displaystyle \frac{1}{2n(2\mathrm{\Delta }d2n)}}\text{ }\text{ }\text{ }_0\varphi _{(2n2)},`$
$`\psi _{(2\mathrm{\Delta }d)}`$ $`=`$ $`{\displaystyle \frac{1}{2(2\mathrm{\Delta }d)}}\text{ }\text{ }\text{ }_0\varphi _{(2\mathrm{\Delta }d2)}={\displaystyle \frac{1}{2^{2k}\mathrm{\Gamma }(k)\mathrm{\Gamma }(k+1)}}(\text{ }\text{ }\text{ }_0)^k\varphi _{(0)},`$ (5.12)
where $`k=\mathrm{\Delta }\frac{d}{2}`$ and $`\text{ }\text{ }\text{ }\text{ }_0`$ is the Laplacian of $`g_{(0)}`$. The regularized action written in terms of the fields at $`\rho =ϵ`$ contains the following explicit logarithmic divergence,
$$S_{\text{M}\text{,reg}}=_{\rho =ϵ}\text{d}^dx\sqrt{\gamma }[\mathrm{log}ϵ(\mathrm{\Delta }\frac{d}{2})\varphi (x,ϵ)\psi _{(2\mathrm{\Delta }d)}(x,ϵ)+\mathrm{}],$$
(5.13)
where the dots indicate power law divergent and finite terms, $`\psi _{(2\mathrm{\Delta }d)}(x,ϵ)`$ is given by (5.12) with $`g_{(0)}`$ replaced by $`\gamma `$ and $`\varphi _{(0)}`$ by $`\varphi (x,ϵ)`$. Using the same argument as in we obtain the matter conformal anomaly,
$$𝒜_\text{M}=\frac{1}{2}\left(\frac{1}{2^{2k2}(\mathrm{\Gamma }(k))^2}\right)\varphi _{(0)}(\text{ }\text{ }\text{ }\text{ }_0)^k\varphi _{(0)}.$$
(5.14)
This agrees exactly with the anomaly calculated in (compare with formulae (10), (37) in ).
### 5.2 Scalars coupled to gravity
In the previous section we ignored the back-reaction of the scalars to the bulk geometry. The purpose of this section is to discuss this issue. The action is now the sum of (2.1) and (5.1),
$$S=S_{\text{gr}}+S_\text{M}.$$
(5.15)
The gravitational field equation in the presence of matter reads
$$R_{\mu \nu }\frac{1}{2}(R+2\mathrm{\Lambda })G_{\mu \nu }=8\pi G_\text{N}T_{\mu \nu }$$
(5.16)
In the coordinate system (2) and with the scalar field having the expansion in (5.5), these equations read
$`\rho [2g_{ij}^{\prime \prime }2(g^{}g^1g^{})_{ij}+\mathrm{Tr}(g^1g^{})g_{ij}^{}]`$ $`+`$ $`R_{ij}(g)(d2)g_{ij}^{}\mathrm{Tr}(g^1g^{})g_{ij}=`$ (5.17)
$`=`$ $`8\pi G_\text{N}\rho ^{d\mathrm{\Delta }1}\left[{\displaystyle \frac{(\mathrm{\Delta }d)\mathrm{\Delta }}{d1}}\varphi ^2g_{ij}+\rho _i\varphi _j\varphi \right],`$
$`_i\mathrm{Tr}(g^1g^{})^jg_{ij}^{}`$ $`=`$ $`16\pi G_\text{N}\rho ^{d\mathrm{\Delta }1}\left[{\displaystyle \frac{d\mathrm{\Delta }}{2}}\varphi _i\varphi +\rho _\rho \varphi _i\varphi \right],`$
$`\mathrm{Tr}(g^1g^{\prime \prime }){\displaystyle \frac{1}{2}}\mathrm{Tr}(g^1g^{}g^1g^{})`$ $`=`$ $`16\pi G_\text{N}\rho ^{d\mathrm{\Delta }2}[{\displaystyle \frac{d(\mathrm{\Delta }d)(\mathrm{\Delta }d+1)}{4(d1)}}\varphi ^2`$
$`+`$ $`(d\mathrm{\Delta })\rho \varphi _\rho \varphi +\rho ^2(_\rho \varphi )^2{\displaystyle \frac{}{}}],`$
If $`\mathrm{\Delta }>d`$, the right-hand side diverges near the boundary whereas the left-hand side is finite. Operators with dimension $`\mathrm{\Delta }>d`$ are irrelevant operators. Correlation functions of these operators have a very complicated singularity structure at coincident points. As remarked in , one can avoid such problems by considering the sources to be infinitesimal and to have disjoint support, so that these operators are never at coincident points. Requiring that the equations in (5.17) are satisfied to leading order in $`\rho `$ yields
$$\varphi _{(0)}^2=0,$$
(5.18)
which is indeed the prescription advocated in .
If $`\mathrm{\Delta }d`$, which means that we deal with marginal or relevant operators, one can perturbatively calculate the back-reaction of the scalars to the bulk metric. At which order the leading back-reaction appears depends on the mass of the field. For fields that correspond to operators of dimension $`\mathrm{\Delta }=dk`$ the leading back-reaction appears at order $`\rho ^k`$, except when $`k=0`$ (marginal operators), where the leading back-reaction is at order $`\rho `$.
Let us see how conformal anomalies arise in this context. The logarithmic divergences are coming from the regulated on-shell value of the bulk integral in (5.15). The latter reads
$`S_{\text{reg}}(\text{bulk})`$ $`=`$ $`{\displaystyle _{\rho ϵ}}\text{d}\rho \text{d}^dx\sqrt{G}[{\displaystyle \frac{d}{8\pi G_\text{N}}}{\displaystyle \frac{m^2}{d1}}\mathrm{\Phi }^2]`$ (5.19)
$`=`$ $`{\displaystyle _{\rho ϵ}}\text{d}\rho \text{d}^dx{\displaystyle \frac{1}{\rho }}\sqrt{g(x,\rho )}[{\displaystyle \frac{d}{16\pi G_\text{N}}}\rho ^{d/2}{\displaystyle \frac{m^2}{2(d1)}}\varphi ^2(x,\rho )\rho ^k]`$
where $`k=\mathrm{\Delta }d/2`$. We see that gravitational conformal anomalies are expected when $`d`$ is even and matter conformal anomalies when $`k`$ is a positive integer, as it should.
In the presence of sources the expectation value of the boundary stress-energy tensor is not conserved but rather it satisfies a Ward identity that relates its covariant divergence to the expectation value of the operators that couple to the sources. To see this consider the generating functional
$$Z_{\text{CFT}}[g_{(0)},\varphi _{(0)}]=<\mathrm{exp}\text{d}^dx\sqrt{g_{(0)}}[\frac{1}{2}g_{(0)}^{ij}T_{ij}\varphi _{(0)}O]>.$$
(5.20)
Invariance under infinitesimal diffeomorphisms,
$$\delta g_{(0)ij}=_i\xi _j+_j\xi _i,$$
(5.21)
yields the Ward identity,
$$^jT_{ij}=O_i\varphi _{(0)}.$$
(5.22)
As we have remark before, $`T_{ij}`$ has a dual meaning , both as the expectation value of the dual stress-energy tensor and as the quasi-local stress-energy tensor of Brown and York. The Ward identity (5.22) has a natural explanation from the latter point in view as well. According to the quasi-local stress-energy tensor is not conserved in the presence of matter but it satisfies
$$^jT_{ij}=\tau _{i\rho }$$
(5.23)
where $`\tau _{i\rho }`$ expresses the flow of matter energy-momentum through the boundary. Evidently, (5.22) is of the form (5.23).
Solving the coupled system of equations (5.17) and (5.4) is straightforward but somewhat tedious. The details differ from case to case. For illustrative purposes we present a sample calculation: we consider the case of two-dimensional massless scalar field ($`d=\mathrm{\Delta }=2,k=1`$).
The equations to be solved are (5.4) and (5.17) with $`d=\mathrm{\Delta }=2`$ and the expansion of the metric and the scalar field are given by (2) and (5.5) (again with $`d=\mathrm{\Delta }=2`$), respectively. Equation (5.4) determines $`\psi _{(2)}`$,
$$\psi _{(2)}=\frac{1}{4}\text{ }\text{ }\text{ }\text{ }_0\varphi _{(0)}.$$
(5.24)
Equations (5.17) determine $`h_{(2)}`$, the trace of the $`g_{(2)}`$ and provide a relation between the divergence of $`g_{(2)}`$ and $`\varphi _{(2)}`$,
$`h_{(2)}=4\pi G_\text{N}\left(_i\varphi _{(0)}_j\varphi _{(0)}{\displaystyle \frac{1}{2}}g_{(0)ij}(\varphi _{(0)})^2\right),`$
$`\mathrm{Tr}g_{(2)}={\displaystyle \frac{1}{2}}R+4\pi G_\text{N}(\varphi _{(0)})^2,`$
$`^ig_{(2)ij}=_i\mathrm{Tr}g_{(2)}+16\pi G_\text{N}\varphi _{(2)}_i\varphi _{(0)}.`$ (5.25)
Notice that $`g_{(2)}`$ and $`\varphi _{(2)}`$ are still undetermined and are related to the expectation values of the dual operators (3.4) and (5.11), respectively. Notice that $`h_{(2)}`$ is equal to the stress-energy tensor of a massless two-dimensional scalar.
Going back to (5.19), we see that the second term drops out (since $`m^2=0`$) and one can use the result already obtained in the gravitational sector,
$$𝒜=\frac{1}{16\pi G_\text{N}}(2a_{(2)})=\frac{1}{16\pi G_\text{N}}(2\mathrm{T}\mathrm{r}g_{(2)})=\frac{1}{16\pi G_\text{N}}R+\frac{1}{2}\varphi _{(0)}\text{ }\text{ }\text{ }\text{ }_0\varphi _{(0)}\frac{1}{2}_i(\varphi _{(0)}^i\varphi _{(0)}),$$
(5.26)
which is the correct conformal anomaly (the last term can be removed by adding a covariant counterterm).
The renormalized boundary stress tensor reads
$$T_{ij}(x)=\frac{1}{8\pi G_\text{N}}\left(g_{(2)ij}+h_{(2)ij}g_{(0)ij}\mathrm{Tr}g_{(2)}\right)(x)$$
(5.27)
Its trace gives correctly the conformal anomaly (5.26). On the other hand, taking the covariant derivative of (5.27) we get
$`^jT_{ij}=O(x)_i\varphi _0(x)`$
$`O(x)=2(\varphi _2(x)+\psi _2(x)).`$ (5.28)
in agreement with equations (5.22) and (5.11).
## 6 Conclusions
Most of the discussions in the literature on the AdS/CFT correspondence are concerned with obtaining conformal field theory correlation functions using supergravity. In this paper we started investigating the converse question: how can one obtain information about the bulk theory from CFT correlation functions? How does one decode the hologram?
Answering these questions in all generality, but within the context of the AdS/CFT duality, entails developing a precise dictionary between bulk and boundary physics. A prescription for relating bulk/boundary observables is already available , and one would expect that it would allow us to reconstruct the bulk spacetime from the boundary CFT. The prescription of , however, relates infinite quantities. One of the main results of this paper is the systematic development of a renormalized version of this prescription. Equipped with it, and with no other assumption (except that the CFT has an AdS dual), we then proceeded to reconstruct the bulk spacetime metric and bulk scalar fields to the first non-trivial order.
Our approach to the problem is to start from the boundary and try to build iteratively bulk solutions. Within this approach, the pattern we find is the following:
$``$ Sources in the CFT determine an asymptotic expansion of the corresponding bulk field near the boundary to high enough order so that all infrared divergences of the bulk on-shell action can be computed. This then allows to obtain a renormalized on-shell action by adding boundary counterterms to cancel the infrared divergences.
$``$ Bulk solutions can be extended one order further by using the 1-point function of the corresponding dual CFT operator.
In the case the bulk field is the metric, our results show that a conformal structure at infinity is not in general sufficient in order to obtain a bulk metric. The first additional information one needs is the expectation value of the boundary stress energy tensor.
As a by-product, we have obtained ready-to-use formulae for the Brown-York quasi-local stress-energy tensor for arbitrary solution of Einstein’s equations with negative cosmological constant up to six dimensions. The six-dimensional result is particularly interesting because, via AdS/CFT, provides new information about the still mysterious $`(2,0)`$ theory. Furthermore, we have obtained the conformal transformation properties of the stress-energy tensors. These transformation rules incorporate the trace anomaly and provide a generalization to $`d>2`$ of the well-known Schwartzian derivative contribution in the conformal transformation rule of the stress-energy tensor in $`d=2`$.
Our discussion extends straightforwardly to the case of different matter. We expect that in all cases obstructions in extending the solution to the deep interior region will be resolved by additional CFT data (including data about non-local observables such as Wilson loops, Wilson surfaces etc.). An interesting case to study in this framework is point particles . Reconstructing the trajectory of the bulk point particle out of CFT data will present a model of how holography works with time dependent processes. Furthermore, following , one could study the interplay between causality and holography. Another extension is to study renormalization group flows using the present formalism. This amounts to extending the discussion in section 5.2 by adding a potential for the scalars. Another application of our results is in the context of Randall-Sundrum (RS) scenarios . Incorporating such a scenario in string theory, in the case the bulk space is AdS, may yield a connection with the AdS/CFT duality . As advocated in , one may view the RS scenario as $`4d`$ gravity coupled to a cut-off CFT. The regulated theory in our discussion provides a dual description of a cut-off CFT. In this context, the counterterms are re-interpreted as providing the action for the bulk modes localized in the brane . We see, for instance, that the counterterms in (5.1) can be re-interpreted as an action for a bulk scalar mode localized on the brane.
## Note added
As this paper was being finalized, appeared with some overlap with the results of section 2.
## Acknowledgments
We would like to thank G. ’t Hooft for reading the manuscript and his useful remarks. This research is supported in part by NSF grants PHY94-07194 and PHY-9802484. KS would like to thank ITP in UCSB for hospitality during initial stages of this work. SS would like to thank the Theory Division at CERN for the hospitality extended to him while this work was in progress.
## Appendix
## A Asymptotic solution of Einstein’s equations
In this appendix we collect the results for the solution of the equations (2.7) up to the order we are interested in.
From the first equation in (2.7) one determines the coefficients $`g_{(n)}`$, $`nd`$, in terms of $`g_{(0)}`$. For our purpose we only need $`g_{(2)}`$ and $`g_{(4)}`$. There are given by
$`g_{(2)}_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{d2}}(R_{ij}{\displaystyle \frac{1}{2(d1)}}Rg_{(0)}{}_{ij}{}^{})`$ (A.1)
$`g_{(4)}_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{d4}}({\displaystyle \frac{1}{8(d1)}}D_iD_jR+{\displaystyle \frac{1}{4(d2)}}D_kD^kR_{ij}`$ (A.5)
$`{\displaystyle \frac{1}{8(d1)(d2)}}D_kD^kRg_{(0)}{}_{ij}{}^{}{\displaystyle \frac{1}{2(d2)}}R^{kl}R_{ikjl}`$
$`+{\displaystyle \frac{d4}{2(d2)^2}}R_i{}_{}{}^{k}R_{kj}^{}+{\displaystyle \frac{1}{(d1)(d2)^2}}RR_{ij}`$
$`+{\displaystyle \frac{1}{4(d2)^2}}R^{kl}R_{kl}g_{(0)}{}_{ij}{}^{}{\displaystyle \frac{3d}{16(d1)^2(d2)^2}}R^2g_{(0)}{}_{ij}{}^{}).`$
The expressions for $`g_{(n)}`$ are singular when $`n=d`$. One can obtain the trace and the divergence of $`g_{(n)}`$ for any $`n`$ from the last two equations in (2.7). Explicitly,
$`\mathrm{Tr}g_{(4)}={\displaystyle \frac{1}{4}}\mathrm{Tr}g_{(2)}^2,\mathrm{Tr}g_{(6)}={\displaystyle \frac{2}{3}}\mathrm{Tr}g_{(2)}g_{(4)}{\displaystyle \frac{1}{6}}\mathrm{Tr}g_{(2)}^3`$
$`\mathrm{Tr}g_{(3)}=0,\mathrm{Tr}g_{(5)}=0,`$ (A.6)
and
$`^ig_{(2)ij}=^iA_{(2)ij},^ig_{(3)ij}=0,^ig_{(4)ij}=^iA_{(4)ij}`$
$`^ig_{(5)ij}=0,^ig_{(6)ij}=^iA_{(6)ij}+{\displaystyle \frac{1}{6}}\mathrm{Tr}(g_{(4)}_jg_{(2)}),`$ (A.7)
where
$`A_{(2)ij}`$ $`=`$ $`g_{(0)ij}\mathrm{Tr}g_{(2)}`$ (A.8)
$`A_{(4)ij}`$ $`=`$ $`{\displaystyle \frac{1}{8}}[\mathrm{Tr}g_{(2)}^2(\mathrm{Tr}g_{(2)})^2]g_{(0)ij}+{\displaystyle \frac{1}{2}}(g_{(2)}^2)_{ij}{\displaystyle \frac{1}{4}}g_{(2)ij}\mathrm{Tr}g_{(2)}`$
$`A_{(6)ij}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(2(g_{(2)}g_{(4)})_{ij}+(g_{(4)}g_{(2)})_{ij}(g_{(2)}^3)_{ij}+{\displaystyle \frac{1}{8}}[\mathrm{Tr}g_{(2)}^2(\mathrm{Tr}g_{(2)})^2]g_{(2)ij}`$
$``$ $`\mathrm{Tr}g_{(2)}[g_{(4)ij}{\displaystyle \frac{1}{2}}(g_{(2)}^2)_{ij}][{\displaystyle \frac{1}{8}}\mathrm{Tr}g_{(2)}^2\mathrm{Tr}g_{(2)}{\displaystyle \frac{1}{24}}(\mathrm{Tr}g_{(2)})^3{\displaystyle \frac{1}{6}}\mathrm{Tr}g_{(2)}^3+{\displaystyle \frac{1}{2}}\mathrm{Tr}(g_{(2)}g_{(4)})]g_{(0)ij}).`$
For even $`n=d`$ the first equation in (2.7) determines the coefficients $`h_{(d)}`$. They are given by
$`h_{(2)ij}=0`$ (A.9)
$`h_{(4)ij}={\displaystyle \frac{1}{2}}g_{(2)ij}^2{\displaystyle \frac{1}{8}}g_{(0)ij}\mathrm{Tr}g_{(2)}^2+{\displaystyle \frac{1}{8}}(^k_ig_{(2)jk}+^k_jg_{(2)ik}^2g_{(2)ij}_i_j\mathrm{Tr}g_{(2)})`$ (A.10)
$`={\displaystyle \frac{1}{8}}R_{ikjl}R^{kl}+{\displaystyle \frac{1}{48}}_i_jR{\displaystyle \frac{1}{16}}^2R_{ij}{\displaystyle \frac{1}{24}}RR_{ij}+({\displaystyle \frac{1}{96}}^2R+{\displaystyle \frac{1}{96}}R^2{\displaystyle \frac{1}{32}}R_{kl}R^{kl})g_{(0)ij}`$
$`h_{(6)ij}={\displaystyle \frac{2}{3}}(g_{(4)}g_{(2)}+g_{(2)}g_{(4)})_{ij}{\displaystyle \frac{1}{3}}g_{(2)ij}^3{\displaystyle \frac{1}{6}}g_{(4)ij}\mathrm{Tr}g_{(2)}+{\displaystyle \frac{1}{6}}g_{(0)ij}(3\mathrm{T}\mathrm{r}g_{(6)}3\mathrm{T}\mathrm{r}g_{(2)}g_{(4)}+\mathrm{Tr}g_{(2)}^3)`$
$`{\displaystyle \frac{1}{12}}[{\displaystyle \frac{1}{4}}_i_j\mathrm{Tr}g_{(2)}^2^k_ig_{(4)jk}^k_jg_{(4)ik}+^2g_{(4)ij}`$
$`+g_{(2)}^{kl}[_l_ig_{(2)jk}+_l_jg_{(2)ik}_l_kg_{(2)ij}]`$
$`+{\displaystyle \frac{1}{2}}^k\mathrm{Tr}g_{(2)}(_ig_{(2)jk}+_jg_{(2)ik}_kg_{(2)ij})`$
$`+{\displaystyle \frac{1}{2}}_ig_{(2)kl}_jg_{(2)}^{kl}+_kg_{(2)il}^lg_{(2)j}{}_{}{}^{k}_kg_{(2)il}^kg_{(2)j}{}_{}{}^{l}].`$ (A.11)
## B Divergences in terms of the induced metric
In this appendix we rewrite the divergent terms of the regularized action in terms of the induced metric at $`\rho =ϵ`$. This is needed in order to derive the contribution of the counterterms to the stress energy tensor.
The coefficients $`a_{(n)}`$ of the divergent terms in the regulated action (3.2) are given by
$`a_{(0)}=2(1d),a_{(2)}=b_{(2)}(d)\mathrm{Tr}g_{(2)},`$ (B.1)
$`a_{(4)}=b_{(4)}(d)[(\mathrm{Tr}g_{(2)})^2\mathrm{Tr}g_{(2)}^2],a_{(6)}=\left({\displaystyle \frac{1}{8}}\mathrm{Tr}g_{(2)}^3{\displaystyle \frac{3}{8}}\mathrm{Tr}g_{(2)}\mathrm{Tr}g_{(2)}^2+{\displaystyle \frac{1}{2}}\mathrm{Tr}g_{(2)}^3\mathrm{Tr}g_{(2)}g_{(4)}\right),`$
where $`a_{(6)}`$ is only valid in six dimensions and the numerical coefficients in $`a_{(2)}`$ and $`a_{(4)}`$ are given by
$$b_{(2)}(d2)=\frac{(d4)(d1)}{d2},b_{(2)}(d=2)=1,b_{(4)}(d4)=\frac{d^2+9d16}{4(d4)},b_{(4)}(d=4)=\frac{1}{2}.$$
(B.2)
Notice that the coefficients $`a_{(n)}`$ are proportional to the expression for the conformal anomaly (in terms of $`g_{(n)}`$) in dimension $`d=n`$ .
The counterterms can be rewritten in terms of the induced metric by inverting the relation between $`\gamma `$ and $`g_{(0)}`$ perturbatively in $`ϵ`$. One finds
$`\sqrt{g_{(0)}}`$ $`=`$ $`ϵ^{d/2}\left(1{\displaystyle \frac{1}{2}}ϵ\mathrm{Tr}g_{(0)}^1g_{(2)}+{\displaystyle \frac{1}{8}}ϵ^2[(\mathrm{Tr}g_{(0)}^1g_{(2)})^2+\mathrm{Tr}(g_{(0)}^1g_{(2)})^2]+𝒪(ϵ^3)\right)\sqrt{\gamma }`$
$`\mathrm{Tr}g_{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{2(d1)}}{\displaystyle \frac{1}{ϵ}}\left(R[\gamma ]+{\displaystyle \frac{1}{d2}}(R_{ij}[\gamma ]R^{ij}[\gamma ]{\displaystyle \frac{1}{2(d1)}}R^2[\gamma ])+𝒪(R[\gamma ]^3)\right)`$
$`\mathrm{Tr}g_{(2)}^2`$ $`=`$ $`{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{1}{(d2)^2}}\left(R_{ij}[\gamma ]R^{ij}[\gamma ]+{\displaystyle \frac{3d+4}{4(d1)^2}}R^2[\gamma ]+𝒪(R[\gamma ]^3)\right)`$ (B.3)
The terms cubic in curvatures in (B) give vanishing contribution in (3.4) up to six dimensions.
Putting everything together we obtain that the counterterms, rewritten in terms of the induced metric, are given by
$$S^{\text{ct}}=\frac{1}{16\pi G_\text{N}}_{\rho =ϵ}\sqrt{\gamma }\left[2(1d)+\frac{1}{d2}R\frac{1}{(d4)(d2)^2}(R_{ij}R^{ij}\frac{d}{4(d1)}R^2)\mathrm{log}ϵa_{(d)}+\mathrm{}\right]$$
(B.4)
where all quantities are now in terms of the induced metric, including the one in the logarithmic divergence. These are exactly the counterterms in except that these authors did not include the logarithmic divergence. Equation (B.4) should be understood as containing only divergent counterterms in each dimension. This means that in even dimension $`d=2k`$ one should include only the first $`k`$ counterterms and the logarithmic one. In odd $`d=2k+1`$, only the first $`k+1`$ counterterms should be included. The logarithmic counterterms appear only for $`d`$ even. The counterterms in (B.4) render the renormalized action finite up to $`d=6`$. This covers all cases relevant for the AdS/CFT correspondence. It is straightforward but tedious to compute the necessary counterterms for $`d>6`$. From (B.4) one straightforwardly obtains (3.7).
## C Relation between $`h_{(d)}`$ and the conformal anomaly $`a_{(d)}`$
We show in this appendix that the tensor $`h_{(d)}`$ appearing in expansion of the metric in (2) when $`d`$ is even is a multiple of the stress tensor derived from the action $`a_{(d)}`$. ($`a_{(d)}`$ is, up to a constant, the holographic conformal anomaly).
This can be shown by deriving the stress-energy tensor of the regulated theory at $`\rho =ϵ`$ in two ways and then comparing the results. In the first derivation one starts from (3) and obtains the regulated stress-energy tensor as in (3.6). Expanding $`T_{ij}^{\text{reg}}[\gamma ]`$ in $`ϵ`$ (keeping $`g_{(0)}`$ fixed) we find that there is a logarithmic divergence,
$$T_{ij}^{\text{reg}}[\gamma ;\mathrm{log}]=\frac{1}{8\pi G_\text{N}}\mathrm{log}ϵ(\frac{3}{2}d1)h_{(d)ij}.$$
(C.1)
On the other hand, one can derive $`T_{ij}^{\text{reg}}[\gamma ]`$ starting from (3.2). One has to first rewrite the terms in (3.2) in terms of the induced metric. This is done in the previous appendix. Once $`T_{ij}^{\text{reg}}[\gamma ]`$ has been derived, we expand in $`ϵ`$. We find the following logarithmic divergence
$$T_{ij}^{\text{reg}}[\gamma ;\mathrm{log}]=\frac{1}{8\pi G_\text{N}}\mathrm{log}ϵ((1d)h_{(d)ij}T_{ij}^a,)$$
(C.2)
where $`T_{ij}^a`$ is the stress-energy tensor of the action $`\text{d}^dx\sqrt{detg_{(0)}}a_{(d)}`$. If follows that
$$h_{(d)ij}=\frac{2}{d}T_{ij}^a$$
(C.3)
We have also explicitly verified this relation by brute-force computation in $`d=4`$.
## D Asymptotic solution of the scalar field equation
We give here the first two orders of the solution of the equation (5.4)
$`\varphi _{(2)}={\displaystyle \frac{1}{2(2\mathrm{\Delta }d2)}}\left(\text{ }\text{ }\text{ }_0\varphi _{(0)}+(d\mathrm{\Delta })\varphi _{(0)}\mathrm{Tr}g_{(2)}\right),`$
$`\varphi _{(4)}={\displaystyle \frac{1}{4(2\mathrm{\Delta }d4)}}(\text{ }\text{ }\text{ }_0\varphi _{(2)}2\mathrm{Tr}g_{(2)}\varphi _{(2)}{\displaystyle \frac{1}{2}}(d\mathrm{\Delta })[\mathrm{Tr}g_{(2)}^2\varphi _{(0)}2\mathrm{T}\mathrm{r}g_{(2)}\varphi _{(2)}]`$
$`{\displaystyle \frac{1}{\sqrt{g_{(0)}}}}_\mu (\sqrt{g_{(0)}}g_{(2)}^{\mu \nu }_\nu \varphi _{(0)})+{\displaystyle \frac{1}{2}}^i\mathrm{Tr}g_{(2)}_j\varphi _{(0)}),`$ (D.1)
where in $`\text{ }\text{ }\text{ }\text{ }_0`$ the covariant derivatives are with respect to $`g_{(0)}`$.
If $`2\mathrm{\Delta }d2k=0`$ one needs to introduce a logarithmic term in order for the equations to have a solution, as discussed in the main text. For instance, when $`\mathrm{\Delta }=\frac{1}{2}d+1`$, $`\varphi _{(2)}`$ is undetermined, but instead one obtains for the coefficient of the logarithmic term,
$$\psi _{(2)}=\frac{1}{4}\left(\text{ }\text{ }\text{ }\text{ }_0\varphi _{(0)}+(\frac{d}{2}1)\varphi _{(0)}\mathrm{Tr}g_{(2)}\right).$$
(D.2) |
warning/0002/astro-ph0002369.html | ar5iv | text | # Breaking the degeneracy of cosmological parameters in galaxy redshift surveys
## 1 Introduction
Galaxy redshift surveys are undoubtedly extremely valuable tools to investigate the evolution of the universe at large scales. The cosmologist’s prerogative is to determine the evolution of the matter density contrast $`\delta `$ and peculiar velocity $`\stackrel{}{v}`$ that yields such cosmic structure, customarily assuming that it formed solely by gravity, and the cosmological parameters that determine their dynamics. In the standard paradigm of a FRW expanding universe, the interplay of both fields is governed by the density parameter $`\mathrm{\Omega }_0`$ and the Hubble parameter $`H_0`$. On the other hand, a relationship between the fields $`\delta `$,$`\stackrel{}{v}`$ and the survey data is established by adopting a bias model that purports the correlation between the $`z`$-space galaxy number-count and the underlying matter field. Devoid of such a relationship, the edifice of measuring cosmological parameters from galaxy redshift surveys has no foundation whatsoever. A standard working hypothesis, that I shall accept throughout this paper, is that of linear bias, i.e. $`b^2P(k)_{\mathrm{gals}}/P(k)_{\mathrm{matter}}`$ (more elaborate bias models are propounded in e.g. Dekel & Lahav 1999). For simplicity we shall also leave out the scale-dependence of $`b`$. Therefore, three relevant parameters that are interesting to pin down from redshift-space samples are in this context $`\mathrm{\Omega }_\mathrm{m}`$,$`H_0`$ and $`b`$. In this paper I shall be chiefly concerned with $`\mathrm{\Omega }_\mathrm{m}`$ and $`b`$ ($`H_0`$ will be scaled out with distance).
Tracing back in time the matter fields takes us to an initial epoch of fluctuations of very small amplitude $`\delta \begin{array}{c}<\hfill \\ \hfill \end{array}10^4`$, seeded by a period of inflationary expansion. At that point the information derived from the galaxy surveys connects with early-universe data such as the spectrum of fluctuations on the CMB. If the matter fields could realistically be traced back to such a primordial stage by integrating the equations of gravitational instability, then the statistics of the $`\delta `$ field would be a potentially key discriminant to rule out cosmological models. For instance, non-gaussianity in the initial $`\delta `$ field rules out most inflationary models, and only those leading to a non-Gaussian primordial spectrum remain acceptable (such models are suggested in e.g. Linde, Sasaki & Tanaka 1999).
Kaiser (1987) proposed measuring cosmological parameters from redshift-space distortions by virtue of the fact that overdense regions appear to be flatter along the line-of-sight in redshift space. This distortion, quantified by the parameter $`\beta =\mathrm{\Omega }_\mathrm{m}^{0.6}/b`$, permits us to solve the equations for $`\delta `$,$`\stackrel{}{v}`$, at least perturbatively (see e.g. Dekel 1994; Coles & Sahni 1995), and measurements of $`\beta `$ have been investigated in much detail in the literature (Strauss & Willick 1995; Dekel 1994,1999a; Dekel, Burstein & White 1997). Also, in view of the fact that the bias parameter is almost certainly dependent on the selected sample, estimates have been computed for $`\beta _{IRAS}`$ given $`b_I`$ for IRAS galaxies (Dekel et al. 1993; Fisher et al. 1995a; Willick et al. 1997a,b; Sigad et al. 1998; more recently from the PSC$`z`$ sample, Canavezes et al. 1998; Tadros et al. 1999; Saunders et al. 2000) and from the Optical Redshift Survey (ORS) (Hudson et al. 1995; Santiago et al. 1995; Baker et al. 1998). The Mark III peculiar velocity survey similarly yields estimates of $`\beta `$ from redshift distortions (Willick et al. 1995,1996,1997a,b; Dekel, Burstein & White 1997; Sigad et al. 1998). It is only beyond the linear approximation (i.e. $`\delta \stackrel{}{v}`$) and, indeed, beyond the assumption of linear bias, that one can break down the degeneracy between $`\mathrm{\Omega }_\mathrm{m}`$ and $`b`$ and estimate these parameters separately, rather than via $`\beta `$ (Fry 1994; Bernardeau et al. 1995). Verde et al. (1998) achieved this by proposing the bispectrum as a measure of cosmological parameters, in a model of non-linear bias. In this paper we also pursue breaking the degeneracy of $`\mathrm{\Omega }_\mathrm{m}`$ and $`b`$ from the redshift-space data and show that by using the least-action framework it is indeed possible to do so within the linear bias model.
The least-action principle (LAP) was first used in the Local Group by Peebles (1989,1990). The trajectories of nearby galaxies were computed subject to two boundary conditions: vanishing initial velocities and fixed present positions. This simple scenario of self-gravitating point-like masses with two boundary conditions produced an estimate of $`\mathrm{\Omega }_\mathrm{m}`$ by fitting to the observations the predicted peculiar velocities of nearby galaxies. The LAP method has also been used as a test of $`\mathrm{\Omega }_\mathrm{m}=1`$ CDM models (Branchini & Carlberg 1994), as well as to integrate the orbits of a significant number of galaxies from partial coverage redshift samples (e.g. Shaya, Peebles & Tully 1995). An equivalent representation of the LAP method in terms of continuous fields, i.e. the density contrast and velocity fields was proposed by Giavalisco et al. (1993), and employed in Susperregi & Binney (1994)(hereafter SB94) and Susperregi (1995) in the reconstruction of $`\mathrm{\Omega }_\mathrm{m}=1`$ simple models, such as exact solutions and Gaussian random fields. More recently, Schmoldt & Saha (1998) proposed a variant of the customary LAP formulation by rewriting the equations motion in redshift space.
The key difference between the variational and perturbative approaches lies on how the errors are spread over the time-reversed evolution. This is qualitatively sketched in Fig. 1. A $`n`$th-order solution differs, in the time-reversed direction, from the true solution by a monotonically growing parameter $`ϵ`$ which sets out from a small value $`ϵ(t_0)`$ (at any rate $`ϵ_0`$ is at least the sum of the systematic and random errors of the dataset) and the conservation of kinematical quantities is preserved up to $`O(ϵ^n)`$. This is adequate within a time span $`t_ct\begin{array}{c}<\hfill \\ \hfill \end{array}t_0`$ where $`ϵ(t_c)1`$, and $`t_c`$ marks a transition into the loss of convergence. The distribution of errors in the LAP method on the other hand, is by construction evenly distributed along the trajectory; the initial and final boundary conditions are fixed, though not without systematic and numerical errors, and the parameter $`ϵ`$ fluctuates along the trajectory between both end-points (Fig. 1b). Hence the solution is well-behaved whether the errors remain within the bound $`ϵ\begin{array}{c}<\hfill \\ \hfill \end{array}1`$ or not. In that respect there is an advantage with respect to perturbative solutions; the downside of it is of course that within the span of time where perturbative solutions are valid, LAP errors may fluctuate with larger amplitude than the perturbative equivalent. The LAP method, in a nutshell, thus consists in finding Ansatze for the matter fields that optimize the distribution of $`ϵ`$ along the phase-space trajectory, and hence minimize the overall departure with respect to the exact solution. The following two difficulties may arise:
* A Finding “dynamically plausible” solutions. If the matter field is sparsely sampled or the errors in the dataset are substantial, then the boundary condition given by the survey, taken at face value, may not correspond to the outcome of gravitational evolution from the initial fluctuations (typically $`\delta 0`$ or vanishing peculiar velocities). The LAP method will in this case find a dynamically plausible fit between the end-points, which will be as faithful a representation of the true evolution as is the quality of the dataset.
* B Formation of multistreams in over-dense regions. Multistreams are characterised by galaxies at the same redshift which are located at different positions along the line of sight and have different infalling velocities. The degeneracy in redshift among streams makes them indistinguishable and hence compatible but inequivalent solutions result, as many as there are streams. The LAP method cannot discriminate among these solutions; multistreams indeed erase the memory of their past evolution.
The second problem can only be overcome by casting aside part of the information contained in the sample and smoothing over the existing non-linearities to transform the multivalued field into a single-valued one, typically with a smoothing length $`5001000`$ $`\mathrm{km}\mathrm{s}^1`$. The resulting smoothed field is clearly a less resolved representation of the underlying galaxy orbits, albeit the only tractable one.
The advent of large galaxy redshift surveys strengthens the motivation to use the LAP method. Near all-sky redshift surveys, e.g. PSC$`z`$, IRAS 1.2 Jy and ORS provide an excellent sky coverage (within a galactic latitude $`|b|\begin{array}{c}>\hfill \\ \hfill \end{array}8^{}`$ for IRAS galaxies and $`|b|\begin{array}{c}>\hfill \\ \hfill \end{array}20^{}`$ for ORS), that may be extended further to cover the Zone of Avoidance via a Wiener reconstruction (Fisher et al. 1995b; Zaroubi et al. 1999). They are therefore a fairly thorough representation of the underlying matter density field. Obviously the greater number of galaxies in the sample the more accurate is the representation of the field, and this is best achieved with a redshift survey. Real-space datasets require Tully-Fisher distance calibrations of individual galaxies, and consequently the end result is a sparser sampling than is achieved with the same computational effort by measuring redshifts and angular coordinates. The goal of this paper is to exploit galaxy redshift surveys to the best effect and extract as much information from them as is possible; the main thesis put forward is the LAP method demonstrably breaks down the degeneracy in the determination of $`\mathrm{\Omega }_\mathrm{m}`$ and $`b`$. This entails very tangible advantages. On the one hand, the freedom to investigate those two parameters separately permits us not to take the idea of bias seriously. A form of bias will certainly always be present in one form or another so that we can make sense of the galaxy number-count with respect to the underlying matter field. However, whether that is a linear or non-linear bias, the more one dissociates this phenomenological relationship from our measurements of $`\mathrm{\Omega }_\mathrm{m}`$, the more credible those measurements will be. This is indeed what LAP does. On the other hand, the LAP method produces a reconstruction on the basis of the redshift-space sample alone, free of any proviso regarding the shape of the power spectrum. Assuming a given shape for $`P(k)`$ unduly overconstrains the system, as will the addition of other datasets.
In this article, I shall mainly apply the LAP method to the IRAS 1.2 Jy survey and study the predicted values of $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$. The reconstructed IRAS 1.2 Jy velocity field is then compared with the Mark III velocity sample to seek a fine-tuning of the parameters. A more thorough undertaking, in terms of the quality of the sample, is to apply the LAP method to PSC$`z`$, which is by a factor of 3 a more densely sampled survey than IRAS 1.2 Jy, and it will be interesting to tackle this in future work. The article is structured as follows: Section 2 describes the LAP method in some detail and how to find solutions that are consistent with a redshift-space dataset; in Section 3 we test the method with several IRAS mock catalogues obtained via $`n`$-body simulations; in Section 4 we apply the method to the IRAS 1.2 Jy galaxy redshift survey, optimizing the predicted velocities with the Mark III dataset; finally, in Section 5 we summarize the main conclusions.
## 2 The LAP method
### 2.1 Redshift-space coordinates
The redshift coordinates of galaxies are defined
$$\stackrel{}{z}=H_0\stackrel{}{r}+\widehat{r}(\widehat{r}\stackrel{}{v}),$$
(1)
where $`\stackrel{}{r}(r,\theta ,\phi )`$ is the physical position, $`H_0`$ is the present value of the Hubble parameter, $`\stackrel{}{v}`$ the peculiar velocity, and $`\widehat{r}`$ a unit vector in the radial (line-of-sight) direction. $`\stackrel{}{z}`$ has units of velocity; its radial component is the redshift $`z_r=cz`$, and the angular components are the same in both x-space and z-space, up to the distance scale. Henceforth we shall measure distances in $`\mathrm{km}\mathrm{s}^1`$, hence $`H_0`$ is scaled out of the equations. In comoving coordinates, (1) reads
$$\stackrel{}{s}=\stackrel{}{x}+\widehat{x}(\widehat{x}_x\alpha ),$$
(2)
where the scale factor of the universe is normalized to $`a(t_0)=1`$; $`\alpha (t,\stackrel{}{x})`$ is the velocity potential, $`\stackrel{}{v}a^1\alpha `$. Hereafter we adopt $`t_0=1`$.
### 2.2 Dynamics
The cosmological perturbations are derived from the action
$$𝒮=_0^1dt_{\mathrm{sample}}d\stackrel{}{x},$$
(3)
where $``$ is given by
$$=\frac{1}{2}(1+\delta )\stackrel{}{v}^2+\alpha \xi \varphi \delta \frac{|\varphi |^2}{3\mathrm{\Omega }_\mathrm{m}a^2};$$
(4)
$`\delta `$ is the density contrast and $`\varphi `$ the gravitational potential caused by the perturbations and
$$\xi \dot{\delta }+\frac{1}{a}[(1+\delta )\stackrel{}{v}]$$
(5)
is the excess flux. The variations $`\delta 𝒮/\delta v_i=\delta 𝒮/\delta \varphi =0`$ yield
$$\stackrel{}{v}=\frac{1}{a}\alpha ,$$
(6)
$$^2\varphi =\frac{3}{2}a^2\mathrm{\Omega }_\mathrm{m}\delta .$$
(7)
Similarly, $`\delta 𝒮/\delta \delta =\delta 𝒮/\delta \alpha =0`$ yield respectively
$$\xi =0,$$
(8)
$$\dot{\alpha }+\frac{|\alpha |^2}{2a^2}+\varphi =0,$$
(9)
where we have eliminated $`\stackrel{}{v}`$ via (6) and we do not consider $`\mathrm{\Omega }_\mathrm{\Lambda }`$. The field equations (8),(9) are subject to the following boundary conditions:
* Homogeneity of the density field at $`t0`$. Density perturbations grow from initial fluctuations of negligible amplitude:
$$\delta (t0,\stackrel{}{x})0.$$
(10)
* Galaxy redshift survey at the present time. The galaxy number-count density $`\rho _s`$ in $`z`$-space constrains the real fields $`\delta (\stackrel{}{x})`$ and $`\alpha (\stackrel{}{x})`$ via
$$\rho _s(\stackrel{}{s})=x^2\frac{N_{\mathrm{gals}}}{V}\left(\frac{1+b\delta }{1+\alpha ^{\prime \prime }}\right),$$
(11)
where the tilde denotes derivation along the radial direction, $`x`$ is the radial comoving distance and $`b`$ is the bias parameter. Condition $`(I)`$ is motivated by the CMB Sachs-Wolfe constraint $`\delta \begin{array}{c}<\hfill \\ \hfill \end{array}10^4`$ over $`r100,000`$ $`\mathrm{km}\mathrm{s}^1`$, so we accept that perturbations are negligible in the limit $`t0`$. A proof for $`(II)`$ is given in Appendix A. In order to solve (8),(9), we construct the trial fields:
$$\delta =\underset{n=0}{\overset{N}{}}f_n(t)\delta _n(\stackrel{}{x}),$$
(12)
$$\alpha =\underset{n=0}{\overset{N}{}}g_n(t)\alpha _n(\stackrel{}{x}),$$
(13)
where the basis functions $`f_n`$,$`g_n`$ are adjusted to numerical convenience. SB94 considered $`f_nD(D1)^n`$, and $`g_n=(\dot{D}/D)f_n`$, where $`D`$ is the linear growth factor, normalized to unity at $`t=1`$, so that the lowest-order series (12),(13) are identical to the perturbative solutions. This is however strictly speaking not a compelling choice, and a sensible choice of orthogonal polynomials leads to an Ansatz of better convergence. As we have discussed in the Introduction (point A), the sparseness of the dataset obscures the dynamical evolution and the LAP method is reduced to a numerical fit of the fields to the truncated equations, that we derive below, subject to (10),(11). In trying to approximate a function $`f(t)`$ by orthogonal polynomials $`P_m(t)`$ in $`0\begin{array}{c}<\hfill \\ \hfill \end{array}t\begin{array}{c}<\hfill \\ \hfill \end{array}1`$, a weight function $`w(t)0`$ tells us the relative importance of the errors spread over the domain. For a uniform $`w`$, $`f_n`$ are the \[spherical\] Legendre polynomials $`L_m(t)`$, whereas for a weight function that is larger at the endpoints (10),(11) than throughout the trajectory, e.g. $`w(t)=(1t^2)^{1/2}`$ (by shifting the domain from $`[0,1]`$ to $`[1,1]`$), the optimal choice are in this case Chebyshev polynomials $`T_n(t)`$. This choice minimizes the errors around the endpoints and it gives a greater weight to the solutions (matching the boundary conditions) in this region. In the analysis that follows, we shall adopt $`f_n=T_n`$ and $`g_n=a^2f_n`$. The fields $`\delta _n`$,$`\alpha _n`$ are expanded in terms of spherical harmonics,
$$\delta _n=\underset{rlm}{}\delta _{rlm}^{(n)}j_l(k_rx)Y_{lm},$$
(14)
$$\alpha _n=\underset{rlm}{}\alpha _{rlm}^{(n)}j_l(k_rx)Y_{lm},$$
(15)
where $`j_l`$ are spherical Bessel functions. Substituting (12),(13) into (6),(7) we get
$$\stackrel{}{v}=a\underset{rlmn}{}\left[\widehat{x}\alpha _{}^{}{}_{rlm}{}^{(n)}j_l(k_rx)+\frac{1}{x}(\widehat{x}\stackrel{}{J}_{lm}^{(n)})\right]T_nY_{lm},$$
(16)
$$\varphi =\frac{3}{2}a^2\mathrm{\Omega }_\mathrm{m}\underset{rlmn}{}\frac{\delta _{rlm}^{(n)}}{k_r^2}T_nj_l(k_rx)Y_{lm};$$
(17)
the coefficients $`\alpha _{}^{}{}_{rlm}{}^{(n)}`$ and $`\stackrel{}{J}_{lm}^{(n)}`$ are given in Appendix B. The boundary conditions (10),(11) then read
$$0=\underset{n=0}{\overset{N}{}}(1)^n\delta _n,$$
(18)
$$\rho _s=x^2\left(\frac{N_{\mathrm{gals}}}{V}\right)\left[1+b\delta (1,\stackrel{}{x})\right]\left[1+\alpha ^{\prime \prime }(1,\stackrel{}{x})\right]^1,$$
(19)
where $`t`$ is rescaled to the interval $`[1,1]`$ for convenience in using $`T_n`$, and in (18) we have used $`T_n(1)=(1)^n`$. The choice of basis functions of SB94 satisfy (18) by construction, and in our choice of basis functions the constraint is less trivial, but still it is easily tackled numerically. If we restrict ourselves to the interval $`0t1`$, then (18) evaluated at $`t=0`$ eliminates all the Chebyshev polynomials of odd order. This is an equivalent approach but we shall adopt the convention above, $`1t1`$. The constraint (19) is the core of the problem as it is where all the information of the dataset is contained. The remainder of the paper will focus on the different ways one can use that constraint.
### 2.3 Finding LAP solutions
Substituting (12)–(15) into equations (8),(9), we get
$$\underset{n=0}{\overset{N}{}}\underset{rlm}{}\left[\dot{T}_n\delta _{rlm}^{(n)}k_r^2T_n\alpha _{rlm}^{(n)}\right]j_l(k_rx)Y_{lm}$$
$$=\underset{p,q=0}{\overset{N}{}}\underset{\genfrac{}{}{0pt}{}{rlm}{r^{}l^{}m^{}}}{}T_pT_q\{\alpha _{}^{}{}_{rlm}{}^{(q)}\delta _{}^{}{}_{r^{}l^{}m^{}}{}^{(p)}j_l(k_rx)j_l^{}(k_r^{}x)$$
(20)
$$+\frac{1}{x^2}[\widehat{x}\stackrel{}{J}_{lm}^{(p)}(\delta )][\widehat{x}\stackrel{}{J}_{l^{}m^{}}^{(q)}(\alpha )]\}Y_{lm}Y_{l^{}m^{}},$$
and
$$\underset{n=0}{\overset{N}{}}\underset{rlm}{}\left[\frac{3}{2}\mathrm{\Omega }_\mathrm{m}k_r^2\delta _{rlm}^{(n)}+\left(\frac{\dot{T}_n}{T_n}+2\frac{\dot{a}}{a}\right)\alpha _{rlm}^{(n)}\right]T_nj_l(k_rx)Y_{lm}$$
$$=\frac{1}{2}\underset{p,q=0}{\overset{N}{}}\underset{\genfrac{}{}{0pt}{}{rlm}{r^{}l^{}m^{}}}{}T_pT_q\{\alpha _{}^{}{}_{rlm}{}^{(q)}\alpha _{}^{}{}_{r^{}l^{}m^{}}{}^{(p)}j_l(k_rx)j_l^{}(k_r^{}x)$$
(21)
$$+\frac{1}{x^2}[\widehat{x}\stackrel{}{J}_{lm}^{(p)}(\alpha )][\widehat{x}\stackrel{}{J}_{l^{}m^{}}^{(q)}(\alpha )]\}Y_{lm}Y_{l^{}m^{}},$$
where the coefficients $`\stackrel{}{J}_{lm}^{(p)}(\delta )`$,$`\stackrel{}{J}_{lm}^{(q)}(\alpha )`$ are defined as in (43) in Appendix B and $`\delta _{}^{}{}_{rlm}{}^{(n)}`$ as in (41) via the trivial substitution $`\alpha \delta `$. By multiplying (20),(21) by $`T_rj_lY_{lm}`$ and integrating over all coordinates, we get
$$\underset{n=0}{\overset{N}{}}T_r\dot{T}_nC_y^\delta \delta _y^{(n)}+\underset{n=0}{\overset{N}{}}T_rT_nC_y^\alpha \alpha _y^{(n)}$$
$$=\underset{p,q=0}{\overset{N}{}}T_rT_pT_q\underset{y^{}y^{\prime \prime }}{}D_{y^{}y^{\prime \prime }}^y\delta _y^{}^{(p)}\alpha _{y^{\prime \prime }}^{(q)},$$
(22)
$$\underset{n=0}{\overset{N}{}}\mathrm{\Omega }_\mathrm{m}T_rT_nS_y^\delta \delta _y^{(n)}+\underset{n=0}{\overset{N}{}}T_r(\dot{T}_n+2\frac{\dot{a}}{a}T_n)S_y^\alpha \alpha _y^{(n)}$$
$$=\underset{p,q=0}{\overset{N}{}}T_rT_pT_q\underset{y^{}y^{\prime \prime }}{}E_{y^{}y^{\prime \prime }}^y\alpha _y^{}^{(p)}\alpha _{y^{\prime \prime }}^{(q)},$$
(23)
where $`y(rlm)`$ and the angle brackets $``$ for the Chebyshev polynomials are defined in Appendix C. In deriving (22),(23), the coefficients $`C_y^\delta `$, $`C_y^\alpha `$, $`S_y^\delta `$, $`S_y^\alpha `$, $`D_{y^{}y^{\prime \prime }}^y`$ and $`E_{y^{}y^{\prime \prime }}^y`$ are calculated via Clebsch-Gordan coefficients for cross-products of $`Y_{lm}`$ and via the standard orthogonality relations for $`Y_{lm}`$ and $`j_l`$, given in Appendix D. Cross-products of $`j_l`$ terms are estimated numerically. We proceed to solve (22),(23) numerically with the following iterative procedure. We first construct an Ansatz of the coefficients $`\delta _y^{(n)}`$,$`\alpha _y^{(n)}`$ that satisfies, to linear order, (22),(23) as well as (18),(19). We start out with the galaxy number-count density $`\rho _s`$. Following its definition in Appendix A, this quantity has units of inverse velocity, and we define its associated $`z`$-space density contrast via
$$\rho _s\frac{4\pi N_{\mathrm{gals}}}{s_{\mathrm{max}}}(1+\delta _s),$$
(24)
where $`s_{\mathrm{max}}cz_{\mathrm{max}}`$ is the maximum redshift in the sample. Our first Ansatz entails $`b=1`$ and linear evolution, so that $`\delta _s^2\alpha `$, and on inverting this relation to obtain the coefficients $`\alpha _y^{(n)}`$, we estimate $`\delta (\stackrel{}{x})\delta _s(\stackrel{}{x}+\widehat{x}\alpha ^{})`$ by using the expression for the radial derivatives (41). This yields a first Ansatz for $`\delta _y^{(n)}`$, $`\alpha _y^{(n)}`$, derived from the dataset, that satisfies the linearized equations, given by the LHS of (22),(23):
$$\left[\begin{array}{cc}\mathrm{C}^\alpha & \mathrm{C}^\delta \\ \mathrm{S}^\alpha & \mathrm{S}^\delta \end{array}\right]\left[\begin{array}{c}\stackrel{}{\alpha }_y\\ \stackrel{}{\delta }_y\end{array}\right]0,$$
(25)
where the column vectors are $`(\stackrel{}{\alpha }_y)_r=\alpha _y^{(r)}`$ and $`(\stackrel{}{\delta }_y)_r=\delta _y^{(r)}`$, with $`r=0,\mathrm{},N`$. The solutions of the homogeneous system are then re-adjusted to satisfy (18),(19) and we use these to construct the quadratic terms on the RHS of (22),(23). This leads to an inhomogeneous system that again we solve for $`\stackrel{}{\delta }_y`$,$`\stackrel{}{\alpha }_y`$. On each iteration we improve the solutions by least-squaring them to satisfy (18),(19) to the best accuracy and we are also free to vary the parameters ($`b`$,$`\mathrm{\Omega }_\mathrm{m}`$) for improved convergence. This procedure is very accurate, as we will show in the next sections, and it permits us to improve the estimate of the mapping $`\stackrel{}{x}\stackrel{}{s}`$ at each iteration using the full non-linear relationship (19). At each iteration, the fields $`\stackrel{}{\delta }_y`$,$`\stackrel{}{\alpha }_y`$ are used to obtain an estimate $`\stackrel{~}{\rho }_s(\stackrel{}{s})`$ of the RHS of (19). We then vary these fields to obtain a minimum of the quantity $`_\stackrel{}{s}(\rho _s\stackrel{~}{\rho }_s)^2`$. Therefore we do not perform a $`j_lY_{lm}`$ expansion of the dataset, and it is very convenient not to do so, as a relationship of this kind between the redshift and real-space coordinates entails that we compare them via a Taylor expansion $`j_l(k_rs)j_l(k_rx)+k_r\alpha ^{}j_l^{}`$; an approximation of this kind $`𝒪(^2j_l)`$ introduces an error of up to 15% for $`l\begin{array}{c}>\hfill \\ \hfill \end{array}10`$ as can be shown from (40) in Appendix B.
### 2.4 Normal modes
We have noted that the linearized equations (25) are a homogeneous matrix system. If the determinant of the matrix is non-zero, then the only possible solution is $`\stackrel{}{\delta }_y=0`$ and $`\stackrel{}{\alpha }_y=0`$. We know however that (25) is also valid for linear fields, and these have non-vanishing coefficients. Therefore we conclude that the determinant of the system vanishes. Such a system of equations is tackled through the Singular Value Decomposition (SVD) procedure. It factorises the singular matrix in (25) in a product of three matrices: two orthogonal matrices U and V, and a diagonal one W, which has one or more vanishing weights along the diagonal. After SVD, (25) reads
$$\mathrm{U}\left[\begin{array}{cccc}0& & & \\ & w_1& & \\ & & w_2& \\ & & & \mathrm{}\end{array}\right]\mathrm{V}\left[\begin{array}{c}\stackrel{}{\alpha }_y\\ \stackrel{}{\delta }_y\end{array}\right]=0,$$
(26)
where the weights $`w_1,w_2,\mathrm{}w_N`$ are non-zero real numbers. Therefore, the vector
$$\stackrel{}{N}_y=\mathrm{V}\left[\begin{array}{c}\stackrel{}{\alpha }_y\\ \stackrel{}{\delta }_y\end{array}\right]$$
(27)
gives a coordinate basis on which the first component, the normal mode, is unconstrained by the system (25). $`N_y^{(0)}`$ is solely determined by (18),(19). The rest of the components of $`\stackrel{}{N}_y`$ (which are identically zero for linear fields) are functions of the normal mode. Therefore, one can rewrite the full non-linear system (22),(23) in terms of the fields $`\stackrel{}{N}_y`$ and this would be strictly speaking the natural basis to investigate the underlying mode coupling induced by gravity. In the Fourier formulation with a set of basis functions like those used in SB94, $`f_n=D(D1)^n`$, it is easy to show numerically that the $`\stackrel{}{k}`$-th normal mode is given by
$$N_\stackrel{}{k}^{(0)}=\delta _\stackrel{}{k}^{(0)}+k^2\alpha _\stackrel{}{k}^{(0)}.$$
(28)
This has a simple physical interpretation: (28) is a vanishing scalar for linear fields and thus its departure from zero gives us a measure of non-linearity. This quantity is determined by the boundary conditions. In the spherical harmonic formulation, the normal modes (equivalent to (28)) are
$$N_y^{(0)}=\underset{n=0}{\overset{N}{}}h_n(\delta _y^{(n)}k_r^2\alpha _y^{(n)}),$$
(29)
where
$$h_n=\frac{\eta }{\pi }_0^1dtw(t)D(t)T_n$$
(30)
where $`\eta =1`$ for $`n=0`$ and $`\eta =2`$ otherwise. The quantity (29) vanishes in the linear regime and, like (28), its departure from zero is a measure of non-linearity.
### 2.5 Using the method in practice
The apparent mathematical complexity of the LAP method has precluded its wider use in practice. The fraction of papers in the literature that employ LAP techniques to investigate large-scale structure is minute in contrast to analyses based on perturbation theory techniques, such as POTENT, VELMOD and others. The latter unquestionably have the virtue of simplicity, and are as efficient as they are easy to implement. However, in practice the method described in this section entails no more complexity than programming an $`n`$-body code; an undertaking that merits the effort, so as to estimate $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$, rather than merely $`\beta `$. The chief difficulty resides in writing an algorithm for an effective numerical resolution of (22),(23). This may be a somewhat arduous task, but at any rate a very straightforward one with a very basic grasp of numerical methods.
The LAP method is very flexible in its implementation. The basic input in the problem are the boundary conditions (10), (11) and the procedure that is to be followed to find a stationary action linking both end-points is largely a matter of numerical convenience. The algorithm used in this section employs Chebyshev polynomials to fit the trial fields $`\delta `$ and $`\alpha `$ to the dynamics. A myriad of other choices (e.g. binomial expansions, Legendre and Hermite polynomials, etc) is also feasible and thus the LAP implementation set out above is by no means a straightjacket recipe (for a more condensed presentation of the algorithm, see Susperregi 2000).
In short, the algorithm can be summarized as follows.
* A galaxy redshift survey is a dataset $`𝒟`$ of points ($`z`$,$`\phi `$,$`\theta `$). Those raw data are transformed to a smoothed redshift-space field $`\rho _s(\stackrel{}{s})`$, given a smoothing length and a window function $`W(k)`$. In this article we shall exclusively implement Gaussian smoothing.
* The name of the game is to compute a fit for $`\delta `$,$`\alpha `$. The starting point is to make a linear Ansatz that is consistent with $`\delta _s`$, which is derived from (24). This is achieved by inverting the relation $`\delta _s^2\alpha `$ and next estimating $`\delta \delta _s(\stackrel{}{x}+\widehat{x}\alpha ^{})`$.
* The linear Ansatz is the first input to be used in equations (22),(23). These yield the homogeneous system (25), which is our second port of call. The solutions $`\delta _y`$,$`\alpha _y`$ obtained are least-square fitted to (18),(19). This requires adopting a value of $`b`$.
* The adjusted values of $`\delta _y`$,$`\alpha _y`$ are brought back to construct the RHS of (22),(23), and from there one computes the new $`\delta _y`$,$`\alpha _y`$ in the LHS of (22),(23). This part of the operation entails an assumed value for $`\mathrm{\Omega }_\mathrm{m}`$. In the normal mode coordinates discussed in 2.4, the modes $`\delta _y`$,$`\alpha _y`$ of the cosmic fields are merely excitation modes of a harmonic oscillator and the terms in the RHS of (22),(23) represent nonlinear perturbations of those excitation modes.
* Successive iterations of the procedure eventually yield the correct values of $`\delta _y`$,$`\alpha _y`$. The values of $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$ are readjusted in the process and their estimated values are those that result in the most rapid convergence of the solutions.
The algorithm thus produces the cosmic fields and an estimate of the cosmological parameters. In the remainder of the article we shall investigate how to make the best use of the procedure and how to quantify the relative likelihood of different values of the cosmological parameters.
## 3 Test of the method
We test the LAP method on mock catalogues derived from $`n`$-body simulations, using a Gaussian smoothing length of 600 $`\mathrm{km}\mathrm{s}^1`$. The IRAS 1.2 Jy power spectrum (Fisher et al. 1993) is adopted as a prior, and the simulations are performed over a periodic box $`L=25,600`$ $`\mathrm{km}\mathrm{s}^1`$ with $`128^3`$ grid-points and $`128^3`$ particles. The simulations are performed from Gaussian initial conditions, for the following values of the parameters: $`b=0.8,1.0,1.2`$ and $`\mathrm{\Omega }_\mathrm{m}=0.3,0.6,1.0`$. The fields are evolved forward in time until $`\sigma _80.7`$ over $`800`$ $`\mathrm{km}\mathrm{s}^1`$, using a Gaussian cutoff. We choose a two-powerlaw functional form of selection function (Yahil et al. 1991):
$$\varphi (rr_s)=\left(\frac{r_s}{r}\right)^{2\alpha }\left(\frac{r_{}^2+r_s^2}{r_{}^2+r^2}\right)^\beta ,$$
(31)
and $`\varphi (rr_s)=1`$, where $`r_s=500`$ $`\mathrm{km}\mathrm{s}^1`$, $`r_{}=5034`$ $`\mathrm{km}\mathrm{s}^1`$, $`\alpha =0.483`$ and $`\beta =1.79`$ (Fisher et al. (1995a); we adopt the estimated central values of these parameters and will not test the fine detail of the variations of $`\varphi (r)`$ due to their errors), and thus we compute the redshift-space dataset over a sphere of radius $`x_{\mathrm{max}}17,000`$ $`\mathrm{km}\mathrm{s}^1`$. The resulting mock catalogue has an effective radius of $`13,000`$ $`\mathrm{km}\mathrm{s}^1`$ beyond which the galaxy number-count is sparse and is cut off for the purpose of the reconstruction. The number of realizations are nine in total, and we denote $`d(b,\mathrm{\Omega }_\mathrm{m})`$ the $`z`$-space mock samples derived in this way. Each dataset $`d(b,\mathrm{\Omega }_\mathrm{m})`$ results from a unique pair of real-space fields $`\delta `$,$`\alpha `$ which are the density contrast and velocity potential that we obtain via the $`n`$-body simulations.
The tests are carried out by using $`d(b,\mathrm{\Omega }_\mathrm{m})`$ as an input dataset in (19) without any prior assumption on the real values of the parameters of the mock sample. We use (19) to solve (22),(23) following the iterative procedure given in §2.3 and derive the estimated fields $`\stackrel{~}{\delta }`$,$`\stackrel{~}{\alpha }`$ for different values of the parameters $`\stackrel{~}{b}`$,$`\stackrel{~}{\mathrm{\Omega }}_\mathrm{m}`$. The likelihood of these parameters is estimated on the basis of the performance of the solutions $`\stackrel{~}{\delta }`$,$`\stackrel{~}{\alpha }`$ in terms of their convergence and ability to satisfy the constraints. We use a likelihood function given by the inverse of the $`\chi `$-squared sum of the differences of the fields between successive iterations in solving (22),(23), i.e.
$$\lambda (b,\mathrm{\Omega }_\mathrm{m})\left[\underset{x}{}\left(\frac{\delta _n\delta _{n1}}{\sigma _\delta }\right)^2+\left(\frac{\alpha _n\alpha _{n1}}{\sigma _\alpha }\right)^2\right]^1,$$
(32)
where $`n25`$, $`\sigma _\delta 0.20`$, $`\sigma _\alpha `$ is a normalization factor that rescales the coefficients $`\sigma _n`$ so that $`\alpha `$ becomes a dimensionless field within the range $`1\begin{array}{c}<\hfill \\ \hfill \end{array}\alpha \begin{array}{c}<\hfill \\ \hfill \end{array}1`$ and we have used $`N=10`$ and $`l15`$ and an initial linear Ansatz. The results are shown in Fig. 2. The likelihood contours are the LAP reconstructions and the crosses on all nine panels of Fig. 2 indicate the values of the real parameters in each mock dataset on the $`(b,\mathrm{\Omega }_\mathrm{m})`$ plane. As can be observed, the likelihood contours are certainly well correlated with the location of the crosses, where the innermost contours mark the level of 95% likelihood, that in all cases lie in the neighbourhood of the real values of the parameters. The likelihood contours show an approximately elliptical shape, with the major semiaxes tilted at approximately 45 degrees, suggesting a correlation between both parameters that merely arises in the numerical computation. The estimates in the reconstruction are fairly good, with a trend in underestimating slightly the values of both parameters. The best reconstructions are for the intermediate value of the density parameter $`\mathrm{\Omega }_\mathrm{m}=0.6`$, shown in the second row. In this case the crosses actually lie within the highest likelihood contours, with very little scatter. Overall, in the nine reconstructions the rms scatter in $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$ lie within the region $`0.26\begin{array}{c}<\hfill \\ \hfill \end{array}\sigma _\mathrm{\Omega }^2\begin{array}{c}<\hfill \\ \hfill \end{array}0.44`$, $`0.15\begin{array}{c}<\hfill \\ \hfill \end{array}\sigma _b^2\begin{array}{c}<\hfill \\ \hfill \end{array}0.32`$. The largest scatter in $`\mathrm{\Omega }_\mathrm{m}`$ is for $`\mathrm{\Omega }_\mathrm{m}=0.6`$, and a similar situation arises with $`b`$, whereby the intermediate value $`b=1.0`$ has the larger error.
The effect of underestimating the true values of the parameters is systematic and can be calibrated. This effect can be largely ascribed to the unconventional choice of likelihood estimator (32). One could argue that, for slowly varying variances, $`\lambda b^2`$ (chiefly from the $`\delta `$ part of the RHS of (32)) and therefore smaller values of the bias factor (and consequently, by correlation, also of $`\mathrm{\Omega }_\mathrm{m}`$) are favoured. However, it is not straightforward to disentangle the dependence of the solutions on the parameters after successive iterations. The likelihood estimator used is thus to some extent biased. However, we find that the criterion of convergence given by the RHS of (32), suitably normalised, is the sharpest discriminator to pin down the best estimates of the cosmological parameters. We have carried out numerous tests with more conventional likelihood estimators (e.g. Fisher likelihood matrix, etc) obtaining much poorer results than with (32).
Fig. 3 shows the density constrast reconstructions for the same datasets $`d(b,\mathrm{\Omega }_\mathrm{m})`$. The reconstructed density contrast $`\delta _{\mathrm{LAP}}`$ is shown on the horizontal axis plotted point-by-point within the selected spherical volume ($`r13,000`$ $`\mathrm{km}\mathrm{s}^1`$) against the real density contrast of the mock datasets. A solid line of slope 1.0 is plotted across each panel that does not correspond to the regression line on each panel though the differences are tiny. The slopes of the regression lines lie within the range $`0.99\pm 0.08`$. The rms value corresponding to the random and numerical errors lies in the range $`0.19\begin{array}{c}<\hfill \\ \hfill \end{array}\sigma _\delta \begin{array}{c}<\hfill \\ \hfill \end{array}0.28`$. The reconstructions in Fig. 3 have been carried out with a prior knowledge of the values of $`b,\mathrm{\Omega }_\mathrm{m}`$ for each dataset. Alternatively, the test can be carried out by putting together the procedure followed to obtain the likelihood in Fig. 2 and investigate the scatter resulting in the plots $`\delta _{\mathrm{mock}}`$ vs. $`\delta _{\mathrm{LAP}}`$ for different values of $`b`$,$`\mathrm{\Omega }_\mathrm{m}`$. Supposedly estimating the values of $`b,\mathrm{\Omega }_\mathrm{m}`$ and finding the optimal correlation between $`\delta _{\mathrm{mock}}`$,$`\delta _{\mathrm{LAP}}`$ ought to be two not unrelated operations. However these two appear to be fairly independent: it turns out that whereas (32) gives us the correct likelihood estimates following the criterion of convergence of the solutions at each iteration, the variations in $`\sigma _\delta `$ for a large range of $`b`$,$`\mathrm{\Omega }_\mathrm{m}`$ are fairly small, and $`\sigma _\delta `$ (as computed from tests such as the nine reconstructions in Fig. 3) is too insensitive to be helpful in the estimate of the parameters. Therefore the tests show that the estimate of the parameters and the reconstruction of the fields are two operations that are to a large extent independent. For an arbitrary sample, one would thus first compute (32), pick the values of $`b,\mathrm{\Omega }_\mathrm{m}`$ at the maximum of the likelihood surface and use these to solve the equations to compute $`\delta `$,$`\alpha `$. Similarly, Fig. 4 shows the comparison of the LAP results with the mock data in the reconstruction of the velocity potential. The values of the fields have been scaled to $`\alpha _{\mathrm{max}}`$ and are dimensionless. It is apparent that the regression line is in all cases slightly greater than unity, with a more accentuated tilt for larger values of ($`b`$,$`\mathrm{\Omega }_\mathrm{m}`$). The smaller values of $`\alpha `$ adjust better to a slope of unity, but with larger scatter than larger $`\alpha `$.
Fig. 5 shows a cross-section on the $`Z=0`$ plane of a particular velocity field reconstruction, that of the dataset $`d(b=1.0,\mathrm{\Omega }_\mathrm{m}=0.3)`$. The figure shows several prominent features of the underlying density field in this case: three overdense regions to which the field vectors converge, on the lower left, middle right and upper left parts of the panel, and two prominent underdense regions, from which the velocities diverge, one at the central region and another one at the middle-left boundary of the circle. It is apparent that the LAP velocities are not vanishing in the normal direction of the boundary surface of the selected subvolume, and therefore the customary Neumann spatial boundary conditions employed on spherical Bessel functions (i.e. vanishing normal velocities at the boundary) do not apply. We note that spatial boundary conditions are unnecessary in the LAP reconstruction, thus we have not brought up the issue in §2. The velocity field agrees within 10% accuracy with the $`n`$-body exact field within 78% of the selected volume, and the remaining 22% differs from the mock sample velocities by an error of $`\begin{array}{c}>\hfill \\ \hfill \end{array}10\%`$ (shown in Fig. 4 by the regions enclosed by the solid curves) and withing this volume 6% differs by an error $`\begin{array}{c}>\hfill \\ \hfill \end{array}20\%`$ (regions enclosed by broken curves). These regions are mostly located in the neighbourhood of peaks, right at the very slopes, where the largest velocities are found. The central regions of peaks and troughs are very accurately reconstructed, and it is indeed the intermediate regions that yield $`\delta `$ points with greater scatter in Fig. 3 and worse velocity reconstructions in Fig. 5.
## 4 Bias and $`\mathrm{\Omega }_\mathrm{m}`$ from IRAS 1.2 Jy
We apply the LAP method to the IRAS 1.2 Jy sample (Strauss et al. 1990,1992; Fisher et al. 1995a) in the same way as we have used it in the reconstruction of the mock catalogs in §3. IRAS 1.2 Jy is not the largest existing near all-sky galaxy redshift catalogue, and it is now superseded by PSC$`z`$ (Canavezes et al. 1998) which contains $`15,000`$ galaxies, so this application is simply an illustration on how the LAP method can be used to break the degeneracy in the estimates of $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$. Other large redshift samples of partial coverage can also be looked at with the LAP method, e.g., Las Campanas and the forthcoming Anglo-Australian 2dF ($`250,000`$ galaxies) and US Sloan Digital Sky Survey (SDSS) ($`10^6`$ galaxies and 25% coverage), with the caveat that boundary regions will be a source of propagating errors in the dynamical evolution. Even so, a large number of galaxies in a redshift survey of limited coverage can provide a good representation of the underlying density field, almost definitely outweighing the disadvantages of sampling a partial region of the sky, and it will be thus predictably worthwhile to apply the LAP method to those surveys. The IRAS 1.2 Jy sample contains 5320 galaxies distributed over 87.6% of the projected celestial sphere. The remaining unsampled 2.4% is an approximately disk-shaped region at a galactic latitude $`|b|\begin{array}{c}<\hfill \\ \hfill \end{array}5^{}`$.
We adopt a Gaussian smoothing length of 1200 $`\mathrm{km}\mathrm{s}^1`$, and make no assumption regarding the power spectrum. The data $`d_{\mathrm{IRAS}}`$ are distributed within a spherical region of radius $`x_{\mathrm{max}}15,000`$ $`\mathrm{km}\mathrm{s}^1`$. We use the dataset in a similar fashion as the mock samples $`d(b,\mathrm{\Omega }_\mathrm{m})`$ in the previous section to derive the $`x`$-space fields $`\delta ,\alpha `$. In §3 we have established that $`\sigma _\delta ,\sigma _\alpha `$ are fairly insensitive to the values of $`b,\mathrm{\Omega }_\mathrm{m}`$. One can thus set out to investigate the likelihood function $`\lambda (b,\mathrm{\Omega }_\mathrm{m})`$ as defined in (32) prior to determining the reconstructed fields. Evidently this is the simplest way to proceed for, unlike in §3, we do not have any clue about the real-space underlying fields (such as $`\delta _{\mathrm{mock}}`$, $`\alpha _{\mathrm{mock}}`$ in §3) to compare them with the reconstructed fields.
The likelihood contour plot is shown in Fig. 6. Clearly the largest values of the likelihood function are centered around $`b1`$ and small $`\mathrm{\Omega }_\mathrm{m}`$. From the test of the LAP method in §3 with $`n`$-body simulations we already know that the likelihood function (32) underestimates both $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$, as is apparent in all nine panels of Fig. 2. We accept this trend is fairly inherent to the numerical application of the method and thus infer that the result presented in Fig. 6 is no different in this respect, and therefore the real values of the parameters are situated somewhat above their maxima in the likelihood function. From Fig. 2 one can quantify these errors to be of the order of $`\mathrm{\Delta }\mathrm{\Omega }_\mathrm{m}0.12`$, $`\mathrm{\Delta }b0.15`$. Therefore, we infer that in Fig. 6 the likelihood maxima and the real values of the parameters are likely to be offset by a similar margin of error. At face value, Fig. 6 estimates that the most likely values of the parameters are $`\mathrm{\Omega }_\mathrm{m}0.18`$ and $`b0.94`$. If we offset these estimates by the errors derived from Fig. 2, then the likely “real” values of the parameters that we obtain for Fig. 6 are $`\mathrm{\Omega }_\mathrm{m}0.31`$ and $`b1.1`$. As a matter of fact, these offset values are still within the region enclosed by the 95% confidence contour.
To put our results in perspective with previous analyses of IRAS 1.2 Jy, we have overlaid on the contour plot of Fig. 6 two previous estimates of $`\beta \mathrm{\Omega }_\mathrm{m}^{0.6}/b`$. An estimate by Willick et al. (1997a) yields $`\beta _I=0.49\pm 0.07`$ (shaded region $`A`$) and an estimate by Sigad et al. (1998) yields $`\beta _I=0.89\pm 0.12`$ (shaded region $`B`$). The estimate of Willick et al. (1997a) is clearly in better agreement with our results as the location of the offset maximum of the likelihood is contained within the shaded region $`A`$ that corresponds to the error margin of their estimate. The estimate given by shaded region $`B`$ is consistent with a scenario $`b1`$, $`\mathrm{\Omega }_\mathrm{m}1`$, which in our analysis falls well outside the 10% likelihood contour.
Fig. 7 shows a $`z`$-space comparison between the reconstructed fields and the dataset. The data on the horizontal axis, $`\delta _{\mathrm{LAP}}^s`$, is obtained from the reconstructed $`x`$-space fields $`\delta ,\alpha `$ via (11). The combination of both fields via the relationship $`\delta _{\mathrm{LAP}}(\stackrel{}{x})\delta _s(\stackrel{}{x}+\widehat{x}\alpha _{\mathrm{LAP}}^{})`$ permits us to reconstruct $`\delta _s`$ which is our only possible point of comparison with $`\delta _{IRAS}`$, and this is shown in Fig. 7. The vertical axis shows the $`z`$-space data points of the smoothed IRAS 1.2 Jy sample. The data are plotted in a point-by-point comparison for all the grid points within the selected subvolume. A solid line of slope 1.0 is plotted across the diagonal of the plot. The slope of the regression line is slightly over the diagonal line, at approximately 1.03. The corresponding rms due to random and numerical errors in the LAP reconstruction is $`\sigma 0.27`$. The values of the parameters that have been used in the reconstruction are $`b=1.0,\mathrm{\Omega }_\mathrm{m}=0.3`$.
### 4.1 Velocity fields
The resulting velocity field for the parameters of Fig. 7 is shown in Fig. 8. The six panels show the reconstructed IRAS 1.2 Jy fields $`\delta _{\mathrm{LAP}}`$ and $`\stackrel{}{v}_{\mathrm{LAP}}`$ in supergalactic coordinates, for three different slices $`Z=0,\pm 2000`$ $`\mathrm{km}\mathrm{s}^1`$. The velocity panels on the right column correspond to the densities on the left, at the same value of $`Z`$. The velocity field follows the main features observed on the $`\delta _{\mathrm{LAP}}`$ field, with a general flow towards the overdense regions and outflow from voids. The largest velocities are located in the intervening regions between overdense and underdense regions, e.g. in $`Z=0`$ (middle panels), large infall velocities are visible in the vicinity of the Comma supercluster (0,80,0), the Hydra-Centaurus (H-C) supercluster (-30,15,0), and Perseus-Pisces (P-P) (50,-5,0). In $`Z=0`$ the largest velocities are located at the lower right region of the H-C overdensity maximum, and also to the left of the P-P maximum. There is a velocity flow from the main void on the lower left of the figure, in the direction of Virgo, and it splits up to left and right, in manner of a ridge, to create an outflow in opposite directions, towards H-C and P-P. In the case of $`Z=2000`$ $`\mathrm{km}\mathrm{s}^1`$(lower row), large velocities are also present around the steeper regions of the prominent overdensities, following a similar pattern as in $`Z=0`$, whereas the field shows more erratic features in $`Z=2000`$ $`\mathrm{km}\mathrm{s}^1`$(upper row), where the outflow from the main void (centre left) shows a general trend towards the main overdense features but is at the same time prone to local variations.
The results presented in Figs. 6-8, can be optimized by using the Mark III velocity redshift survey to pin down $`b`$,$`\mathrm{\Omega }_\mathrm{m}`$ more accurately. We shall pursue this and look for the optimal values of $`b`$,$`\mathrm{\Omega }_m`$ by computing the LAP solutions that satisfy
$$\delta (\stackrel{}{v}_{\mathrm{LAP}}\stackrel{}{v}_{\mathrm{MarkIII}})^2=0,$$
(33)
where $`\delta `$ denotes a variation, not the density contrast. In practice, this is achieved as follows. One adds (33) to the two already existing constraints of the LAP method (18),(19). Those are tackled in the manner summarized in §2.5. In actual terms, it’s far more practical to deal with (33) in terms of the velocity potential, so what we have done in the present analysis is in reality to compute $`\alpha _{\mathrm{MarkIII}}`$ from the smoothed observed velocity field, and thus used (33) in the manner of a second constraint on $`\alpha `$.
The comparison with the $`\stackrel{}{v}_{\mathrm{MarkIII}}`$ data sets further constraints on the likelihood contours of Fig. 6 as is shown below. Mark III contains approximately 3,400 galaxies, which are compiled from several sets of elliptical and SO galaxies (Willick et al. 1995,1996,1997a). The sample spans out to $`6000`$ $`\mathrm{km}\mathrm{s}^1`$, though in some directions it is irregularly sampled to $`x_{\mathrm{max}}8000`$ $`\mathrm{km}\mathrm{s}^1`$and $`x_{\mathrm{min}}4000`$ $`\mathrm{km}\mathrm{s}^1`$. The distances are inferred via forward Tully-Fisher and $`D_n\sigma `$ distance indicators which may entail an error in the region 17-21%. Mark III predicts a bulk flow $`v_B194\pm 32`$ $`\mathrm{km}\mathrm{s}^1`$towards the Shapley concentration (Zaroubi, Hoffman & Dekel 1999)(for a low-resolution Gaussian smoothing $`1200`$ $`\mathrm{km}\mathrm{s}^1`$, within a sphere $`r6000`$ $`\mathrm{km}\mathrm{s}^1`$), in contrast to $`v_B250400`$ $`\mathrm{km}\mathrm{s}^1`$ that is estimated in most other samples, including PSC$`z`$ (a compilation of $`v_B`$ estimates is summarized in Dekel 1999b). $`\delta _{IRAS}`$ and $`\delta _{\mathrm{MarkIII}}`$ are consistent with mildly non-linear gravitational instability and linear bias (Sigad et al. 1998), though there are some differences, e.g. the Mark III sample appears to show a strong shear across the Hydra-Centaurus complex that is absent in IRAS 1.2 Jy (as indeed also in ORS). Recent papers have studied in detail the differences between the IRAS 1.2 Jy and Mark III velocity and density fields (Sigad et al. 1998; also Dekel et al. 1999 following an improved version of POTENT).
We consider the Mark III sample with a Gaussian smoothing length of 1200 $`\mathrm{km}\mathrm{s}^1`$. The data are carefully corrected for Malmquist biases (following the recipe set out in Sigad et al. (1998) for the preparation of the data), and the distances of 1,241 objects are modified as a result. The LAP method is solved for IRAS 1.2 Jy within spherical volume of radius $`x_{\mathrm{max}}15,000`$ $`\mathrm{km}\mathrm{s}^1`$, and the minimization fit with Mark III (33) is done within a spherical subvolume of radius $`x6000`$ $`\mathrm{km}\mathrm{s}^1`$. Therefore most of the volume of the LAP solutions remains free of the constraint (33) and the fraction of the volume where $`\stackrel{}{v}_{\mathrm{LAP}}`$ is least-squared to $`\stackrel{}{v}_{\mathrm{MarkIII}}`$ is only 0.064. Naturally such a small fraction forecasts an almost negligible impact in the fine-tuning of the parameters, unless the fields differred drastically to start with, which they do not. The $`\stackrel{}{v}_{\mathrm{LAP}}`$ solution in the remainder of the volume is indirectly affected by this fit, and the variations in modulus $`\mathrm{\Delta }v_{\mathrm{LAP}}`$ outside the comparison subvolume are $`\begin{array}{c}<\hfill \\ \hfill \end{array}12`$%.
Fig. 9 shows the likelihood contours for ($`b`$,$`\mathrm{\Omega }_\mathrm{m}`$) computed via the adjustment entailed in (33). The solid contours are the purely IRAS 1.2 Jy prediction, as in Fig. 6, and the dotted contours are the result of the comparison with Mark III. The contours are ever so slightly shifted towards greater values of the parameters and, as expected, the effect is small. The shift towards larger $`b`$,$`\mathrm{\Omega }_\mathrm{m}`$ is not in fact an altogether undesirable modification, as we have already discussed that the LAP solutions are found to be per se offset to smaller values than their “real” values. The important conclusion to be drawn from Fig. 9 is that the comparison with Mark III is entirely consistent with the predictions for $`b`$ and $`\mathrm{\Omega }_\mathrm{m}`$ extracted from the IRAS 1.2 Jy sample alone.
## 5 Conclusions
The LAP method provides a practical means to break the degeneracy between $`\mathrm{\Omega }_\mathrm{m}`$ and $`b`$ in galaxy redshift surveys. The method is employed in the manner of a nonlinear constraint on the redshift-space dataset and, although in formulation it comes across as algebraically cumbersome, it is of considerable simplicity and efficiency from the numerical point of view. The method is sound in that it does not require an a priori approximation of the map $`\stackrel{}{x}\stackrel{}{s}`$ to pin down the solution and it provides considerable freedom to ascribe relative importance to the data available, i.e. the initial and final endpoints, to which we wish to invariably assign greater weight than intermedate stages of which little or no data are available.
The method can prove significant to measure $`\mathrm{\Omega }_\mathrm{m}`$ in the latest largest samples, and extract the most accurate information prior to comparison with other datasets, such as the CMB radiation power spectrum and SN data. One important challenge for the future is to attain a better grasp of the concept of bias and this will be probably achieved via microlensing data and $`n`$-body simulations of the formation of galaxies and clusters from primordial fluctuations, rather than from galaxy redshift surveys. Once a model of bias is adopted on a sound footing, then clearly the LAP model is impeccable in producing an estimate of $`\mathrm{\Omega }_\mathrm{m}`$. In the simple linear bias model we have employed we have totally relegated any consideration of scale-dependence in $`b`$. This is a point I have deliberately omitted for simplicity. Thus, the estimates computed in this paper ought to be regarded qualitatively as weighted averages of the “real” $`b`$ over different scales, if indeed scale-dependent bias models are to be believed.
In this paper, we have employed the likelihood function (32) to investigate the values of $`b`$,$`\mathrm{\Omega }_\mathrm{m}`$. Clearly this is not a unique choice. However, our choice is guided by the argument of relative convergence of the solutions, which is justifiably a reasonable criterion to get close to the “real” solutions. In view of the performance of the $`\lambda `$ function in the reconstruction of the mock samples, this choice does not appear to be totally off the mark. A potential reason for concern could be the offset observed between the maxima of the likelihood functions and the real values of the parameters in the $`n`$-body simulations. However the recurrence of this offset in a predictable manner lends strength to the argument that it arises as a numerical fault that is easy to account for systematically in the analysis of the datasets. The reconstructions of the fields are, on the other hand, of considerable accuracy and no numerical defficiency or hindrance is observed. The application of the method to IRAS 1.2 Jy predicts the parameters to be fairly accurately located in the immediate neighbourhood of the maxima $`\mathrm{\Omega }_\mathrm{m}0.3`$ and $`b1.1`$, which is found to be most compatible with the estimate of $`\beta `$ given by Willick et al. (1997a). In a flat universe such predicted values are perfectly consistent with a non-vanishing cosmological constant or a quintessence scalar field component. The likelihood examined in this way is only very slightly modified when the velocities predicted via the LAP method are finely-tuned with data from the Mark III sample. The shift of the predicted values is towards slightly greater values of the parameters but it remains comfortably consistent with the results obtained from IRAS 1.2 Jy.
## Acknowledgments
The numerical analysis has been performed in the Starlink facilities at Queen Mary & Westfield College (London), and thanks are due to a number of people for supplying and lending a helping hand with technicalities regarding the data. The author thanks an anonymous referee for insightful comments that have greatly improved the readability of the manuscript. This research has been funded in part at EHU by research grant UPV172.310-G02/99.
## Appendix A Orbit-crossing in redshift space
We shall prove the boundary condition (11). The number-counts of galaxies $`n`$ in $`x`$-space and $`z`$-space satisfy, by conservation of the number of galaxies:
$$\mathrm{d}n(\stackrel{}{s})=\underset{\mathrm{streams}}{}\mathrm{d}n(\stackrel{}{x}_i),$$
(34)
for all streams at the same redshift, $`\stackrel{}{s}=\stackrel{}{x}_i+\widehat{x}(\widehat{x}_x\alpha _i)`$. In our analysis we shall only consider single-valued solutions, and therefore there is just one stream only in (34), i.e. $`\mathrm{d}n(\stackrel{}{s})=\mathrm{d}n(\stackrel{}{x})`$. Hence
$$\rho _s(\stackrel{}{s})\mathrm{d}\Omega \frac{\mathrm{d}n(\stackrel{}{s})}{\mathrm{d}s}=x^2(1+b\delta )\frac{N_{\mathrm{gals}}}{V}\frac{\mathrm{d}x}{\mathrm{d}s}\mathrm{d}\Omega ,$$
(35)
where $`n(\stackrel{}{s})`$ is the galaxy number-count, $`\mathrm{d}\Omega `$ a solid angle element and the $`x`$-space selected volume of the sample is $`V\frac{4}{3}\pi x_{\mathrm{max}}^3`$, and
$$s\widehat{x}\stackrel{}{s}=x+\alpha ^{}.$$
(36)
Therefore
$$\frac{\mathrm{d}x}{\mathrm{d}s}=\frac{1}{1+\alpha ^{\prime \prime }},$$
(37)
and substituting this in (35), we get
$$\rho _s=x^2\frac{N_{\mathrm{gals}}}{V}\left(\frac{1+b\delta }{1+\alpha ^{\prime \prime }}\right).$$
(38)
In the case of multistreams, the RHS of (38) is integrated over all streams, bearing in mind that turn-around regions, which occur at $`\delta 1`$ and for which $`\mathrm{d}s/\mathrm{d}x=0`$, are excluded from the sum. An example of such a region in shown in Fig. A1. An initial saddle-point $`\mathrm{d}s/\mathrm{d}x=0`$ on the $`s(x)`$ curve starts the creation of a turn-around region. At the stage shown in Fig. A1, both points $`A`$ and $`B`$ satisfy this condition and obviously they departed from an initial saddle-point $`A=B`$. The region spanning between $`A`$ and $`B`$ is three-valued (each redshift in the interval $`z_B<z<z_A`$ corresponds to three $`x`$ positions), whereas $`z_A`$ and $`z_B`$ are bivalued. To make such an scenario tractable, we need to replace $`s(x)`$ over the interval $`z_B<z<z_A`$ by a monotonic curve that matches the existing curve at $`z_B`$ and $`z_A`$ and its first derivative. This is obviously tantamount to applying a larger smoothing length than the existing one to erase the overdense region that is the cause of the turn-around.
## Appendix B Evaluation of radial derivatives
The radial derivative of the velocity potential coefficients (15) can be written as
$$\frac{\mathrm{d}}{\mathrm{d}x}\alpha _n=\underset{rlm}{}\alpha _{}^{}{}_{rlm}{}^{(n)}j_l(k_rx)Y_{lm},$$
(39)
where, using the equality
$$\frac{\mathrm{d}}{\mathrm{d}u}j_l(u)=(2l+1)^1\left[lj_{l1}(u)(l+1)j_{l+1}(u)\right],$$
(40)
we have
$$\alpha _{}^{}{}_{rlm}{}^{(n)}=k_r\left[\frac{(l+1)}{(2l+3)}\alpha _{r(l+1)m}^{(n)}\frac{l}{(2l1)}\alpha _{r(l1)m}^{(n)}\right].$$
(41)
Similarly
$$\alpha _{}^{\prime \prime }{}_{rlm}{}^{(n)}=k_r^2\{\frac{(l+1)}{(2l+3)}\frac{(l+2)}{(2l+5)}\alpha _{r(l+2)m}^{(n)}$$
$$\left[\frac{(l+1)^2}{(2l+3)(2l+1)}+\frac{l^2}{(2l1)(2l+1)}\right]\alpha _{rlm}^{(n)}$$
(42)
$$+\frac{l}{(2l1)}\frac{(l1)}{(2l3)}\alpha _{r(l2)m}^{(n)}\}.$$
On the other hand, the coefficients $`\stackrel{}{J}_{lm}^{(n)}`$ given in (16) are
$$\stackrel{}{J}_{lm}^{(n)}=\underset{r}{}[\frac{\alpha (l,m+1)}{2}\alpha _{rl(m+1)}^{(n)}(i\widehat{x}_1\widehat{x}_2)$$
$$+\frac{\beta (l,m1)}{2}\alpha _{rl(m1)}^{(n)}(i\widehat{x}_1+\widehat{x}_2)+im\alpha _{rlm}^{(n)}\widehat{x}_3],$$
(43)
where
$$\alpha (l,m)=\left[l(l+1)m(m1)\right]^{1/2},$$
(44)
$$\beta (l,m)=\left[l(l+1)m(m+1)\right]^{1/2}.$$
(45)
## Appendix C Chebyshev polynomials
The Chebyshev polynomials are defined $`T_n(\mathrm{cos}\theta )\mathrm{cos}(n\theta )`$ (following the normalization of Abramowitz & Stegun 1972). We define the angle brackets $`,`$ according to the orthogonality properties of $`T_n`$ (e.g. Courant & Hilbert 1989):
$$u_1^1dtw(t)u(t),$$
(46)
where $`w(t)=(1t^2)^{1/2}`$ is a weight function and therefore
$$T_nT_m=\delta _{nm}\frac{\pi }{2},$$
(47)
for $`n0`$ and $`T_0^2=\pi `$. In (22)(23) we encounter two types of angle brackets to evaluate (other than (47): $`T_n\dot{T}_m`$ and $`T_nT_mT_r`$ (we have deliberately omitted $`\mathrm{\Omega }_\mathrm{m}T_nT_m`$, by approximating $`\mathrm{\Omega }_\mathrm{m}`$ by a constant, and ditto for $`H`$. The second type of product is trivially transformed into (47) via
$$2T_nT_m=T_{n+m}+T_{nm}$$
(48)
for $`nm`$, and the first requires a little numerical manipulation using the relation
$$(1t^2)\dot{T}_n=ntT_n+T_{n1}.$$
(49)
## Appendix D Orthogonality relations
The orthogonality relations for the spherical harmonics and the Bessel functions are respectively
$$_0^{2\pi }d\phi _0^\pi \mathrm{d}(\mathrm{cos}\theta )Y_{lm}Y_{l^{}m^{}}=\delta _{ll^{}}\delta _{mm^{}},$$
(50)
$$_0^1dxx^2j_l(k_rx)j_l(k_sx)=\frac{1}{2k_rk_s}\left[j_l(k_rx)+xj_l^{}(k_rx)\right]^2\delta _{rs}.$$
(51) |
warning/0002/cond-mat0002380.html | ar5iv | text | # Snake orbits and related magnetic edge states
## I Introduction
The transport properties of a two-dimensional electron gas (2DEG) subjected to a nonhomogeneous perpendicular magnetic field (periodically modulated or not) have been the focus of a great deal of research in recent years. Current fabrication technologies permit to create nonhomogeneous magnetic fields on a nanometer scale by deliberately shaping or curving the 2DEG, or by integration of superconducting or ferromagnetic materials on top of the 2DEG. This will add a new functional dimension to the present semiconductor technology and will open avenues for new physics and possible applications.
Theoretically, the effect of nonhomogeneous magnetic fields on a 2DEG have been studied both in the ballistic and the diffusive regime. The resulting perpendicular magnetic field can act as a scattering centre, but can also bind electrons, and so influence the transport properties of the 2DEG. In transport calculations one needs the electron states, which are obtained by solving the Schrödinger equation.
Müller studied theoretically the single particle electron states of a 2DEG in a wide quantum waveguide under the application of a nonuniform magnetic field and showed that in the case of a magnetic field modulation in one direction, transport properties also become one dimensional and electron states propagate perpendicularly to the field gradient.
Making use of this decoupling, the electron states for different nonhomogeneous magnetic field profiles along one dimension were investigated, i.e. for a periodically modulated magnetic field , for magnetic quantum steps, barriers and wells in an infinite 2DEG and in a narrow waveguide.
In this paper we consider an infinite 2DEG subjected to a step-like magnetic field, i.e abruptly changing in magnitude or polarity at $`x=0`$, in one dimension (taken to be the $`x`$-direction). Preliminary results were presented in Ref. . First the situation for two opposite homogeneous magnetic fields with the same strength will be considered. The classical trajectories correspond to snake orbits and were already used in the seventies to describe electron propagation parallel to the boundary between two magnetic domains. Back then, one was interested in understanding the electron transport through multi-domain ferromagnets and it turned out to be more convenient to work with the classical trajectrories than with the corresponding electron states, which allows one to use a semi-classical theory which reduces the complexity of the theory considerably. We are interested in transport through a 2DEG situated in a semiconductor in which the Fermi energy is orders of magnitude smaller than in the metallic systems of Ref. .
We will study thoroughly the quantum mechanics of such electron states in a 2DEG subjected to this step magnetic field profile, and we will compare them with their classical counterpart. We will discuss the energy spectrum and the corresponding electron states, and derive several properties. We will show the existence of states wich have a velocity opposite to the expected classical orbits. Additionally, we will show that adding a background magnetic field modifies the spectrum and the states considerably.
The paper is organized as follows. In Sec. II we present our theoretical approach. In Sec. III we calculate the energy spectrum, the wavefunctions and their corresponding group velocity, and compare this with their quantum mechanical counterpart. In Sec. IV we study the influence of a background magnetic field on the quantum mechanical and classical behaviour. In Sec. V we focus on the negative velocity state, and finally, in Sec. VI, we construct time dependent states, and interpret them classically for several magnetic field profiles.
## II Theoretical approach
We consider a system of noninteracting electrons moving in the $`xy`$-plane in the absence of any electric potentials. The electrons are subjected to a magnetic field profile $`\stackrel{}{B}=(0,0,B_z\left(x\right))`$. First, we will study the electronic states near the edge of two magnetic fields with opposite strength
$$B_z\left(x\right)=B_0\left[2\theta (x)1\right],$$
(1)
which is independent of the $`y`$-coordinate. Next, we will consider the influence of a background magnetic field $`B`$ on these states, which results in the magnetic field profile
$$B_z\left(x\right)=B_0\left[2\theta (x)1\right]+B.$$
(2)
In the following we will use $`B^l=B_z\left(x<0\right)`$ and $`B^r=B_z\left(x>0\right)`$, to denote respectively the magnetic field on the left and the right hand side of the magnetic edge.
The one-particle states in such a 2DEG are described by the Hamiltonian
$$H=\frac{1}{2m_e}p_x^2+\frac{1}{2m_e}\left[p_y\frac{e}{c}A\left(x\right)\right]^2.$$
(3)
Taking the vector potential in the Landau gauge,
$$\stackrel{}{A}=(0,xB_z\left(x\right),0),$$
(4)
we arrive at the following 2D Schrödinger equation
$`\left\{{\displaystyle \frac{^2}{x^2}}+\left[{\displaystyle \frac{}{y}}+ixB_z(x)\right]^2+2E\right\}\psi (x,y)=0,`$
where the magnetic field is expressed in $`B_0`$, all lengths are measured in the magnetic length $`l_B=\sqrt{\mathrm{}c/eB_0}`$, energy is measured in units of $`\mathrm{}\omega _c`$, with $`\omega _c=eB_0/m_ec`$ the cyclotron frequency and the velocity is expressed in units of $`l_B\omega _c`$.$`H`$ and $`p_y`$ commute due to the special form of the gauge, and consequently the wavefuntion becomes
$$\psi (x,y)=\frac{1}{\sqrt{2\pi }}e^{iky}\varphi _{n,k}(x),$$
(5)
which reduces the problem to the solution of the 1D Schrödinger equation
$$\left[\frac{1}{2}\frac{d^2}{dx^2}+V_k(x)\right]\varphi _{n,k}(x)=E_{n,k}\varphi _{n,k}(x),$$
(6)
where it is the $`k`$-dependent effective potential
$$V_k\left(x\right)=\frac{1}{2}\left[xB_z\left(x\right)k\right]^2,$$
(7)
which contains the two dimensionality of the problem. We will solve Eq. (6) numerically by use of a discretization procedure. In some limiting cases analytical results can be obtained.
## III In the absence of a background magnetic field
Let us first consider the case when no background magnetic field is present. The situation is then symmetric, and more easily to solve. The effective potential for this case is shown in Fig. 1(a) for $`k=2`$ (dotted curve) and $`k=2`$ (solid curve). We notice from Eq. (7) that this potential is built from two parabolas, with minima situated at $`x^l=k`$, and $`x^r=k`$, thus respectively on the left and right hand side of the magnetic edge. The total potential has for $`k>0`$ two local minima respectively at $`x=k`$ and $`x=+k`$, while for $`k<0`$ it has only one minimum at $`x=0`$. Before we describe the energy spectrum of the snake orbits and their corresponding properties, we first discuss the limiting behaviour.
### A Limiting behaviour for $`k\pm \mathrm{}`$
For $`k\mathrm{}`$, the minima of the parabolas are situated far from each other. The electrons are in the Landau states of two opposite magnetic fields, one on the left, the other on the right, and they are not interacting with each other. The electron wavefunctions are given by $`L|x=C_mH_m(x+k)e^{\left(x+k\right)^2/2}`$, and $`R|x=C_mH_m(xk)e^{\left(xk\right)^2/2}`$, respectively, where $`H_m\left(x\right)`$ is the Hermite polynomial. For decreasing $`k`$ the parabolas shift towards each other, and the electrons will start to “feel” each other. In terms of wavefunctions, this results in a parabolic cylinder function $`\varphi (x)=D_{E\frac{1}{2}}\left[\sqrt{2}\left(xk\right)\right]`$ matched at $`x=0`$, with the condition that $`\frac{d}{d\alpha }D_{E\frac{1}{2}}\left(\alpha \right)|_{\alpha =\sqrt{2}k}=0`$ or $`D_{E\frac{1}{2}}(\sqrt{2}k)=0`$, for the symmetric and the antisymmetric wavefunction, respectively. This leads to a change in energy of the electron states, which can be understood as a lifting of the degeneracy of the two original electron wavefunctions. The energy can then be written as $`E_\pm (k)=L\left|H\right|L\left(k\right)\pm L\left|H\right|R\left(k\right)=E\left(k\right)\pm \mathrm{\Delta }E\left(k\right)`$ with the corresponding wavefunctions $`|\varphi =|R|L`$. One can see that the presence of the second parabola results in two effects: (1) a decrease of$`E\left(k\right)`$ due to the finite presence of the wave function in the other parabola, and (2) a splitting of the energy level due to the overlap, i.e. one level ($`E_+`$) shifts upwards, while the other ($`E_{}`$) shifts down. For $`k\mathrm{}`$, this results in the following first-order approximation to $`E`$ and $`\mathrm{\Delta }E`$:
$`E_m\left(k\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}+m{\displaystyle \frac{2^{m1}}{m!\sqrt{\pi }}}k^{2m1}e^{k^2},`$ (9)
$`\mathrm{\Delta }E_m\left(k\right)`$ $`=`$ $`{\displaystyle \frac{2^m}{m!\sqrt{\pi }}}e^{k^2}k^{2m+1}.`$ (10)
In the other limit $`k\mathrm{}`$, the effective potential can be approximated by a triangular well $`V\left(x\right)=k^2/2kB_0x`$. Solutions for this potential consist of Airy functions, again matched at $`x=0`$ with the condition that $`\varphi ^{}\left(0\right)=0`$ or $`\varphi \left(0\right)=0`$ which results, respectively into the anti-symmetric wavefunction $`\varphi _{2m}\left(x\right)=C_{2m}\left(k\right)\left(\left|x\right|/x\right)Ai\left[z_{Ai^{},m+1}+\left(2k\right)^{1/3}\left|x\right|\right]`$ and a symmetric one $`\varphi _{2m+1}\left(x\right)=C_{2m+1}\left(k\right)Ai\left[z_{Ai,m+1}+\left(2k\right)^{1/3}\left|x\right|\right]`$, respectively with energy
$`E_{2m}\left(k\mathrm{}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[k^2z_{Ai^{},m+1}\left(2\left|k\right|\right)^{2/3}\right],`$ (12)
$`E_{2m+1}\left(k\mathrm{}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[k^2z_{Ai,m+1}\left(2\left|k\right|\right)^{2/3}\right],`$ (13)
where $`z_{Ai,n}`$ ($`=2.338`$, $`4.088`$, $`5.521`$, $`\mathrm{}`$, $`\left[3\pi \left(4n1\right)/8\right]^{2/3}`$) and $`z_{Ai^{},n}`$ ($`=1.019,`$ $`3.248,`$ $`4,820,`$ $`\mathrm{},`$ $`\left[3\pi \left(4n3\right)/8\right]^{2/3}`$), denote respectively the $`n^{th}`$ $`\left(n=1,2,3,\mathrm{},\mathrm{}\right)`$ zero of the Airy function and of its derivative. One can see that for increasing negative $`k`$, the difference between the two energy branches increases, which is to first order linear in $`\left|k\right|`$. Namely the more negative $`k`$, the narrower the well, thus the more the energy levels are shifted from each other.
### B Spectrum and velocity
Solving Eq. (6) numerically gives rise to the energy spectrum shown (solid curves) in Fig. 2(a). For $`k=\mathrm{}`$, we obtain the earlier mentioned Landau levels, which are labelled with the quantum number $`m`$. Each level is twofold degenerate. For decreasing $`k`$, the degeneracy is lifted and they separate into two different branches with eigenstates $`|2m`$ and $`|2m+1`$ and eigenvalues $`E_{2m}`$ and $`E_{2m+1}`$, and corresponding quantum numbers $`n=2m`$ and $`n=2m+1`$. This quantum number $`n`$ does not only result from arranging the levels according to their lowest energy, starting with $`n=0`$, but it also reflects the number of nodes of the corresponding wavefunction. Notice that the levels have now a non-zero derivative, i.e. electrons propage in the $`y`$-direction, and their group velocity is given by $`v\left(k\right)=E\left(k\right)/k`$. (The minus sign appears here because in Eq. (5) we took $`k_y=k`$). This group velocity is plotted (solid curves) in Fig. 3(a) for the 6 lowest levels. For $`k=\mathrm{}`$ electrons are in a Landau level, and consequently there is no net current in the $`y`$-direction. Decreasing $`k`$, results in a net current in the $`y`$-direction, which is positive for the upper branches $`\left(2m+1\right)`$, but is initially negative for the branches $`\left(2m\right)`$. For more negative values of $`k`$ it increases almost linearly with increasing $`\left|k\right|`$, which becomes the first order analytical result $`v_m(k\mathrm{})=k`$, obtained by differentiating (12) and (13).
### C Classical picture
The center of the classical orbit corresponds to a zero in the effective potential. The energy spectrum can be divided up into three regions which can be classicaly understood by the electron orbits drawn in Fig. 1(a). In region (A) the electrons move in closed orbits either in the magnetic field on the left or on the right hand side. Since its cyclotron radius is smaller than the distance to the magnetic field discontinuity, they feel a homogeneous magnetic field. There is no net velocity. In region (B) the cyclotron radius intersects the magnetic field discontinuity slightly, i.e. in such a way that the moving electron and the center of its orbit are on the same side. The electron is nevertheless able to penetrate in the opposite magnetic field region, which results in a (rather small) propagation in the $`y`$ direction $`v_y>0`$. For $`k=0`$ the center of the orbit is exactly on the edge between the two opposite magnetic fields. In region (C) the center is located in the opposite magnetic field region of which the electron is moving in, resulting in a faster propagation of the electron in the $`y`$-direction. These different regions are also indicated in Fig. 2(a).
We can also make a quantitative classical study of the velocity, starting from the quantum mechanical energy spectrum. Since in classical mechanics there is no quantization, we make use of the obtained quantum energy spectrum in order to find the classical energy and thus the radius of the cyclotron orbit. Classically, the energy is contained in the circular velocity $`v_\phi `$ through $`E\left(k\right)=v_\phi ^2\left(k\right)/2`$. For any given quantum mechanical $`E(k)`$-value we obtain classically the circular velocity $`v_\phi (k)=\sqrt{2E(k)}`$. Now if we consider $`x_0=\pm k`$ to be the center of the electron orbit, we can calculate for every $`k`$-value the classical velocity $`v_y(k)`$, since we also know $`v_\phi (k)`$ and the cyclotron radius $`R(k)=v_\phi (k)`$. Using geometric considerations, we obtain the following relation
$$v_y(k)=v_\phi (k)\sqrt{1[k/v_\phi (k)]^2}/\mathrm{arccos}[k/v_\phi (k)],$$
(14)
which is shown in Fig. 3(a) by the dotted curves. Comparing this with its quantum mechanical counterpart, we notice that for $`k<0`$ values good agreement is found, but for $`k>0`$ there is a large discrepancy. Moreover, one can see that negative velocities cannot exist classically.
The critical $`k`$-value, $`k^{}`$, for which no classical propagating states can exist, i.e. the electron describes just a circular orbit in a homogeneous magnetic field, has to be equal to the cyclotron radius $`k^{}=R\left(k^{}\right)=v_\phi \left(k^{}\right)=\pm \sqrt{2E\left(k^{}\right)}`$, which leads to the boundary drawn in Fig. 2(a) (dotted parabola).
## IV With a background magnetic field
With a background magnetic field three different configurations: a) $`0<B<B_0`$, b) $`B=B_0`$, and c) $`B_0<B`$ must be considered. In the following we will study the snake orbits in these configurations.
### A $`0<B<B_0`$
Applying a background magnetic field $`0<B<B_0`$, results in a situation which is very similar to the previous one. Again the two magnetic fields have opposite sign, but in this case they also have a different strength, i.e. $`B^l=B^r/p`$. Again we can calculate analytically the correction to the energy in the limit $`k\mathrm{}`$. For an electron on the right hand side in the $`m^{th\text{ }}`$Landau state of a magnetic field with strength $`B^r=B_0`$, the deviation from the Landau energy due to the presence of the other parabola in the effective potential is given by the following matrix element, which to second-order reads
$`E_m(k\mathrm{})`$ $`=`$ $`\left|R\left|H\right|R\left(k\right)\right|`$ (15)
$`=`$ $`\left[m+{\displaystyle \frac{1}{2}}\right]{\displaystyle \frac{2^{m2}}{m!\sqrt{\pi }}}\left(1+{\displaystyle \frac{1}{p}}\right)k^{2m1}e^{k^2}.`$ (16)
For an electron on the left hand side, i.e. in the smaller magnetic field $`B^l=B_0/p`$ region, in the $`m^{th}`$ Landau level, this results in
$`E_m(k\mathrm{})`$ $`=`$ $`\left|L\left|H\right|L\left(k\right)\right|`$ (17)
$`=`$ $`{\displaystyle \frac{1}{p}}\left[m+{\displaystyle \frac{1}{2}}\right]{\displaystyle \frac{2^{m2}}{m!\sqrt{\pi }}}\left(1+{\displaystyle \frac{1}{p}}\right)k^{2m1}e^{k^2}.`$ (18)
Also in this case the energy is smaller then the corresponding Landau energy. The downward energy shift decreases for increasing $`p`$.
If $`p`$ is an integer, Landau states on the left and right hand side, respectively with quantum number $`pm`$ and $`m`$, coincide for $`k\mathrm{}`$. As a consequence these states have an overlap, which reads to first order
$`L\left|H\right|R`$ $`=`$ $`\left(1\right)^{m+1}2^{m\left(p+1\right)/2}\left({\displaystyle \frac{1}{\left(pm\right)!m!\pi }}\right)^{1/2}`$ (20)
$`\times p^{pm/2}e^{k^2\left(1+p\right)/2}k^{m\left(p+1\right)+1}.`$
One can see that for decreasing magnetic field, i.e. increasing $`p`$, this function decreases because of the exponential factor. The electron wavefunction in the lower magnetic field region is extended over a larger region, and further away from the other (center at $`kp`$). The overlap therefore decreases with increasing $`p`$. As a result of this, the energy for $`k\mathrm{}`$ and $`p>1`$ is given by $`R\left|H\right|R`$ and $`L\left|H\right|L`$.
For $`p=1`$, we obtain the previous result, but for increasing $`p`$, the second order term in Eqs. (16) and (18) becomes more important than Eq. (20), because of the exponential factor. The splitting is lifted, and the main contribution to the negative velocity for $`k\mathrm{}`$ arises from Eq. (20) due to the finite extend of the wavefunction in the other parabola.
As an example we studied numerically the case when a background magnetic field $`B=B_0/2`$ is applied, i.e. $`B^l=B_0/2`$ and $`B^r=3B_0/2`$. As one can see in Fig. 1(b), this results in two parabolas with different minima and confinement strength. The resulting spectrum (see Fig. 2(b)) is very similar to the one of the previous case, but unlike the previous symmetrical case, not all states are twofold degenerate for $`k\mathrm{}`$. We now obtain two different sets of Landau states, corresponding to electrons moving in different magnetic field regions with different strength. In this case the second Landau level on the left coincides with the first on the right. The classical picture for the three different regions corresponds to the one drawn in Fig. 1(b), and is also similar to the previous case, except for the different cyclotron radii. With this picture in mind, one can again calculate the classical velocity, which turns out to be identical to Eq. (14). From Fig. 3(b) we notice that again we obtain good agreement for $`k>0`$, but for $`k<0`$, there is a large discrepancy. The negative velocity can also in this case not be explained classically.
The critical $`k`$-value $`k^{}=\sqrt{2E\left(k^{}\right)}`$ for which snake orbits are classicaly possible are indicated by the parabola in Fig. 2(b).
### B $`B=B_0`$
When a background magnetic field $`B=B_0`$ is applied, we obtain the magnetic barrier studied in Ref., where the magnetic field is different from zero only in the region $`x>0`$, i.e. $`B^l=0`$ and $`B^r=2B_0`$. From Fig. 1(c) one can see that in this case the potential is made up of only one parabola and on the left side it is a constant $`k^2/2`$. The energy spectrum and corresponding velocities for this particular case are shown in Fig. 2(c) and 3(c), respectively. We notice that for $`k\mathrm{}`$ we again obtain Landau states, which correspond to bound states on the right hand side of the magnetic edge. Consistent, as being a limiting case of the former magnetic field states, i.e. $`p=\mathrm{}`$, the energy decreases with decreasing $`k`$ and there is no splitting of the energy levels. Thus now we only have states which propagate with negative velocity to which we cannot assign a classical interpretation.
Also in this case we can divide up the spectrum into three regions: (A) the electrons move in closed orbits in the magnetic field region on the right hand side, (B) electrons are free, propagate forward and are reflected on the barrier and (C) electrons are free, propagate backward and are reflected on the magnetic edge. Notice that for a free electron, the energy is larger than $`k^2/2`$, since now the electron also propagates in the $`x`$-direction and consequently has an additional kinetic energy $`k_x^2/2`$.
Classically, propagating states in the magnetic field region do not exist, only Landau states do. The boundary where these classical trajectories are possible is again given by $`k^{}=\sqrt{2E\left(k^{}\right)}`$.
### C $`B_0>B`$
By applying a background magnetic field with strength larger than $`B>B_0`$, we arrive at the situation where $`0<B^l<B^r`$. The magnetic fields on the left and the right hand side have the same sign, but a different strength, i.e. $`B^l=B^r/p.`$
To obtain the energy in the limits $`k\pm \mathrm{}`$, we again can approximate the wavefunction as being in a Landau state in the corresponding magnetic field. We found
$`E(k\mathrm{})`$ $`=`$ $`R\left|H\right|R\left(k\right)`$ (21)
$`=`$ $`\left[m+{\displaystyle \frac{1}{2}}\right]`$ (23)
$`{\displaystyle \frac{2^{m2}}{m!\sqrt{\pi }}}\left(1+{\displaystyle \frac{1}{p}}\right)k^{2m1}e^{k^2},`$
for an electron on the right hand side in the $`m^{th\text{ }}`$Landau state of a magnetic field with strength $`B^r=B_0`$. For an electron on the left hand side, in the smaller magnetic field $`B^l=B_0/p`$ in the $`m^{th}`$ Landau level, we have
$`E(k\mathrm{})`$ $`=`$ $`L\left|H\right|L\left(k\right)`$ (24)
$`=`$ $`{\displaystyle \frac{1}{p}}\left[m+{\displaystyle \frac{1}{2}}\right]+{\displaystyle \frac{2^{m2}}{m!\sqrt{\pi }}}\left(1+{\displaystyle \frac{1}{p}}\right)k^{2m1}e^{k^2},`$ (25)
which results in a negative velocity.
The energy spectrum and the velocity of these eigenstates for the case when $`B=3B_0/2`$, i.e. $`B^l=B_0/2`$, $`B^r=5B_0/2`$, are plotted respectively in Fig. 2(d) and 3(d). The center of the orbit is situated on the right side for $`k>0`$, for $`k<0`$ it is on the left side. For $`k\pm \mathrm{}`$, the electrons move in a homogeneous magnetic field (on the left ($`k\mathrm{}`$) or right ($`k+\mathrm{}`$) hand side of $`x=0`$), and thus $`v_y=0`$.
From Fig. 1(d) one notices that there is only one minimum in the effective potential because the minima of both parabolas are now situated on the same side. The trajectories corresponding with regions (A), (B), and (C) are depicted in Fig. 1(d). The trajectories in region (A’) are similar to those in (A) but now for a magnetic field on the left hand side, i.e. with smaller strength.
Geometrical considerations yield the following classical velocity
$`v_y(k)`$ $`=`$ $`2v_\phi (k)\sqrt{1\left[k/v_\phi (k)\right]^2}`$ (28)
$`\times \{B^l\mathrm{arccos}[k/v_\phi (k)]`$
$`+B^r\mathrm{arccos}[k/v_\phi (k)]\}^1,`$
which is plotted in Fig. 3(d) as dotted curves together with the quantum mechanical group velocity. One can see that, in contrast to the previous cases, the negative velocity can be understood as classical snake orbits, but these snake orbits all run in the same $`y`$-direction and now there are no states with $`v_y>0.`$
Notice that the quantum mechanical velocity exhibits a small oscillatory behaviour on top of a uniform profile. These whiggles can be understood from the electron distribution over the two parabolas (see Fig. 4). With increasing $`k`$, the electron distribution is shifted from the left parabola to the right one. Due to the wavelike character of this distribution, the probability for an electron to be in the right parabola (integrated solid region in the inset of Fig. 4) exhibits whiggles as function of $`k`$, with $`n`$ maxima as shown in Fig. 4. Energetically it is favourable for an electron state to have as much as possible electron probability in the lower potential region. Consequently, when the electron probability in the lower potential region attains a maximum, a maximum downward energy shift will be introduced on top of the overall energy change, and this will result in a maximum in the group velocity.
## V Negative velocity state
Formally, the existence of the quantum mechanical negative velocity state can be attributed to the fact that shifting two one dimensional potential wells towards each other results in a significant rearrangement of the energy levels in the composite potential well. Because the composed well is broader, some states, p.e. the ground state, have an energy which is lower than in each of the individual narrower wells. In this particular case, the wave vector $`k`$ measures the distance of the two wells to each other, and consequently this energy decrease results in a negative group velocity $`E/k`$. In this section we focus on these negative group velocity states.
Since the negative velocity states are present for any background magnetic field $`B`$, but can only be understood classically in the situation $`B>B_0`$, we will investigate the group velocity for a fixed $`k`$ value with varying background magnetic field. In Fig. 5 the spectrum is plotted as function of the applied background field $`B`$. We have chosen $`k=1.5`$ because in this case a large negative velocity is obtained for the lowest level when $`B=0`$. Notice that for $`B<B_0`$: (1) almost all levels decrease in energy with increasing background field; (2) there is an anti-crossing for $`E/E_0=0.4+0.458B/B_0`$ (dotted line). This anti-crossing occurs when $`B/B_0=n/(n+1)`$, with $`n`$ the Landau level index. For this condition some of the Landau levels are degenerate in the limit $`k\mathrm{}`$ (see Fig. 2(b) for the case $`n=0`$); and (3) for $`BB_0`$ the separation between the levels decreases to zero and a continuous spectrum is obtained with a separate discrete level at the anti-crossing line. The continuous spectrum for $`B=B_0`$ results from the scattered states in the potential of Fig. 2(c), while the discrete state is the bound state in this potential. For the considered $`k`$-value, i.e. $`k=1.5`$, only one bound state is found for $`B=B_0`$.
The corresponding group velocity $`v_y=E/k`$ is shown in Fig. 6. Notice that the maximum negative velocity is obtained near the anti-crossings in the energy spectrum (Fig. 5). Near $`B/B_0=n/\left(n+1\right)`$ the splitting in the energy spectrum (see Fig. 2) is largest and as a consequence one of the levels is pushed strongly down in energy and consequently $`v_y`$ becomes strongly negative. Notice that: (1) every level has some $`B/B_0`$ region at which $`v_y<0`$, and (2) for $`B/B_01`$ the velocity $`v_y0`$, while (3) the enveloppe of $`\left(v_y\right)_{\mathrm{min}}`$ in Fig. 6 reaches for $`B/B_0`$ the $`v_y<0`$ value of the $`B=B_0`$ state. For $`B>B_0`$ we have $`v_y<0`$ for all states.
Using expressions (14) and (28), we can also calculate the classical velocity corresponding to the energy spectrum in Fig. 5. This is shown in Fig. 7. We notice that for $`B/B_0<1`$, the classical velocity has a similar behaviour as the quantum mechanical one, except for the anti-crossings and the lack of negative velocities, which do not have a classical analogon. These negative velocities appear suddenly for $`B/B_0>1`$ and exhibit more or less the same behaviour.
As was already apparent from the above study a necessary condition for the existence of the non-classical edge states is the presence of two local minima in the effective potential. In the limiting case $`B=B_0`$ the second minima is the limiting case of a flat region in $`V_k\left(x\right)`$ for $`x<0`$. But not all these states have a negative velocity. How can we classify them?
From Fig. 2(a,b) one notices that initially (for rather small $`B`$ values) the parabola $`E=k^2/2`$ separates the region where only states with positive group velocity exist, from the region where also negative velocity states are present. This is due to the fact that the value of this parabola equals the barrier height between the two parabolic wells for the corresponding $`k`$-value. When the energy exceeds this barrier, the shape of the wavefunction is not determined anymore by the separate parabolas, but by the overall composite well width. For decreasing $`k`$ the well is squeezed, and thus all the energy levels are pushed upwards, resulting in a positive group velocity. Although this is not an exact rule which cannot be extended rigorously throughout the $`B<B_0`$ regime, it nevertheless provides insight into the $`k`$-values (or $`B`$ values) for which these negative velocities states arise.
Inspection of the wavefunctions shows that there is a feature which marks the negative velocity states, and which relates indirectly to the presence of the different potential wells. It turns out that if the wavefunction or its first derivative exhibits a dip at some $`x`$ which satisfies $`\varphi (x)\varphi ^{\prime \prime }(x)>0`$ and $`\varphi ^{}(x)=0`$, or $`\varphi ^{}(x)\varphi ^{\prime \prime \prime }(x)>0`$ and $`\varphi ^{\prime \prime }(x)=0`$ and the condition that $`\varphi (x)0`$, then the state has a non-classical negative velocity. This is true for every $`k`$ value, as long as $`B<B_0`$. This is illustrated in Fig. 8 where we plot the wavefunction for $`k=1.5`$, $`n=2`$, with background magnetic field $`B/B_0=0.68`$ and $`0.73`$. The above dip in the wavefunction or its derivative (indicated by the dashed circle in Fig. 8) is a result of the different potential wells, which have their separate influence on the shape of the wavefunction, and therefore hamper the matching. The difference in $`\varphi (x)`$ being zero (or not), can be interpreted as a generalization of matching the individual states in an asymmetric (symmetric way) when the Landau states are degenerate at $`k\mathrm{}`$.
## VI Time dependent classical Interpretation
One can make different attempts to link a classical picture to quantum mechanics. Often the comparison starts with a schematic classical picture which is then supported by comparing the quantum mechanical probability density with the classical one, obtained through calculation of the classical electron trajectory solving Newtons equation. For a 1D problem one can also verify the classical motion by inspection of the velocity parallel to the edge. For a cilindrical symmetric problem, the classical electron motion can be inferred from the magnetic moment or the circular current distribution of the electron state. In this paper, a quantitative comparison was made by use of a quantum mechanical velocity parallel to the 1D magnetic field discontinuity. In the following, we will try a different approach where we will construct time dependent states, and in doing so we will introduce another feature, i.e. the oscillation frequency perpendicular to the magnetic edge.
### A $`B=0`$
We already mentioned before that the solutions for this kind of problem are the parabolic cylinder functions $`\varphi (x)=D_{E\frac{1}{2}}\left[\sqrt{2}\left(xk\right)\right]`$, matched in such a way that we have symmetric and anti-symmetric wavefunctions as is shown for $`k=2`$ in Fig. 1(a). At $`k=\mathrm{}`$, these symmetric and antisymmetric states are twofold degenerate (see the two wavefunctions corresponding to the solid square in Fig. 1(a)). Due to this degeneracy any linear combination of these states is also an eigenstate. If we take the following linear combination
$$\begin{array}{c}|m_+=\frac{1}{\sqrt{2}}\left(|2m+|2m+1\right),\hfill \\ |m_{}=\frac{1}{\sqrt{2}}\left(|2m|2m+1\right),\hfill \end{array}$$
(29)
we arrive at the well known Landau states, i.e. wavefunctions of electrons located in two different homogeneous magnetic field profiles. One electron is moving clockwise, while the other is moving counterclockwise. For decreasing $`k`$ this degeneracy is lifted. Although taking linear combinations of states with a different energy yields a time dependent solution, we will extrapolate this picture towards all the other states. We choose a new orthonormal but time dependent basis:
$`|m_+`$ $`=`$ $`\left(e^{iE_{2m}t}|2m+e^{iE_{2m+1}t}|2m+1\right)/\sqrt{2}`$ (30)
$`|m_{}`$ $`=`$ $`\left(e^{iE_{2m}t}|2me^{iE_{2m+1}t}|2m+1\right)/\sqrt{2}`$ (31)
with $`E_{m_+}=E_m_{}=(E_{2m}+E_{2m+1})/2`$. The resulting energy spectrum is shown in Fig. 2(a) by the dashed curves. The corresponding velocities are plotted in Fig. 9(a). For every branch there are two states $`|m_+`$ and $`|m_{}`$.
With these new quantum states much better agreement is obtained with the corresponding classical results (dotted curves in Fig. 9(a)). Because of the addition of the two eigenstates the negative velocity almost disappeared. Only the lowering of the energy, as was mentioned in the limiting case (i.e. $`k\mathrm{}`$), results in a small negative velocity, which can’t be understood even in this picture. Also the boundary which indicates when classical states propagate is in much better agreement now.
Since we now have time dependent states, we can calculate a new feature: the oscillation frequency $`\omega _x`$ in the $`x`$-direction. The time dependent probability densities of the $`|m_+`$ and $`|m_{}`$ states have the following form:
$`|m_+|x(t)|^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(|2m|x|^2+|2m+1|x|^2`$ (33)
$`+2\mathrm{cos}[\omega _xt]2m|x2m+1|x),`$
$`=`$ $`|m_{}|x(t+\pi /\omega _x)|^2,`$ (34)
where $`\omega _{x,m}=(E_{2m+1}E_{2m})/\mathrm{}`$ is the quantum mechanical oscillator frequency in the $`x`$-direction.
Classically we can calculate this frequency $`\omega _x`$ again, using simple geometrical considerations, which results in
$$\omega _x(k)=\frac{\pi }{2\mathrm{arccos}(k/v_\phi (k))}.$$
(35)
Both results are plotted in Fig. 9(b), and we obtain reasonably good agreement between the quantum (solid curves) and classical (dotted curves) results. Notice that for $`\left|k\right|>k^{}`$, classically $`\omega _x=0`$, which means that the electron does not oscillate between the two different magnetic field regions (i.e. it is not a snake orbit state), but it oscillates in a homogeneous magnetic field and consequently we obtain the time independent eigenstates corresponding to the Landau levels.
Of course this approach is only useful if proper linear combinations are possible. Unfortunately this is not the case when a background magnetic field $`BB_0`$ is applied.
### B $`B>B_0`$
The above approach is also fruitful in the case when $`0<B^l<B^r`$. We can again add adjacent levels, two by two, similar as described before. We can repeat exactly as was done before, and we arrive again at the time dependent states of Eq. (31). The energy spectrum of these states when $`B=3B_0/2`$ is shown in Fig. 2(d), by dashed curves. From Fig. 10(a), we notice that the classical velocity is in better agreement then before, since the amplitude of the whiggles is lowered, due to the summation. The quantum mechanical oscillation frequency in the $`x`$-direction is again given by $`\omega _{x,m}=(E_{2m+1}E_{2m})`$ and plotted in Fig. 10(b). We notice that since there are no degenerate states, we always have oscillating electrons. For $`k\mathrm{}`$ the electron oscillates with frequency $`\omega _x=2.5\omega _c`$, while for $`k\mathrm{}`$ the electron oscillates with frequency $`\omega _x=0.5\omega _c`$, i.e. the electrons circle around in their seperate homogeneous magnetic fields. This can also be seen from the classical oscillation frequency in the $`x`$-direction, which is given by
$$\omega _x(k)=\pi \left[\frac{\mathrm{arccos}(k/v_\phi (k))}{B^l}+\frac{\mathrm{arccos}(k/v_\phi (k))}{B^r}\right]^1.$$
(36)
Notice that also here whiggles in $`v_y`$ are present (see Fig. 10(b)) which are not present in the classical results. It is clear that proper linear combinations can always be made, as long as$`B>B_0`$.
## VII Conclusions
We studied the electron states near discontinuities in the magnetic field. Different 1D magnetic field profiles, i.e. steps, were considered. The quantum mechanical energy spectrum was obtained and the group velocity of the states was calculated. Their corresponding classical orbits were found and the propagating states which are located at the magnetic field discontinuity correspond to snake orbits. Quantum mechanical magnetic edge states were found which move along the magnetic field step in opposite direction to the classical snake orbits and which cannot be understood classically. We were able to construct non stationary quantum mechanical states which closely approximate the classical solution for the symmetrical case $`B^l=B^r`$ and for the more general case $`B^r>B^l>B_0`$.
###### Acknowledgements.
This work was partially supported by the Inter-university Micro-Electronics Center (IMEC, Leuven), the Flemish Science Foundation (FWO-Vl), BOF-GOA and the IUAP-IV. J.R. was supported by “het Vlaams Instituut voor de bevordering van het Wetenschappelijk & Technologisch Onderzoek in de Industrie” (IWT) and F. M. P. is a research director with the FWO-Vl. We acknowledge fruitful discussions with A. Matulis, P. Vasilopoulos, S. Badalian and J. A. Tyszynski. |
warning/0002/hep-th0002150.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Classification of type II BPS states with different supersymmetries have been discussed in several papers. They form $`U`$-duality multiplets of stings in various dimensions. The corresponding black hole entropies are also given by duality invariant expressions. However, its statistical understanding requires the degeneracy of states. It has been suggested that string networks may provide the degeneracy of such states, when considered on tori.
Network configurations have also been an important topic of discussion from various other angles in type II string theory, $`N=4`$ gauge theories and non-commutative geometry. In string theory, they provide $`1/4`$, $`1/8`$ and other lower supersymmetric states. Moreover type IIB planar string networks which end on D3-branes represent the $`1/4`$ BPS dyonic states of $`N=4`$ gauge theories in four dimensions. Network-like configurations, have also appeared in other supersymmetric field theories, but are likely to have connections with those mentioned above. Recently, extension of network configurations to strings carrying 1-form electric charges (per unit length) and currents was also presented. These can be of interest from the point of view of cosmological applications.
In this paper we study the application of string networks to $`D=5`$ BPS states. It is known that one can have $`1/8`$ supersymmetric particle-like states in $`D=5`$. We give a string network representation of such states, by compactifying periodic $`D=8`$ non-planar networks of an $`SL(3,Z)`$-multiplet of type II strings on $`T^3`$. We also generalize the results to the $`SL(5,Z)`$ U-duality in seven dimensions.
Our exercise can also be used to write down mass formula of $`1/8`$ BPS states, in certain world-volume theories of 2-branes. These branes are themselves identified as $`U`$-duality branes obtained from toroidal compactification of ten-dimensional branes. For the eight dimensional case, one notices that there exist 2-branes which are invariant under the $`SL(3,Z)`$ part of $`U`$-duality (for the purpose of this paper, $`SL(3,Z)`$ is the only relevant part of $`U`$-duality in eight dimensions). Masses of states, which can be identified as a string-junction ending on these branes can then be found from the above exercise. Such an analysis for D3-branes has been performed in great detail, including for the case of non-abelian gauge groups etc.. We expect that, same should be possible for these $`U2`$-branes as well.
We now start by describing the periodic non-planar $`(p,q,r)`$ string networks of our interest, built out of basic structures as shown in figures-1. For convenience, we first consider the string networks, whose basic building blocks are 4-string junctions as in figure-1(a). The existence of such a junction can be seen from the general structure of non-planar string networks of . In the construction of , the basic building blocks of the networks are 3-string junctions whose 3-prongs lie in a specific two dimensional plane in a three dimensional space, now identified as $`T^3`$. However, different junctions, including the adjacent ones can have their 3-prongs in different two dimensional planes, giving them a non-planar form (see figure-1(b)). Then by shrinking the length of the intermediate links of such adjacent junctions, which is a free parameter in these BPS constructions, one gets a 4-string junction. Such objects have also been studied in .
A periodic structure of such 4-string junctions can then be constructed as a three dimensional generalization of the string network lattice in . However one now needs four strings with $`SL(3)`$ charges: $`(p_I,q_I,r_I)`$, $`I=1,2,3,4`$ to construct such 4-string junctions. By fixing the lengths of these string-links to $`l_I`$, $`(I=1,2,3,4)`$, and by imposing charge conservation condition on junctions: $`_{I=1}^4p_I`$ = $`_{I=1}^4q_I`$ = $`_{I=1}^4r_I`$ = $`0`$, one obtains a periodic lattice. Although we do not present a pictorial representation of such 3-D periodic networks, their existence is guaranteed from the existence of the three dimensional lattice vectors $`\stackrel{}{a}`$, $`\stackrel{}{b}`$ and $`\stackrel{}{c}`$ given below in terms of the ‘link-vectors’ $`\stackrel{}{l}_I`$. These link-vectors themselves are given by the lengths of the prongs mentioned above and their orientation is given as in in order to preserve $`1/8`$ supersymemtry. More precisely, these orientations for a string with $`SL(3,Z)`$ quantum numbers $`(P_I)_i(p_I,q_I,r_I)`$ are given in terms of components $`(X_I)_a`$, $`(a=1,2,3)`$ of a vector in real space (now identified as $`T^3`$):
$$\stackrel{}{V_I}=(X_I)_a\widehat{e}_a,(a=1,2,3),$$
(1)
with $`(X_I)_a`$ given by:
$$(X_I)_a=(\lambda ^1)_{ai}(P_I)_i.$$
(2)
$`\widehat{e}_a`$’ in this paper always denotes orthogonal set of unit vectors in $`T^3`$, although its index $`\mathrm{`}a^{}`$ is chosen to be same as that of an internal $`SO(3)`$ vector. $`\lambda ^1`$ in the above equation denotes the vielbein corresponding to the $`SL(3)/SO(3)`$ moduli:
$`G=\left(\begin{array}{cc}g+a^Tae^\varphi & e^\varphi a^T\\ ae^\varphi & e^\varphi \end{array}\right),`$ (3)
with $`g`$ being a $`2\times 2`$ matrix:
$`g=\left(\begin{array}{cc}e^{(\varphi +\alpha )}+\chi ^2e^\alpha & e^\alpha \chi \\ \chi e^\alpha & e^\alpha \end{array}\right).`$ (4)
Then one has
$`\lambda ^1=\left(\begin{array}{ccc}e^{(\varphi +\alpha )/2}& \chi e^{(\varphi +\alpha )/2}& e^{(\varphi +\alpha )/2}a_1+\chi e^{(\varphi +\alpha )/2}a_2\\ 0& e^{\alpha /2}& e^{\alpha /2}a_2\\ 0& 0& e^{\varphi /2}\end{array}\right).`$ (5)
To summarize these definitions, $`P_I`$’s defined above are internal $`SL(3,Z)`$ vectors associated with the $`I`$’th string, whereas $`X_I`$’s are internal $`SO(3)`$ vectors constructed by contracting $`P_I`$’s with the vielbein $`\lambda `$. This $`SO(3)`$ is the maximal compact subgroup of $`SL(3)`$. Finally $`\stackrel{}{V}_I`$’s in eqn.(1) are vectors in $`T^3`$, due to their dependence on unit vector $`\widehat{e}_a`$. Identification of its components with those of $`X_I`$’s in eqns. (1), (2) is a property of the string networks, as the spatial and internal orientations of the links in a network are always aligned in a specific manner.
Major exercise now is to start from the expression of the mass associated with the above 4-string junction defining the unit cell of a periodic network lattice and to show that these can be rewritten in terms of three independent $`SL(3,Z)`$ charges $`(P_I)_i(p_I,q_I,r_I)`$, $`(I=1,2,3)`$, $`SL(3)/SO(3)`$ moduli $`G`$, and three-dimensional vectors: $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$ defined in terms of the lengths $`l_I`$’s of the four legs of the string junction as well as the unit vectors along these legs, $`\widehat{n}_I\stackrel{}{V}_I/|V_I|`$:
$$\stackrel{}{a}=\stackrel{}{l}_1\stackrel{}{l}_4,\stackrel{}{b}=\stackrel{}{l}_2\stackrel{}{l}_4,\stackrel{}{c}=\stackrel{}{l}_3\stackrel{}{l}_4.$$
(6)
In fact, as we will see below various combination of $`(\stackrel{}{a},\stackrel{}{b},\stackrel{}{c})`$ provide additional moduli in the lower dimensional theory, after quantum numbers $`(p_4,q_4,r_4)`$ are eliminated in favor of the remaining ones, using charge conservation conditions.
Technical non-triviality of our exercise, with respect to the one performed in for the planar network, is in dealing with the 3-dimensional problem in our case, compared to the 2-dimensional one in . To perform this exercise explicitly, we first consider the case when $`SL(3)/SO(3)`$ moduli, $`G`$, have a diagonal form. It will be observed that the final expression that we derive, easily generalize to the most general moduli as well. For the diagonal case, $`G`$ has a form: $`G=diag.(e^{\varphi +\alpha },e^\alpha ,e^\varphi )`$. Moreover for this case, the string tension is given as: $`T_I=|X_I|=[e^{(\varphi +\alpha )}p_I^2+e^\alpha q_I^2+e^\varphi r_I^2]^{\frac{1}{2}},(I=1,2,3,4)`$.
We now use the above expressions to compute the mass of the BPS state, given by the string network configuration built by the above 4-string junctions. It is given by
$$m_{BPS}^2=(l_1T_1+l_2T_2+l_3T_3+l_4T_4)^2.$$
(7)
Now, to eliminate the lengths of the link-vectors of the strings in favor of the generating vectors $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$, we use expressions of various scalar and vector combinations formed from these by taking their dots and cross products: $`a^2=l_1^2+l_4^22\stackrel{}{l}_1.\stackrel{}{l}_4,\stackrel{}{a}.\stackrel{}{b}=\stackrel{}{l}_1.\stackrel{}{l}_2\stackrel{}{l}_1.\stackrel{}{l}_4\stackrel{}{l}_2.\stackrel{}{l}_4+l_4^2,\stackrel{}{a}\times \stackrel{}{b}=\stackrel{}{l}_1\times \stackrel{}{l}_2\stackrel{}{l}_1\times \stackrel{}{l}_4\stackrel{}{l}_4\times \stackrel{}{l}_2`$ etc.. Moreover, these expressions can be rewritten in terms of quantum numbers $`(p_I,q_I,r_I)`$, $`(I=1,2,3)`$, moduli fields ($`\varphi `$, $`\alpha `$), lengths of the links $`l_I`$, and their string-tensions $`T_I`$, by using relations:
$`\stackrel{}{l}_1\times \stackrel{}{l}_2={\displaystyle \frac{l_1l_2}{T_1T_2}}[\widehat{e}_1(q_1r_2q_2r_1)e^{(\varphi +\alpha )/2}+`$ (8)
$`\widehat{e}_2(r_1p_2r_2p_1)e^{\alpha /2}+\widehat{e}_3(p_1q_2p_2q_1)e^{\varphi /2}],`$ (9)
and two other expressions obtained by taking cyclic permutations in indices $`(1,2,3)`$ and $`(p,q,r)`$. Similarly,
$`\stackrel{}{l}_1\times \stackrel{}{l}_4={\displaystyle \frac{l_1l_4}{T_1T_4}}[\widehat{e}_1(q_1(r_2+r_3)(q_2+q_3)r_1)e^{(\varphi +\alpha )/2}+`$ (10)
$`\widehat{e}_2(r_1(p_2+p_3)(r_2+r_3)p_1)e^{\alpha /2}+\widehat{e}_3(p_1(q_2+q_3)(p_2+p_3)q_1)e^{\varphi /2}],`$ (11)
and again two others obtained by the above cyclic permutations.
With the help of above expressions, and after some algebra, one can show that the mass of the BPS state, after $`T^3`$ compactification can be written as:
$`m_{BPS}^2=[(\stackrel{}{V}_1.\stackrel{}{V}_1)a^2+(\stackrel{}{V}_2.\stackrel{}{V}_2)b^2+(\stackrel{}{V}_3.\stackrel{}{V}_3)c^2+2(\stackrel{}{V}_1.\stackrel{}{V}_2)(\stackrel{}{a}.\stackrel{}{b})+2(\stackrel{}{V}_1.\stackrel{}{V}_3)(\stackrel{}{a}.\stackrel{}{c})+`$ (12)
$`2(\stackrel{}{V}_2.\stackrel{}{V}_3)(\stackrel{}{b}.\stackrel{}{c})]+2[(\stackrel{}{a}\times \stackrel{}{b}).(\stackrel{}{V}_1\times \stackrel{}{V}_2)+(\stackrel{}{a}\times \stackrel{}{c}).(\stackrel{}{V}_1\times \stackrel{}{V}_3)+(\stackrel{}{b}\times \stackrel{}{c}).(\stackrel{}{V}_2\times \stackrel{}{V}_3)],`$ (13)
$`m_1^2+m_2^2,`$ (14)
where $`m_1^2`$ and $`m_2^2`$ correspond to the terms in the two square brackets in eqn.(14). This equation is one of the main result of this paper. It gives the BPS mass in terms of nine integers $`(p_I,q_I,r_I)`$’s, moduli $`(\varphi ,\alpha )`$ (through their appearance in $`\stackrel{}{V_I}`$), as well new set of moduli formed out of $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$. The generalization of the result, to the case when the full set of $`SL(3)/SO(3)`$ moduli are turned on, is straight-forward. In that case, mass formula remains same as (14). However $`\stackrel{}{V}_I`$’s and $`X_I`$’s involve general $`SL(3)/SO(3)`$ moduli through their dependence on the vielbein in eqn.(5).
We now show that the above mass formula has an $`SL(3,Z)_U\times SL(3,Z)_u`$ symmetry. The first $`SL(3,Z)`$ is essentially the $`U`$-duality symmetry of type II strings in eight dimensions. The second $`SL(3,Z)`$ comes from the compactification of the network on $`T^3`$.
We first show the $`SL(3,Z)\times SL(3,Z)`$ invariance of the terms in the first square bracket in eqn.(14), identified as $`m_1^2`$. These terms can be rewritten as:
$$m_1^2=P^TMP,$$
(15)
where $`P`$ is $`9\times 1`$ column vector with entries:
$`P=\left(\begin{array}{c}P_1\\ P_2\\ P_3\end{array}\right),`$ (16)
and $`M`$ is a matrix:
$`M=\left(\begin{array}{ccc}a^2G^1& (\stackrel{}{a}.\stackrel{}{b})G^1& (\stackrel{}{a}.\stackrel{}{c})G^1\\ (\stackrel{}{b}.\stackrel{}{a})G^1& b^2G^1& (\stackrel{}{b}.\stackrel{}{c})G^1\\ (\stackrel{}{c}.\stackrel{}{a})G^1& (\stackrel{}{c}.\stackrel{}{b})G^1& c^2G^1\end{array}\right).`$ (17)
One can then write down the action of two $`SL(3)`$’s mentioned above on charges and moduli, including the ones constructed out of vectors $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$. $`SL(3)_U`$ has the identical action as in eight dimensions. This leaves any $`T^3`$ vectors such as $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$ etc. invariant and acts on $`P`$ through a diagonal action on $`P_I`$’s as:
$$P_I\mathrm{\Lambda }_UP_I,G^1\mathrm{\Lambda }_{U}^{1}{}_{}{}^{T}G^1\mathrm{\Lambda }_U^1.$$
(18)
with $`\mathrm{\Lambda }_U`$ being $`SL(3,Z)`$ matrices. Second symmetry, namely $`SL(3)_u`$ acts on $`P`$ as:
$`\left(\begin{array}{c}P_1\\ P_2\\ P_3\end{array}\right)\mathrm{\Lambda }_u\left(\begin{array}{c}P_1\\ P_2\\ P_3\end{array}\right),`$ (19)
namely, it mixes the indices $`(1,2,3)`$ associated with $`SL(3)`$ charges $`(p,q,r)`$ of various strings among themselves. In addition, one has
$`\stackrel{}{A}\left(\begin{array}{c}\stackrel{}{a}\\ \stackrel{}{b}\\ \stackrel{}{c}\end{array}\right)\mathrm{\Lambda }_{u}^{1}{}_{}{}^{T}\left(\begin{array}{c}\stackrel{}{a}\\ \stackrel{}{b}\\ \stackrel{}{c}\end{array}\right).`$ (20)
Due to the action of the symmetry group defined above, $`9\times 9`$ moduli matrix in the compactified theory $`(M)`$, transforms under $`SL(3)`$’s as $`M(\mathrm{\Lambda }_U^{1T}I_3)M(\mathrm{\Lambda }_U^1I_3),M(I_3\mathrm{\Lambda }_u^{1T})M(I_3\mathrm{\Lambda }_u^1)`$.
We have therfore shown an explicit $`SL(3,Z)_U\times SL(3,Z)_u`$ invariance of the first part of the BPS mass formula (15). We also observe that by factoring out the volume of the polyhedron formed out of vectors $`(\stackrel{}{a},\stackrel{}{b},\stackrel{}{c})`$ from the matrix $`M`$ in (17), it can be identified with a matrix parameterizing $`[SL(3/SO(3)]\times [SL(3)/SO(3)]`$ moduli.
We now see that the terms in the second square bracket of (14), namely $`m_2^2`$, are also invariant under the transformations (18), (19) and (20). We first consider $`SL(3)_U`$, after writing vectors $`\stackrel{}{V}_I`$’s as: $`\stackrel{}{V}_I=T_I\widehat{n}_I,(I=1,2,3)`$, with $`T_I`$ being the tensions of the strings and $`\widehat{n}_I`$ being the unit vectors along them. $`SL(3,Z)_U`$ invariance of $`m_2^2`$ then follows from the fact that it acts on various quantities inside second square bracket in (14) only through terms in the expressions of string tensions.
$`SL(3,Z)_u`$ symmetry of $`m_2^2`$ is also clear by noticing that although $`\stackrel{}{V}_I`$’s are spatial (or $`T^3`$) vectors, they transform under $`SL(3,Z)_u`$ due to its action on quantum numbers $`p_I,q_I,r_I`$’s in a similar manner as $`P_I`$’s mentioned above in (19). Then, using the definition $`\stackrel{}{A}`$ as in (20), the invariance of $`m_2^2`$ can be seen by writing it as:
$$m_2^2=(\stackrel{}{A}_I\times \stackrel{}{A}_J).(\stackrel{}{V}_I\times \stackrel{}{V}_J).$$
(21)
We have therefore shown an $`SL(3,Z)\times SL(3,Z)`$ invariance of the mass formula obtained from a periodic network of 4-string junctions in eight dimensions.
A similar analysis goes through for more general periodic string-networks constructed out of 4-prong structures as shown in figure-1(b). One now has lattice vectors defined as:
$$\stackrel{}{a}=\stackrel{}{l}_1(\stackrel{}{l}_4+\stackrel{}{l}_5),\stackrel{}{b}=\stackrel{}{l}_2(\stackrel{}{l}_4+\stackrel{}{l}_5),\stackrel{}{c}=\stackrel{}{l}_3\stackrel{}{l}_4.$$
(22)
The mass of the $`1/8`$ supersymmetric BPS state associated with the compactified string network is now given by the expression:
$$m_{BPS}^2=(\underset{i=1}{\overset{5}{}}l_IT_I)^2,$$
(23)
with lengths and tensions now being associated with the string-links in figure-1(b). This expression, after similar algebra as above, can now be written as:
$$m_{BPS}^2=P^TMP+\underset{I=1}{\overset{5}{}}(\stackrel{}{l}_I\times \stackrel{}{l}_J).(\stackrel{}{V}_I\times \stackrel{}{V}_J).$$
(24)
Then using the definitions of the lattice vectors in (22), and charge-conservations on vertices $`O_1`$ and $`O_2`$ in figure-1(b), one can show that the final $`1/8`$ BPS mass formula is once again given by eqns. (14), (15) and (21).
We also obtain the $`SL(2,Z)\times SL(2,Z)`$ invariant formula of by turning off appropriate moduli and charges. For example, when only nonzero $`SL(3)`$ charges are: $`(p_I,q_I)`$, $`(I=1,2)`$, then by setting $`\varphi =a_1=a_2=0`$ in (14) one reproduces exactly the same expression as in . This can be seen from the form of $`m_1^2`$ in (15), which reduces to the first term in eqn.(17) of in these limits. Moreover, for $`m_2^2`$ only one of the term in the second square bracket in (14) is nonzero and gives precisely the second term in eqn.(17) of .
We now comment on the connection of these results with $`U`$-duality in $`D=5`$. The full $`U`$-duality symmetry in $`D=5`$ is $`E_{6(6)}`$ and gauge charges are in its 27-dimensional representation. $`E_{6(6)}`$ however has an $`SL(6)`$ subgroup whose origin can be seen from the interpretation of the $`D=5`$ theory as $`T^5`$ compactified M-theory. This $`SL(6)`$, in turn, has an $`SL(3)\times SL(3)`$ subgroup which can be identified with $`SL(3)_U\times SL(3)_u`$ mentioned above. Nine charges represented by $`p_I,q_I,r_I`$, $`(I=1,2,3)`$ are within 27 of $`E_6`$, as can be seen by decomposing this under $`SL(6)`$ and identifying them to lie within $`\mathrm{𝟏𝟓}`$ of $`SL(6)`$.
The generalization of the result to $`SL(5,Z)`$ $`U`$-duality (in D=7) is also straightforward. One can analogously consider the case of periodic network lattice involving 6-string junctions (as well as other similar structures) and define a set of five vectors, $`\stackrel{}{\stackrel{~}{A}}_I`$ $`(I=1,..,5)`$ similar to $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$ defined earlier. Similarly one has a set of other five vectors, $`\stackrel{}{\stackrel{~}{V}}_I`$, whose components are given in terms of quantum numbers $`p_I,q_I,r_I..`$ etc., as well as $`SL(5)/SO(5)`$ moduli. The final mass formula, now with $`1/32`$ supersymmetry has a form:
$$M_{BPS}^2=(\stackrel{}{\stackrel{~}{V}}_I.\stackrel{}{\stackrel{~}{V}}_J)(\stackrel{}{\stackrel{~}{A}}_I.\stackrel{}{\stackrel{~}{A}}_J)+(\stackrel{}{\stackrel{~}{V}}_I\times \stackrel{}{\stackrel{~}{V}}_J).(\stackrel{}{\stackrel{~}{A}}_I\times \stackrel{}{\stackrel{~}{A}}_J).$$
(25)
To conclude the discussion of the compactified non-planar networks as lower dimensional BPS states, we like to point out that several other possibilities of network compactification can be discussed by restricting to smaller subgroups of $`U`$-duality. For example, one can construct planar periodic networks of $`(p,q,r)`$-strings in eight dimensions, by considering $`SL(2,Z)`$ subgroups of $`SL(3,Z)`$. One then has $`1/4`$ supersymmetric BPS states in six dimensions after compactifying these networks on $`T^2`$. It however remains to be seen whether one can obtain complete multiplets of the full duality symmetry, by combining various such possibilities of compactified networks.
We now discuss the application of the results to certain world-volume theories of branes, following a similar exercise for the case of planar IIB string networks. The planar IIB configurations are of interest from the point of view of $`1/4`$ BPS dyon solutions of $`N=4`$ gauge theory. These $`N=4`$ theories in turn are considered to be the linearized approximation of the world-volume theories of D3-branes that are invariant under the $`SL(2,Z)`$ duality of the IIB theory. Moreover electric and magnetic charges also transform under this $`SL(2,Z)`$. In eight dimensional type II theories, a similar role is played by 2-branes which are invariant under $`SL(3,Z)`$ and are known as $`U2`$-branes. From the point of view of branes, this $`SL(3)`$ acts on charges originating from three different components of $`N=1,D=10`$ gauge fields defining the world-volume theory. As an example, such charges can be identified in a 2-brane of this type by compactifying D4-branes on $`T^2`$. The $`SL(3,Z)`$ then acts on three charges, originating from the two internal components of the D4-brane gauge fields $`(A_3,A_4)`$ and a third one obtained by a Hodge-dualization of the three-dimensional gauge fields $`A_\mu `$. In a theory of parallel multi-branes, these fields are expected to form appropriate adjoint representations of the enhanced symmetries. Existence of the $`SL(3,Z)`$ symmetry on the world-volume can also be argued from the point of view of heterotic strings in $`D=3`$. The full duality symmetry of heterotic strings in $`D=3`$ is known to be $`O(8,24,Z)`$. The above $`SL(3,Z)`$, in this picture, then belongs to the $`SL(8,Z)`$ subgroup of $`O(8,24,Z)`$ , which transforms various components of ten-dimensional gauge fields, once again after Hodge-dualizations, in vector representations.
Then, to generalize the results of we consider a configuration of four such branes and above configuration of 4-string junction is formed by strings ending on these $`U2`$-branes. For example, in figure-1(a), points $`(A,B,C,D)`$ can be identified with the positions of these branes. The world-volume of the branes is orthogonal to the three dimensional space of strings and junctions. Vectors $`\stackrel{}{a},\stackrel{}{b},\stackrel{}{c}`$ described above parameterize the vacuum expectation values of the adjoint Higgs fields in the resulting $`N=8`$ supersymmetric theory. Now, to give a mass formula for such states while making connection with the work of , we choose special values for $`p_I,q_I,r_I`$ to be: $`(p_1,q_1,r_1)=(1,0,0)`$, $`(p_2,q_2,r_2)=(0,1,0)`$, $`(p_3,q_3,r_3)=(0,0,1)`$. In this case, BPS mass formula (14) reduces to
$`M^2=e^{(\varphi +\alpha )}|\stackrel{}{a}|^2+e^\alpha |\stackrel{}{b}|^2+e^\varphi |\stackrel{}{c}|^2+2[(\stackrel{}{a}\times \stackrel{}{b}).\widehat{e}_3e^{\varphi /2}+`$ (26)
$`(\stackrel{}{a}\times \stackrel{}{c}).\widehat{e}_2e^{\alpha /2}+(\stackrel{}{b}\times \stackrel{}{c}).\widehat{e}_1e^{(\varphi +\alpha )/2}].`$ (27)
Following , we now interpret this as the mass of a $`1/8`$ supersymmetric bound state in the world-volume theory described above. For this we define charges, similar to those in :
$$\stackrel{}{Q}_1=e^{(\varphi +\alpha )/2}\stackrel{}{a},\stackrel{}{Q}_2=e^{\alpha /2}\stackrel{}{b},\stackrel{}{Q}_3=e^{\varphi /2}\stackrel{}{c}.$$
(28)
Since the role of the couplings on the world-volume theory is played by the space-time moduli $`\varphi ,\alpha `$ etc., and $`(\stackrel{}{a},\stackrel{}{b},\stackrel{}{c})`$ define the vacuum expectation values of scalars in this world-volume theory, $`\stackrel{}{Q}_i`$’s above can be identified with the physical charges in this theory. The subscripts in expression of charges now have their origins in fields $`A_4,A_5,A_\mu `$ mentioned above, whereas vector sign above them can be interpreted to be along certain $`R`$-symmetry directions. Similar form of energy expressions, for $`1/2`$ BPS states involving these fields, were observed in . However it is of interest to verify these results directly from the point of view of world-volume gauge theories, following a similar exercise in $`D=4`$ in and to further examine the properties of such bound states.
Acknowledgement: I would like to thank Aalok Misra for useful communications. |
warning/0002/cond-mat0002336.html | ar5iv | text | # Untitled Document
TRAPPED BOSE–CONDENSATE IN GRAVITY FIELD
D.B.Baranov, V.S.Yarunin
Joint Institute for Nuclear Research, Dubna, 141980, Russia
yarunin@thsun1.jinr.ru
> The $`1D`$ and $`2D`$ Bose-condensation of trapped atoms in a gravitational field are considered. The deformation of the finite parabolic potential in this field is modeling via the combination of two rectangular $`1D`$ and $`2D`$ traps, for which the cut-off and the re-definition of spectrum are taken into account. A Bose-condensation $`T_c`$ shift by the gravity is calculated. A sign and a magnitude of it in a deformed trap depends on the order of including the gravitation field. The special choice of this order may describe three consistent Bose–condensations with different temperatures. These transitions may be associated with a transportation of a trap on the cycle (I) Earth–(II) Space–(III) Earth.
>
> PACS: 05.70.Jk, 67.40.-w
The news on the critical temperature $`T_c`$ of the Bose-Einstein condensation (BEC) entered the last decade. It is known, that the critical $`T_c`$ of the superfluid <sup>4</sup>He (the oldest – though inderect – example of BEC) is the same for any tube with a liquid helium. The decreasing of <sup>4</sup>He $`T_c`$ in porous glasses was found experimentaly and was explained theoreticaly , and this shift of a critical temperature depends on the parameters of a porous glasses . Meanwhile the calculation of BEC $`T_c`$ in an infinite volume for bosons with a weak interaction between them showes the increasing of $`T_c`$ .
As to the ideal gas, it was found , that in an external field, restricting an infinite volume with this gas, $`T_c`$ depends on the form of a field. The experimental discovery of alkali atoms BEC in magnetic traps gives the opportunity to discuss the same properties in finite systems. A finite volume assumes mesoscopic boundary conditions, which require a new formulation of the thermodynamical limit in theory . The decreasing of an ideal gas BEC $`T_c`$ was shown as a finite size correction $`N\mathrm{}N<\mathrm{}`$ for the ”trap version” thermodynamical limit $`\mathrm{}\omega /T0`$, $`N\mathrm{}`$ calculations.
Here we intend to look for a change of BEC critical temperature $`T_c`$ via the influence of a gravitational field on $`T_c`$ of a BEC gas in traps. A problem of the gravity influence on the critical temperature of <sup>4</sup>He arises in the connection of Space experiments on this subject, and the same experiments on BEC in trapped atomic gases are expected in the nearest future . So it is interesting to predict theoretically the contribution of gravity field to the critical BEC temperature in traps. We deal with the $`1D`$ and $`2D`$ theory of spin-less non-interacting atoms in a finite trap. The basic parameters are taken from the BEC of alkali atoms experiments . We start with the isotropic parabolic trap with a potential barrier $`U_0`$, frequency $`\omega `$ and a height $`h`$ and introduce two rectangular traps in order to describe the main part of gravity effect via the deformation of a trap potential.
The potentials of a parabolic trap are different for the cases of a zero- and non-zero gravity field in direction $`z`$. The potential inside of a trap may be written as
$$U(z)=U_g+(m\omega ^2/2)\left(z+\mathrm{\Delta }\right)^2,h/2<z<h/2,\mathrm{\Delta }=g/\omega ^2,$$
(1)
where $`m`$ is the mass of an atom, $`g`$ is the free fall acceleration constant. The dashed line is an initial parabolic potetial, a solid line is a parabolic potential, deformed by a gravity field. The gravity field shifts the minimum of potential from its center $`z=0`$ by the value $`\mathrm{\Delta }`$ and down by the value $`U_g=mg^2/2\omega ^2`$. The difference between new (shifted) potential barriers $`U_+`$ (on the right) and $`U_{}`$ (on the left)
$$U_\pm =\frac{m\omega ^2}{2}\left(\mathrm{\Delta }\pm \frac{h}{2}\right)^2+U_g,U_0U_{}=U_+U_0=\frac{1}{2}mgh$$
is equal to $`mgh`$. It follows from these formulas that the condition of macroscopic stability of the trap in gravity field is $`g<\omega ^2h/2`$. It is fulfilled for the experiments and the last inequality will be satisfied for the frequencies more than 100 $`Hz`$. For all that the quantum dynamics of atoms undergoes nonperturbate gravitational disturbance because of the big shift of the trap potential $`\left(U_g\mathrm{}\omega \right)`$ in comparison with its frequency.
One-particle boson eigenfunctions $`u_n`$ and their energies $`\epsilon _n>0`$ for the $`2D`$-case
$$\left\{^2+\frac{2m}{\mathrm{}^2}\left[ϵ_nU(r)\right]\right\}u_n=0,r=\{x,z\},u_n(r)=u_{n_x}(x)u_{n_z}(z),$$
(2)
$$ϵ_n=\frac{1}{2m}u_nu_ndr+U(r)u_nu_n𝑑r,n=\{n_x,n_z\}$$
are the starting points of the quantum analisys of any finite trap $`U(r)`$, a ground level of it is $`\epsilon _00`$. In order to introduce the Gibbs statistics we use the path-integral device, applied earlier for the BEC of interacting bosons.
The system of $`N`$ bosons (2) in a volume $`V`$ may be represented by the Hamilton function $`H`$
$$\widehat{H}H=𝑑r\left[\frac{\mathrm{}^2}{2m}\left(\psi ^{}(r,t)\psi (r,t)\right)+\psi ^{}U(r)\psi \right].$$
Here $`\psi ^{},\psi `$ are the path-integral trajectories with periodical boundary conditions on $`[0,\beta ]`$. Following , we separate the trajectories $`b_0^{},b_0`$ as the ”slow” variables relative to the ”fast” trajectories $`b_n^{},b_n,n0`$ for $`T<T_c`$
$$\psi ^{}(r,t)=\frac{1}{\sqrt{V}}\left(\underset{n0}{}b_n^{}(t)u_n(r)+b_0^{}(t)u_0(r)\right),|b_0|=\sqrt{N_0}1,$$
so that the function $`N_0(t)`$ is the path-integral image of the ”Bose-condensate” particles in Bogoliubov theory . The ”broken gauge symmetry” is known as a property of the bose-gas with the separated condensate fraction. Still the quasiclassical integral of motion $`\overline{N}=N_0+\overline{N_1}`$ for the system (2) may be expressed by the equation
$$\frac{d\overline{N}}{dt}\{\widehat{H},N_0\}+i[\widehat{H},\widehat{N}_1]=0,\widehat{N}_1=\underset{n0}{}b_n^+b_n.$$
(3)
Here $`N_1`$ is a number of non-condensate bosons, $`\{,\}`$ is the Poisson bracket in $`b_0,b_0^{}`$ classical amplitudes and $`[,]`$ is the quantum commutator via operators $`b_n^+,b_n^{}`$, $`n,n^{}0`$. The action $`S`$ with the kinetic term $`K`$ for $`N`$ non-interacting bosons in any volume $`V`$ looks like
$$S=_0^\beta (KH)𝑑t=_0^\beta \left[b_0^{}\frac{db_0}{dt}+\epsilon _0b_0^{}b_0+\underset{n0}{}\left(b_n^{}\frac{db_n}{dt}+\epsilon _0b_n^{}b_n\right)\right]𝑑t.$$
An effective action $`S_{ef}(\rho ,\mu )`$ for a condensate density $`\rho =N_0/V`$ and a chemical potential $`\mu `$ is defined via the partition function with constraint
$$Q=e^{S_{ef}(\rho ,\mu )}𝑑\rho 𝑑\mu =Sp\left[\mathrm{exp}(\beta H)\delta _{R,R_1+\rho }\right],R=N/V,R_1=N_1/V.$$
The equation (3) is the background for the definition of $`\mu `$, and by the use of Fourier decomposition for $`\delta _{R,R_1+\rho }`$ with $`\mu =iy/\beta `$ we represent $`Q`$
$$Q=Db_0^{}Db_0\underset{n0}{}Db_n^{}Db_n_\pi ^\pi 𝑑y\mathrm{exp}[iy(R_1+\rho R)+S/\mathrm{}]$$
as the integral over $`y`$, the functional integral over the ”slow” $`b_0,b_0^{}`$ and ”fast” functional integral trajectories $`b_n,b_n^{}^{}`$ with the periodical boundary conditions in \[$`0,\beta `$\]. The calculation of $`Q`$ over ”fast” variables gives the effective action $`S_{ef}`$ of the condensate
$$Q=Db_0^{}Db_0𝑑\mu \mathrm{exp}S_{ef},$$
$$S_{ef}=_0^\beta \left[b_0^{}\frac{db_0}{dt}+\rho (\mu \epsilon _0)\mu R\right]𝑑t\underset{n0}{}\mathrm{ln}\left(1\mathrm{exp}[\beta (\epsilon _n\mu )]\right).$$
It leads to the variational equations $`\delta S_{ef}(b_0,b_0^{},\mu )=0`$ in the form
$$\frac{S_{ef}}{\beta \mu }=\rho +\frac{1}{V}\mathrm{\Sigma }R=0,\mathrm{\Sigma }=\underset{n_x,n_z0}{}f(ϵ_{n_x,n_z},\beta ,\mu ),$$
(4)
$$f(ϵ_{n_x,n_z},\beta ,\mu )f=\frac{1}{\mathrm{exp}[\beta (\epsilon _n\mu )]1},n=\{n_x,n_z\},$$
$$\frac{S_{ef}}{\beta b_0^{}}=\frac{db_0}{dt}b_0(\epsilon _0\mu )=Lb_0=0,\mu =\epsilon _0kT\mathrm{ln}\left(1+\frac{1}{N_0}\right).$$
(5)
Equation (4) means the balance of condensate and non-condensate particles. The solution of equation (5) follows via the calculation
$$(\beta L)^1=[\mathrm{exp}\beta (ϵ_0\mu )1]$$
of the operator $`L^1`$ with the periodical boundary conditions. The equations (4,5) determine an BEC solution for the model (2,3) below the critical temperature for the case of translation-noninvariant systems. So, the critical temperature $`T_c`$ for a finite-size trap is determined by the equations $`\mathrm{\Sigma }(T_c)=N`$, $`\rho =0`$, while at the temperarure $`T0`$ the equations $`\rho R,R_10,\mu \epsilon _0`$ are present. Note, that if we take $`b_0^{},b_0`$ as $`constants`$, we get only $`T=0`$ limit of solutions for $`\mu `$, it is a property of the ideal Bose-gas. The formulas (3-5) are valid for any trap $`U(r)`$. We apply them for a parabolic trap in formulas (6-8).
The best way to catch a mesoscopic effect is to look for the shift of a critical temperature $`T_c`$ of BEC in a trap. It is obvious, that a finite size $`h`$ for a parabolic trap manifests itself in the following mesoscopic properties:
($`i`$) a finite number of levels $`n_{max}=U_0/\mathrm{}\omega `$ of the initial symmetrical parabolic trap,
($`ii`$) a diminishing of the left-side potential barrier $`U_0(U_{}U_g)<U_0`$ by the gravitational field. The re-definition of one-particle matrix elements $`ϵ_n`$ must be done in both these cases. It means that we must ($`i`$) cut-off an upper limit $`n<n_{max}`$ for barrier $`U_0`$ and calculate energies $`ϵ_n^0`$, as well as ($`ii`$) cut-of an upper limit in (2) for barrier $`U_{}`$ and calculate energies $`ϵ_n^\pm `$. So the further calculation will be done with the formula
$$N=N_0+N_1,N_1=\underset{n0}{\overset{n_{max}}{}}f_n,\epsilon _n=\{\begin{array}{cc}\epsilon _n^0,& g=0,\\ \epsilon _n^\pm ,& g0.\end{array},n=\{n_x,n_y\}.$$
(6)
of the full number $`N`$ of atoms in a trapped ideal bose gas. The level of $`n_{max}`$ defines the upper bound states of atoms. The numerical parameters of the experiments for a parabolic trap in an isotropic version may be reduced approximately to the values
$$U_010^9eV,h=2mm,U_+U_{}=mgh,\omega 10^{13}eV.$$
(7)
We suppose the critical $`3D`$ peak density $`R10^{14}cm^3`$ to be independent on time as the evaporation cooling is neglected here. The $`nD`$ scaling for the number of atoms in a parabolic trap with a volume $`V`$, over wich the system is confined, looks like
$$R|_{nD}=\frac{N}{V},V\left(\frac{m\omega ^2}{T_c}\right)^{n/2},n=1,2,3.$$
(8)
In order to symplify a problem we would like to represent the parabolic trap by the rectangular trap. It means, that we describe the deformation of the parabolic potential in the gravitation field by transforming the symmetrical rectangular trap with potential $`U_0`$ to the asymmetrical one with barriers $`U_+U_{}=mgh`$. Our purpose is to estimate the shift of a critical temperature in the asymmetrical rectangular trap relative the symmetrical one as a measure of the gravity field influence on BEC, tunneling effects are not considered. The values of $`R`$ in a $`1D`$ and $`2D`$ parabolic traps are $`10^5cm^1`$ and $`10^{10}cm^2`$ in the view of (8). We keep these $`R`$ the same for a rectangular traps. So a number of atoms in a $`1D`$ rectangular trap is $`N_1|_{1D}10^5cm^1\times 0.2cm210^4`$, and in a $`2D`$ rectangular trap is $`N_1|_{2D}10^{10}cm^2\times (0.2cm)^2410^8`$. The equation for a spectrum $`\epsilon _nE_k`$ for the symmetrical $`1D`$ and $`2D`$ traps looks like
$$tg(kh)=\frac{2k\mathrm{}}{p_k^0}\left[\frac{k^2\mathrm{}^2}{(p_k^0)^2}1\right]^1,E_k=\frac{\mathrm{}^2(k_x^2+k_z^2)}{2m},p_k^0=\sqrt{2m(U_0E_k)},k=\{k_x,k_z\}$$
and for the asymetrical (deformed) trap for $`k=k_z`$ looks like
$$tg(kh)=k\mathrm{}\left(\frac{1}{p_k^+}+\frac{1}{p_k^{}}\right)\left(\frac{k^2\mathrm{}^2}{p_k^+p_k^{}}1\right)^1,p_k^\pm =\sqrt{2m(U_\pm E_k)}.$$
The distances between the energy levels $`E_k`$ of an initial rectangular trap with the given parameters (7) are $`10^{19}eV`$ at a bottom and $`10^{13}eV`$ at a top barrier $`U_0`$. It has $`10^5`$ levels (contrary $`10^4`$ in experiments ), so the average energy distance between levels $`710^{14}eV`$ is the same as the frequency $`\mathrm{}\omega 610^{14}eV`$ of the paprabolic trap.
In order to get more systematic picture we start with 1D case with the $`R10^5cm^1`$ density. The critical temperature $`T_c^0|_{1D}=2.4110^{10}K`$ for an initial rectangular trap $`U_0`$ is found from the equation (6) with the use of calculated spectrum $`E_k`$
$$N_0=0,N_1^0|_{1D}=N|_{1D}=\underset{k_z0}{}f(T_c^0|_{1D},E_{k_z}^0)=210^4,k<k_{max}.$$
A suitable accuracy of calculations may be received for $`k_{max}1000`$ in the last sum in $`k_z`$ variable. The deformed rectangular trap $`U_{}10^8U_0`$ is determined via the parameters $`h`$ and $`U_+U_{}=mgh`$ in (7). It contains only $`10`$ levels, which are occupied by
$$N_1^\pm |_{1D}=\underset{k_z0}{\overset{10}{}}f(T_c^0|_{1D},E_{k_z}^0)=1.810^4<N|_{1D}$$
atoms. These atoms undergoes a new BEC transition in the system of re-defined levels with a critical temperature in a deformed trap $`T_c^\pm |_{1D}=2.2710^{10}`$, that is found from the equation
$$N_0=0,N_1^\pm |_{1D}=\underset{k_z0}{\overset{10}{}}f(T_c^\pm |_{1D},E_{k_z}^\pm )=1.810^4,k<k_{max}.$$
So the shift of the critical temperature is $`\mathrm{\Delta }(T_c)|_{1D}=T_c^\pm |_{1D}T_c^0|_{1D}=0.1410^{10}`$. The rest part $`0.210^4`$ of atoms becomes non–trapped (continuous spectrum).
As to the $`2D`$ case, the sum over the variables $`k_x,k_z`$ is taken within the constraint condition <sup>1</sup><sup>1</sup>1This condition is trivial for an initial trap, but rather effective for a deformed trap. $`E_k<E(n|_{max})`$, $`k=\{k_x,k_z\}`$. A temperature $`T^0|_{2D}=5.0510^7K`$ for an initial rectangular trap $`U_0`$ is found from the equation (6)
$$N_0=0,N_1|_{2D}=N|_{2D}=\underset{k_x,k_z0}{}f(T_c^0|_{2D},E_{k_xk_z}^0)=410^8,k<k_{max}.$$
The latter sum includes all the items of the partition function within the possible degeneracy of levels. A suitable accuracy of calculations may be received for $`k_{max}400`$ in the last sum in $`k_x,k_z`$ variables. The deformed rectangular trap $`U_{}10^8U_0`$ contains $`10`$ levels, which are occupied by
$$N_1^\pm |_{2D}=\underset{k_x,k_z0}{\overset{10}{}}f(T_c^0|_{2D},E_{k_xk_z}^0)=0.9610^8<N|_{2D}$$
atoms. Just like in $`1D`$ case, we look for a transition temperature in a deformed trap using an equation
$$N_0=0,N_1^\pm |_{2D}=\underset{k_x,k_z0}{\overset{10}{}}f(T_c^\pm |_{2D},E_{k_xk_z}^0).$$
The result of calculation is $`T_c^\pm |_{2D}=5.0310^7K`$, so the shift of a critical temperature is
$$\mathrm{\Delta }(T_c)|_{2D}=T_c^\pm |_{2D}T_c^0|_{2D}=0.0210^7K.$$
(9)
Now we describe the $`T_c`$ shift in a system with the same number of atoms, $`410^8`$ for example (for $`2D`$-case in our theory). The critical temperature for the symmetrical trap was found above as $`T_c^0|_{2D}=5.0510^7K`$. The critical temperature for the asymmetrical trap $`\left(T_c^\pm |_{2D}\right)^{}=2.0210^6K`$ is found from the equation
$$N_0=0,N_1|_{2D}=N|_{2D}=\underset{k_x,k_z0}{\overset{10}{}}f[\left(T_c^\pm |_{2D}\right)^{},E_{k_xk_z}^\pm ]=410^8,$$
$$\mathrm{\Delta }(T_c)=\left(T_c^\pm |_{2D}\right)^{}T_c^0|_{2D}=1.5210^6K.$$
(10)
Note, that $`\left(T_c^\pm |_{2D}\right)^{}`$ in (10) does not coinside with the $`T_c^\pm |_{2D}`$ in (9), the reason is that the processes (9) and (10) are different. This important fact will be discussed in the end of the paper.
The above calculations concern a rectangular trap. As to an initial parabolic trap, it has the same barriers $`U_\pm `$ with the same shift $`U_+U_{}=mgh`$, but differs in spectrum. It is reasonable to suppose, that the same scale of a BEC critical temperature shift will be valid for the parabolic trap in the gravitation field, so the qualitative correspondence between the critical temperatures $`T_c^0T|_{g=0}`$, $`T_c^\pm T|_{q0}`$ for a rectangular and parabolic traps is present. Meanwhile the tunneling effect is not reproduced by a rectangular trap. As it is seen from the formula (1), the macroscopic stability of a parabolic trap in a gravity field is expressed by the unequality $`\omega ^2>2g/h`$. It is clear, that this condition is trivial in the absence of gravity field or for the trap on the whole axe $`\mathrm{}<z<\mathrm{}`$. The finite size trap effect in the case of a deformed rectangular trap $`n_{max}10`$, attributed to the most noticable gravity shift of $`T_c`$, may be reproducrd for the parabolic trap via the condition
$$(\omega ^2h2g)+0.$$
This condition means, that the largest shift of $`T_c`$ may be noticed just before the destruction of the parabolic trap.
Let us turn in a conclusion to the non-ideal gas with an interaction $`G(r,r^{})`$ beteen atoms. In the 1D case the effective action in (4,5) for interacting bosons is generalized in the Hartree-Fock-Bogoliubov approximation as
$$S_{ef}(\rho ,\mu )=\beta \rho ^2\gamma _0V+\beta \mu (\rho R)V\beta \rho \epsilon _0+$$
$$\frac{1}{2}\underset{nn^{}0}{}\left[\beta (W_{nn^{}}\mu \delta _{nn^{}})4lnsh\frac{\beta E_{nn^{}}}{4}+\beta A_{nn^{}}\right],A_{nn^{}}=\frac{V\varphi _n^{}M_{nn^{}}m\varphi _n^{}}{E_{nn^{}}^2},$$
$$\rho =\frac{|b_0|^2}{V},m=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right),M=\left(\begin{array}{cc}W_{nn^{}}\mu \delta _{nn^{}}& 2\gamma _{nn^{}}\rho \\ 2\gamma _{nn^{}}\rho & W_{nn^{}}+\mu \delta _{nn^{}}\end{array}\right),$$
$$W_{nn^{}}=w_{nn}/V+2\rho \gamma _{nn^{}},E_{nn^{}}^2=(M|_{11})^2(M|_{12})^2,\varphi _n^{}=(b_0^{},b_0)(\rho \gamma _{0n}).$$
Matrix elements of $`G(z,z^{})`$ in $`\delta `$–approximation $`G(z,z^{})G\delta (zz^{})`$
$$\gamma _0=\frac{G}{2}𝑑zu_0^4,\gamma _{0n}=\frac{G}{2}𝑑zu_0^3u_n,\gamma _{nn^{}}=\frac{G}{2}𝑑zu_0^2u_nu_n^{}$$
are detrmined for the functions of harmonic oscillator in $`z`$ direction with the potential barrier $`(U_{}U_g)<U_0`$ at the left endpoint of the trap. Only the cut-off spectrum contribution (without the re-definition of $`ϵ_n`$) is taken into account. Equations $`\delta S_{ef}(\rho ,\mu )=0`$ for a trap with the constant density of atoms $`R=N/V`$ may be evaluated approximately by dividing the sum over full number of trap levels into two terms with $`nn_0`$ and $`nn_0`$. Parameter $`n_0100`$ is defined under the condition
$$n_0\mathrm{}\omega \rho \gamma _{n_0n_0}GN_0,N_0N,$$
(11)
following from the unequality
$$10^{13}eV\mathrm{}\omega GN_010^{11}eV,G=\frac{4\pi \mathrm{}^2a}{m},$$
where the energy of interaction between atoms is taken from the estimation for the scattering length $`a=4.9nm`$ with density of atoms $`R10^{10}cm^3`$. Thus, a density of condensate is written in the form $`\rho =\rho ^<+\rho ^>`$, where $`\rho ^<=\rho |_{nn_0}`$, $`\rho ^>=\rho |_{nn_0}`$. We can estimate the ratio $`Y`$ of Bose-condensate densities in a trap with $`\rho |_g`$ and without gravity $`\rho |_0`$
$$Y=\frac{\rho |_g}{\rho |_0}=\frac{(\rho ^<+\rho ^>)|_g}{(\rho ^<+\rho ^>)|_0}\frac{\rho ^<|_g}{\rho ^<|_0}\left(1\frac{\rho ^>|_0}{\rho ^<|_0}+\frac{\rho ^>|_g}{\rho ^<|_g}\right),\rho ^<|_{0,g}\rho ^>|_{0,g}.$$
In the case of a ”strong” gravity field $`(U_{}U_g)U_0`$ the inequality $`\rho ^<|_g>\rho ^<|_0`$ follows after the evaluations. In the case of a ”weak” gravity field $`(U_{}U_g)U_0`$ the equations $`\rho ^<|_g=\rho ^<|_0`$ and $`\rho ^>|_g>\rho ^>|_0`$ are valid. The inequalities
$$Y_{strong}=\frac{\rho ^<|g}{\rho ^<|_0}\left(1\frac{\rho ^>|_0}{\rho ^<|_0}\right)>1,Y_{weak}=\left(1\frac{\rho ^>|_0\rho ^>|_g}{\rho ^<|_0}\right)>1$$
are proved. Taking into account, that the large Bose-condensate density corresponds the large critical temperature, we note, that the gravitional field increases the critical temperature $`T_c`$ of Bose-condensation for non-ideal gas in a trap. The parameter $`n_0`$ in (11) corresponds the parameter $`n_{max}`$ in (6) as the border of BEC occupated states. The tunneling effects are also neglected here.
We would like to look for a correspondence between formulas (9,10) and the probable motion of a trap between the Earth and Space. We measure the $`\left(T_c^\pm |_{2D}\right)^{}=2.0210^6K`$ temperature for the $`N|_{2D}`$ atoms on the Earth in a deformed trap in the begining. Then the trap goes to Space, and the same number of atoms $`N|_{2D}`$ undergoes the phase transition (in a non–deformed trap) at the temperature $`T_c^0|_{2D}=5.0510^7K`$ with a shift (10). If the trap goes back to the Earth, some part of atoms becomes non–trapped due to its deformation. The rest part $`N_1^\pm |_{2D}`$ of them undergoes a phase transition at the temperature $`T_c^\pm |_{2D}=5.0310^7K`$, and a temperature shift is given by (9).
So we have three transition temperatures (with the same trap!) in different cases I, II, III
$$\left(T_c^\pm |_{2D}\right)^{}=2.0210^6KT_c^0|_{2D}=5.0510^7KT_c^\pm |_{2D}=5.0310^7K,$$
(12)
start on the Earth (I) $``$ in Space(II) $``$ finish on the Earth (III).
This is a picture for an ideal Bose–gas. The sign of $`T_c`$ shift for the non–ideal gas on the stage I$``$ II was performed in the last formulas of the paper, but the re–definition of levels was not taken into account. It seems, that the $`3D`$ calculation will increase the magnitudes of the critical temperatures in (12).
The mesoscopic nature of the problem considered above is seen in the strong dependence of the calculated $`T_c`$ shifts on the initial parameters, and the way of their association with experimental prototypes.
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warning/0002/hep-th0002203.html | ar5iv | text | # 1 Introduction
## 1 Introduction
There seem to be two kinds of field theories: one good and the other bad. Good theories are renormalizable theories, such as $`\varphi ^4`$ theory, QED, and QCD, which are well defined at all energy scales, and for which everything can be calculated in terms of a finite number of parameters. Bad theories are non-renormalizable theories, such as the four-Fermi weak theory and non-linear sigma model, which are well defined only below a finite cutoff energy (or momentum), and for which more and more parameters must be introduced as the order of approximations increases. Theories can be classified using the standard power counting rule: if the lagrangian of a theory contains only fields with dimension four or less, it is renormalizable, and otherwise it is non-renormalizable. This classification into renormalizable and non-renormalizable theories is so simple and quite popular, but we know it is wrong.
The physical meaning of renormalization and renormalizability cannot be understood without the renormalization group (to be abbreviated as RG) introduced by K. G. Wilson. Using the renormalization group we also distinguish two classes of theories, but this time one is the class of $`\varphi ^4`$-like theories, and the other is that of QCD-like theories.
A $`\varphi ^4`$-like theory or a non-asymptotic free theory is characterized by an IR fixed point of the renormalization group. If the relevant (or mass) parameters are taken to zero, the theory approaches the IR fixed point along a RG trajectory. With non-vanishing masses, the theory is driven away from the fixed point. The effects of irrelevant and marginal parameters become smaller along a RG trajectory: the effect of an irrelevant parameter, which is of order $`1`$ at the cutoff energy scale $`\mathrm{\Lambda }`$, is suppressed by a positive power of $`\frac{\mu }{\mathrm{\Lambda }}`$ at an arbitrary low energy scale $`\mu `$, but the effect of a marginal parameter of order $`1`$ at $`\mathrm{\Lambda }`$ is only suppressed as $`\frac{1}{\mathrm{ln}\frac{\mathrm{\Lambda }}{\mu }}`$ at $`\mu `$. We cannot take a true continuum limit of this kind of theories; as we take $`\mathrm{\Lambda }\mathrm{}`$, the effects of marginal parameters vanish, and we are left with a free massive theory.<sup>2</sup><sup>2</sup>2We assume that the fixed point is a free massless theory. This is called “triviality” in the literature. (See Fig. 1.)
In contrast the other class of theories, QCD-like or asymptotic free theories, admit a true continuum limit. The fixed point of a QCD-like theory is an UV fixed point. We take the relevant parameters to zero as we take the continuum limit. The continuum limit obtained this way is well-defined at all energy scales, and it is parameterized by mass and relevant marginal parameters which drive the theory away from the fixed point along the RG trajectories. With a cutoff large but finite, the theory differs from the continuum limit due to the effects of irrelevant parameters, which are suppressed by positive powers of $`\frac{\mu }{\mathrm{\Lambda }}`$.
Thus, any theory has good low energy behaviors as long as we fine-tune the mass parameters. Fine-tuning is necessary, since the natural mass scale is $`\mathrm{\Lambda }`$, and the mass parameter must be fine-tuned, typically to the order of $`\frac{\mu ^2}{\mathrm{\Lambda }^2}`$, to attain a finite physical mass much smaller than the cutoff.<sup>3</sup><sup>3</sup>3In pure QCD we only need to tune but not fine-tune the gauge coupling constant to zero, since it is only marginally relevant. With massive quarks, the quark masses must be fine-tuned to zero to the order $`\frac{\mu }{\mathrm{\Lambda }}`$. A theory is either QCD-like or $`\varphi ^4`$-like: if it is not QCD-like, it is $`\varphi ^4`$-like, and with a large but finite cutoff $`\mathrm{\Lambda }`$ and fine-tuning of mass parameters the theory describes interactions of order $`\frac{1}{\mathrm{ln}\frac{\mathrm{\Lambda }}{\mu }}`$ at a low energy scale $`\mu `$.
The above RG viewpoint was clearly stated in the first of refs. regarding the equivalence between the non-renormalizable Nambu-Jona-Lasinio (NJL) model and the renormalizable Yukawa theory. Despite their difference in appearance, the two theories are both $`\varphi ^4`$-like, and they describe the same physics if we ignore the irrelevant differences suppressed by positive powers of $`\frac{\mu }{\mathrm{\Lambda }}`$. Similarly, the $`O(N)`$ linear- and non-linear sigma models are both $`\varphi ^4`$-like, and they are equivalent up to differences suppressed by positive powers of $`\frac{\mu }{\mathrm{\Lambda }}`$. As these two concrete examples show, the appearance of a theory at the cutoff scale is misleading. It is the fixed point of the RG which dictates the low energy behaviors of a theory.
The purpose of the present paper is to extend the work of refs. and apply the above RG viewpoint to explore the possibility of constructing gauge theories using apparently non-renormalizable lagrangians without elementary gauge fields. The dynamical generation of gauge symmetries have already been discussed in the literature. Immediately after the work on the NJL model (the first of refs. ), purely fermionic construction of QCD was attempted in ref. . Even earlier, in generalizing the idea of non-linear realizations of symmetries, dynamical generation of gauge symmetries, called hidden local symmetries, was shown to be possible. A non-perturbative study of QCD constructed as an induced gauge theory has been also given in ref. .
Our work differs from these earlier works in two aspects: first we will study $`\varphi ^4`$-like gauge theories from the RG viewpoint given above. We are not introducing new theories. Rather we are showing that certain non-renormalizable theories, often perceived as undesirable, are really the same as perturbatively renormalizable gauge theories. Second we give a detailed analysis of the Ward identities. A renormalizable theory with a vector field is not necessarily a gauge theory. To be a gauge theory, Ward identities must be satisfied. We will show it possible to tune marginally irrelevant parameters to satisfy the necessary Ward identities. In this paper we only consider two theories with abelian gauge symmetries.
A comment is in order regarding the use of $`1/N`$ expansions in our work. We emphasize that all our results are supported by the RG viewpoint given above, and that they are valid for any $`N`$ starting from $`1`$. We use the $`1/N`$ expansions not to study the theories non-perturbatively. We are only interested in perturbation theory with respect to small coupling constants such as the fine structure constant and scalar self-coupling. For $`\varphi ^4`$-like theories, the so-called non-perturbative effects are all cutoff dependent: their contributions are suppressed by positive powers of $`\frac{\mu }{\mathrm{\Lambda }}`$ which we ignore in our study of non-renormalizable theories. Here we use the $`\frac{1}{N}`$ expansions to get a small coupling constant of order $`\frac{1}{\mathrm{ln}\frac{\mathrm{\Lambda }}{\mu }}`$ from loop corrections. Naïve perturbative expansions in powers of a bare coupling does not work, since the coupling is of order $`1`$.
This paper is organized as follows. In sect. 2 we review the equivalence of the $`O(N)`$ non-linear sigma model with the linear sigma model at energies much below the cutoff. This is to remind the reader how misleading the appearance of a lagrangian can be and to emphasize the usefulness of the $`1/N`$ expansions. In sect. 3 we study a non-renormalizable fermionic model with a current-current interaction and show its equivalence to the standard (massive) QED. This is followed by the study of a little more non-trivial scalar model which also has a current-current interaction in sect. 4. The model is shown to be equivalent to the (massive) scalar QED. Finally, the paper is concluded in sect. 5. For the convenience of the reader, we summarize the Ward identities of QED and scalar QED in appendix A. A table of integrals with a momentum cutoff is given in appendix B.
Throughout the entire paper we will work in the four dimensional euclidean space. Our convention is that the weight of a euclidean functional integral is given by $`\mathrm{exp}[S]=\mathrm{exp}[d^4x]`$ where $`S`$ is a euclidean action, and $``$ is a euclidean action density.<sup>4</sup><sup>4</sup>4We call it an action density rather than a lagrangian density to avoid potential confusion about signs.
## 2 Review of the $`O(N)`$ non-linear sigma model
We begin with a brief review of the equivalence between the $`O(N)`$ linear and non-linear sigma models. We wish to explain quantitatively how the apparently non-renormalizable non-linear sigma model can be physically equivalent to the perturbatively renormalizable linear sigma model. We first summarize the relevant results on the linear sigma model. The model is defined for $`N`$ real scalar fields $`\varphi ^I`$ ($`I=1,\mathrm{},N`$) by the following action density
$$=\frac{1}{2}_\mu \varphi ^I_\mu \varphi ^I+\frac{m^2}{2}\varphi ^I\varphi ^I+\frac{\lambda }{8N}\left(\varphi ^I\varphi ^I\right)^2$$
(1)
with a momentum cutoff $`\mathrm{\Lambda }_\mathrm{L}`$. To leading order in $`\frac{1}{N}`$, the renormalized squared mass $`m_r^2`$ and the renormalized self-coupling $`\lambda _r`$ are obtained as<sup>5</sup><sup>5</sup>5The terms of order $`\frac{m_r^2}{\mathrm{\Lambda }_\mathrm{L}^2}`$ and less are ignored.
$`m^2`$ $`=`$ $`m_r^2+{\displaystyle \frac{\lambda }{2(4\pi )^2}}\left(\mathrm{\Lambda }_\mathrm{L}^2+m_r^2\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }_\mathrm{L}^2}{m_r^2}}\right)`$ (2)
$`{\displaystyle \frac{1}{\lambda _r}}`$ $`=`$ $`{\displaystyle \frac{1}{\lambda }}+{\displaystyle \frac{1}{2(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }_\mathrm{L}^2}{\mu ^2}}{\displaystyle \frac{1}{2(4\pi )^2}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }_0^2}{\mu ^2}}`$ (3)
where $`\mu `$ is an arbitrary renormalization scale. The correlation functions of $`\varphi ^I`$ are made UV finite in terms of $`m_r^2`$ and $`\lambda _r`$. This is the standard renormalization. In Eq. (3) we have introduced the Landau scale $`\mathrm{\Lambda }_0`$ at which the bare coupling $`\lambda `$ diverges for a given $`\lambda _r`$. The Landau scale gives the largest energy scale beyond which the theory is not defined. The linear sigma model is thus characterized uniquely by two parameters: $`m_r^2`$ and $`\mathrm{\Lambda }_0`$.
Let us now take a look at the non-linear sigma model. It is defined by the action density
$$=N\left[\frac{v^2}{2}_\mu \mathrm{\Phi }^I_\mu \mathrm{\Phi }^I+\frac{a}{8}\left(_\mu \mathrm{\Phi }^I_\mu \mathrm{\Phi }^I\right)^2\right]$$
(4)
with a momentum cutoff $`\mathrm{\Lambda }_{\mathrm{NL}}`$. Due to the non-linear constraint
$$\mathrm{\Phi }^I\mathrm{\Phi }^I=1,$$
(5)
the theory is not renormalizable by the usual power counting.
Some comments on the four-derivative term in Eq. (4) are in order. At first sight it seems totally irrelevant; if we expand the action density naïvely in terms of the unconstrained fields $`\pi ^i=\sqrt{N}v\mathrm{\Phi }^i`$ ($`i=1,\mathrm{},N1`$), the four-derivative term gives rise to an interaction term suppressed by $`\frac{1}{v^2}`$ which is of order $`\frac{1}{\mathrm{\Lambda }_{\mathrm{NL}}^2}`$. The four-derivative term is indeed irrelevant but not for that reason; it can give marginal contributions, i.e., contributions not suppressed by inverse powers of the momentum cutoff, at low energies. Here, the role of $`a`$ is merely to rescale the momentum cutoff, and therefore it is physically irrelevant.
It is the existence of a critical value $`v_c^2`$ which is the key to the renormalizability of the non-linear model and its equivalence to the linear model. At the critical point $`v^2=v_c^2`$ the theory becomes a theory of $`N`$ free massless scalars. For $`v^2<v_c^2`$ the larger fluctuations of the fields are encouraged, and the $`O(N)`$ symmetry is fully restored. For $`v^2>v_c^2`$, however, the fluctuations are discouraged, and the $`O(N)`$ is spontaneously broken to $`O(N1)`$. By fine-tuning the parameter $`v^2`$ near $`v_c^2`$, the non-linear sigma model gives the same physics as the linear sigma model; all the differences are suppressed by positive powers of $`\frac{\mu ^2}{\mathrm{\Lambda }_{\mathrm{NL}}^2}`$ where $`\mu `$ is an arbitrary but finite renormalization scale.
The large $`N`$ calculations for the non-linear model is well known. For simplicity, we restrict ourselves to the symmetric phase. To leading order in $`\frac{1}{N}`$, the critical value $`v_c^2`$ is given by
$$\frac{v_c^2}{z}\frac{\mathrm{\Lambda }_{\mathrm{NL}}^2}{(4\pi )^2},$$
(6)
where
$$z1(4\pi )^2\frac{a}{4}.$$
(7)
We take $`a<\frac{4}{(4\pi )^2}`$ so that $`z>0`$.
By straightforward calculations, we can verify that the non-linear model is equivalent to the linear sigma model with renormalized squared mass $`m_r^2`$ and self-coupling $`\lambda _r`$, if we choose $`v^2`$ by
$$\frac{v^2v_c^2}{z}=\frac{m_r^2}{(4\pi )^2}\mathrm{ln}\frac{\mathrm{\Lambda }_{\mathrm{NL}}^2}{m_r^2}<0$$
(8)
and choose the cutoff $`\mathrm{\Lambda }_{\mathrm{NL}}`$ by
$$\mathrm{ln}\frac{\mathrm{\Lambda }_{\mathrm{NL}}^2}{\mathrm{\Lambda }_0^2}=2(4\pi )^2\frac{a}{4}$$
(9)
where $`\mathrm{\Lambda }_0`$ is related to $`\lambda _r`$ by Eq. (3). We see that the constant $`a`$ merely changes the ratio of $`\mathrm{\Lambda }_{\mathrm{NL}}`$ to $`\mathrm{\Lambda }_0`$ by a finite amount.
A similar analysis can be given to the Nambu-Jona-Lasinio model. The Nambu-Jona-Lasinio model with a non-renormalizable Fermi interaction is equivalent to a perturbatively renormalizable Yukawa theory if we ignore contributions suppressed by negative powers of the momentum cutoff.
## 3 QED with electrons
In this and next sections we analyze models with non-renormalizable current-current interactions. We first consider a purely fermionic theory defined by the following action density<sup>6</sup><sup>6</sup>6Our convention for euclidean fermionic fields might differ somewhat from the standard convention. Replace $`\overline{\psi }`$ by $`i\overline{\psi }`$ to get a more familiar kinetic term $`\overline{\psi }(\text{/}M)\psi `$. The hermitian gamma matrices $`\gamma _\mu `$ satisfy the Clifford algebra $`\{\gamma _\mu ,\gamma _\nu \}=2\delta _{\mu \nu }`$ as usual.
$$=\overline{\psi }^I\left(\frac{1}{i}\text{/}+iM\right)\psi ^I\frac{1}{2Nv^2}J_\mu J_\mu ,$$
(10)
where $`I`$ runs from $`1`$ to $`N`$, and the current $`J_\mu `$ is defined by
$$J_\mu \overline{\psi }^I\gamma _\mu \psi ^I.$$
(11)
To define a theory we introduce a momentum cutoff $`\mathrm{\Lambda }`$. It is essential to use a momentum cutoff as opposed to the dimensional regularization. We are interested in the dependence of the theory on the UV cutoff, and the dimensional regularization is not suitable for this purpose, since it automatically gives the limit of an infinite cutoff.
It turns out that the theory defined by the action density (10) and the current (11) is missing one marginal parameter.<sup>7</sup><sup>7</sup>7Similarly, in ref. a dimension eight field was introduced to the action density to account for a missing marginal parameter in the naïve NJL model. Instead of the current given by Eq. (11), we consider a more general
$$J_\mu \overline{\psi }^I\gamma _\mu \psi ^I+\frac{h}{\mathrm{\Lambda }^2}\overline{\psi }^I\stackrel{}{\text{/}}\gamma _\mu \text{/}\psi ^I$$
(12)
The momentum cutoff does not respect gauge invariance, and we need to adjust the coefficient $`h`$ for the Ward identities.<sup>8</sup><sup>8</sup>8See Appendix A for a summary of the Ward identities for QED.
It is important to observe that the action density $``$ is invariant under the charge conjugation $`𝒞`$ defined by
$$\psi C\overline{\psi }^T,\overline{\psi }\psi ^TC^1$$
(13)
where the four-by-four matrix $`C`$ satisfies
$$C^1\gamma _\mu C=\gamma _\mu ^T.$$
(14)
The current $`J_\mu `$ is odd under $`𝒞`$.
It is not necessary but helps our calculations to introduce a vector auxiliary field $`A_\mu `$. We then rewrite the action density as
$``$ $`=`$ $`\overline{\psi }^I\left({\displaystyle \frac{1}{i}}\text{/}+iM\right)\psi ^I{\displaystyle \frac{1}{2Nv^2}}J_\mu J_\mu +{\displaystyle \frac{1}{2}}\left(vA_\mu +{\displaystyle \frac{1}{\sqrt{N}v}}J_\mu \right)^2`$ (15)
$`=`$ $`\overline{\psi }^I\left({\displaystyle \frac{1}{i}}\text{/}+iM\right)\psi ^I+{\displaystyle \frac{1}{2}}v^2A_\mu ^2`$
$`+{\displaystyle \frac{1}{\sqrt{N}}}A_\mu \left(\overline{\psi }^I\gamma _\mu \psi ^I+{\displaystyle \frac{h}{\mathrm{\Lambda }^2}}\overline{\psi }^I\stackrel{}{\text{/}}\gamma _\mu \text{/}\psi ^I\right)`$
What do we expect at energy scales much below the cutoff $`\mathrm{\Lambda }`$? We ignore the effects suppressed by the inverse powers of the cutoff $`\mathrm{\Lambda }`$ but keep those only suppressed by negative powers of the logarithm of $`\mathrm{\Lambda }`$. If we fine-tune the mass parameter $`v`$ so that the mass scale of the theory remains UV finite, we expect that the above theory becomes equivalent to a theory which is renormalizable by power counting. The smallest renormalizable theory with fermions $`\psi ^I`$ and a real vector field $`A_\mu `$ is given by the following action density:
$$_{ren}=\frac{1}{4}F_{\mu \nu }^2+\frac{1}{2\xi _0}(_\mu A_\mu )^2+\frac{m_0^2}{2}A_\mu ^2+\frac{\lambda _0}{8N}(A_\mu ^2)^2+\overline{\psi }^I\left(\frac{1}{i}\text{/}+\frac{e_0}{\sqrt{N}}A\text{/}+iM\right)\psi ^I$$
(16)
where we have imposed the $`𝒞`$ invariance. At tree level the Ward identity demands that the parameter $`\lambda _0`$ vanish. But renormalizability alone allows an arbitrary $`\lambda _0`$.
To leading order in $`\frac{1}{N}`$ the fermion-photon vertex receives no radiative correction, and we only need to calculate the one-loop contributions to the two- and four-point proper vertices of the photon field $`A_\mu `$. The three-point vertex vanishes due to the $`𝒞`$ invariance. If the theories defined by (15) and (16) have the same two- and four-point vertices, the two theories are equivalent, since any higher point functions can be constructed out of the fermion-photon vertex and the photon two- and four-point vertices independently of the cutoff $`\mathrm{\Lambda }`$.
The one-loop calculations with a momentum cutoff $`\mathrm{\Lambda }`$ are straightforward, and we only write the results here. The inverse propagator is calculated from the top one-loop diagram in Fig. 2 as follows:
$`\mathrm{\Pi }_{\alpha \beta }(k^2)={\displaystyle \frac{1}{e^2}}[m_\gamma ^2\delta _{\alpha \beta }+{\displaystyle \frac{1}{\xi }}k_\alpha k_\beta `$ (17)
$`+(k^2\delta _{\alpha \beta }k_\alpha k_\beta )(18{\displaystyle \frac{e^2}{(4\pi )^2}}{\displaystyle _0^1}dxx(1x)\mathrm{ln}{\displaystyle \frac{M^2+x(1x)k^2}{\mu ^2}})]`$
where we define
$`{\displaystyle \frac{(4\pi )^2}{e^2}}`$ $``$ $`{\displaystyle \frac{4}{3}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}1+{\displaystyle \frac{4}{3}}h{\displaystyle \frac{1}{9}}h^2`$ (18)
$`{\displaystyle \frac{(4\pi )^2}{e^2}}m_\gamma ^2`$ $``$ $`(4\pi )^2v^2+\left(2+4h{\displaystyle \frac{2}{3}}h^2\right)\mathrm{\Lambda }^2+(212h)M^2`$ (19)
$`{\displaystyle \frac{(4\pi )^2}{e^2}}{\displaystyle \frac{1}{\xi }}`$ $``$ $`{\displaystyle \frac{1}{3}}(1+4hh^2)`$ (20)
Eq. (19) implies that we must fine-tune the squared mass parameter $`v^2`$ so that the squared mass $`m_\gamma ^2`$ of the photon is finite and positive. This is fine-tuning as opposed to tuning, since $`\frac{v^2}{\mathrm{\Lambda }^2}`$ must be tuned to an order $`1`$ quantity to the accuracy of $`\frac{M^2}{\mathrm{\Lambda }^2}`$. With this fine-tuning, the above photon two-point function is identical to the one in the massive QED with the running gauge coupling constant $`e`$, photon mass $`m_\gamma `$, and gauge fixing parameter $`\xi `$. The mass parameter $`\mu `$ is an arbitrary renormalization scale.
The four-point function at zero external momenta is obtained from the bottom one-loop graphs in Fig. 2, and it depends on $`h`$:
$$\mathrm{\Pi }_{\alpha \beta \gamma \delta }=\frac{1}{N}\frac{1}{(4\pi )^2}\left(\delta _{\alpha \beta }\delta _{\gamma \delta }+\delta _{\alpha \gamma }\delta _{\beta \delta }+\delta _{\alpha \delta }\delta _{\beta \gamma }\right)\left(\frac{4}{3}16h+24h^2\frac{16}{3}h^3\right).$$
(21)
For the theory to be the massive QED, the Ward identity for the four-point function must be satisfied. Hence, we must choose the parameter $`h`$ such that the above four-point function vanishes:
$$\frac{4}{3}16h+24h^2\frac{16}{3}h^3=0.$$
(22)
This equation has three real roots. We can choose, for example, $`h0.097`$. We note that the choice of $`h`$ is a tuning, but not a fine-tuning; it must be tuned relative to order $`1`$, but not $`\frac{\mu ^2}{\mathrm{\Lambda }^2}`$.
The result corresponding to $`m_\gamma =0`$ can be obtained from the induced QED which is defined by the action density (15) with a specific choice $`v=h=0`$.<sup>9</sup><sup>9</sup>9In addition we must use a cutoff which respects the gauge invariance. Otherwise the electron loops will not be gauge invariant. We have chosen to study the more general (15) because our interest is to verify the equivalence of the non-renormalizable theory to the renormalizable theory which is defined by Eq. (16) with arbitrary parameters $`m_0^2,\lambda _0`$.
We have chosen the interaction to have the current-current form, but it was not necessary. Instead of modifying the current by a dimension $`5`$ term proportional to $`h`$, we could have introduced a counterterm
$$\mathrm{const}\times \left(\frac{1}{N\mathrm{\Lambda }^4}\left(\overline{\psi }^I\gamma _\mu \psi ^I\right)^2\right)^2$$
(23)
in the action density. This is how we introduce missing marginal parameters for scalar QED in the next section.
Before closing this section, we make a remark on the sign of the interaction term in the action density (10). Since the current $`J_\mu `$ defined by Eq. (12) is a real field, the current-current interaction term in (10) is negative definite. We believe it does not invalidate the theory. Our reasoning goes as follows. First we note that the theory defined by (10) is equivalent to the massive QED, which is a stable theory, modulo irrelevant differences of order $`\frac{\mu ^2}{\mathrm{\Lambda }^2}`$. Therefore, if the theory is unstable, the effects of the instability must be suppressed by positive powers of $`\frac{\mu ^2}{\mathrm{\Lambda }^2}`$. This suggests that a potential instability can only arise from the large fluctuations of the fields, for example $`A_\mu `$ of order $`\mathrm{\Lambda }`$. If this is the case, stability will be assured by redefining the theory by the action density (15) where the auxiliary field $`A_\mu `$ is restricted within a finite range $`|A_\mu |<\mathrm{\Lambda }`$. The effects of this modification are suppressed by positive powers of $`\frac{\mu ^2}{\mathrm{\Lambda }^2}`$.
## 4 Scalar QED
The fermionic theory may be a little too simple. It resembles the induced QED too much. If we had used a regularization which allows shifts of momentum such as the dimensional regularization, the vacuum polarization would have come out transverse, and the photon four-point function would have vanished at zero external momenta. The Ward identities are then satisfied automatically. Let us introduce a more non-trivial example of a purely bosonic theory in this section.
The theory is defined by the following action density:
$``$ $`=`$ $`_\mu \varphi ^I_\mu \varphi ^I+m^2\varphi ^I\varphi ^I+{\displaystyle \frac{\lambda }{4N}}\left(\varphi ^I\varphi ^I\right)^2`$ (24)
$`{\displaystyle \frac{v^2}{2}}\left({\displaystyle \frac{i}{\sqrt{N}v^2}}\varphi ^I\stackrel{}{_\mu }\varphi ^I\right)^2+\mathrm{\Delta },`$
where the asterisk $``$ denotes complex conjugation, and the counterterm $`\mathrm{\Delta }`$ is defined by
$`\mathrm{\Delta }`$ $``$ $`{\displaystyle \frac{a}{N}}(\varphi ^I\varphi ^I){\displaystyle \frac{1}{2}}\left({\displaystyle \frac{i}{\sqrt{N}v^2}}\varphi ^I\stackrel{}{_\mu }\varphi ^I\right)^2`$ (25)
$`+{\displaystyle \frac{b}{N(4\pi )^2}}{\displaystyle \frac{1}{8}}\left\{\left({\displaystyle \frac{i}{\sqrt{N}v^2}}\varphi ^I\stackrel{}{_\mu }\varphi ^I\right)^2\right\}^2.`$
We have introduced enough number of parameters so that the theory is equivalent to the theory defined by the following renormalizable action density:
$`_{ren}`$ $`=`$ $`_\mu \varphi ^I_\mu \varphi ^I+m_0^2\varphi ^I\varphi ^I+{\displaystyle \frac{\lambda _0}{4N}}(\varphi ^I\varphi ^I)^2`$ (26)
$`+{\displaystyle \frac{1}{4}}F_{\mu \nu }^2+{\displaystyle \frac{1}{2\xi _0}}(_\mu A_\mu )^2+{\displaystyle \frac{m_{\gamma ,0}^2}{2}}A_\mu ^2`$
$`+{\displaystyle \frac{e_0}{\sqrt{N}}}A_\mu i\varphi ^I\stackrel{}{_\mu }\varphi ^I+{\displaystyle \frac{\gamma }{2N}}\varphi ^I\varphi ^IA_\mu ^2+{\displaystyle \frac{\delta }{8N}}{\displaystyle \frac{1}{(4\pi )^2}}(A_\mu ^2)^2`$
Just as the Ward identities can determine the constants $`\gamma `$ and $`\delta `$ uniquely, we will be able to fix the coefficients $`a,b`$ of the counterterms by imposing the Ward identities. We should note that unlike the fermionic theory discussed in the previous section, the interaction of this theory is not solely current-current type due to the counterterms.
To facilitate the $`\frac{1}{N}`$ expansions, we introduce scalar and vector auxiliary fields $`\alpha ,A_\mu `$ and rewrite the action density as follows:
$``$ $`=`$ $`_\mu \varphi ^I_\mu \varphi ^I+m^2\varphi ^I\varphi ^I+{\displaystyle \frac{\lambda }{4N}}\left(\varphi ^I\varphi ^I\right)^2{\displaystyle \frac{v^2}{2}}\left({\displaystyle \frac{i}{\sqrt{N}v^2}}\varphi ^I\stackrel{}{_\mu }\varphi ^I\right)^2`$ (27)
$`+{\displaystyle \frac{1}{2}}\left(\alpha +i\sqrt{{\displaystyle \frac{\lambda }{N}}}\varphi ^I\varphi ^I\right)^2+{\displaystyle \frac{1}{2}}\left(vA_\mu +{\displaystyle \frac{i}{\sqrt{N}v}}\varphi ^I\stackrel{}{_\mu }\varphi ^I\right)^2+\mathrm{\Delta }`$
$`=`$ $`_\mu \varphi ^I_\mu \varphi ^I+{\displaystyle \frac{1}{2}}\alpha ^2+i\sqrt{{\displaystyle \frac{\lambda }{N}}}\alpha \varphi ^I\varphi ^I`$
$`+{\displaystyle \frac{v^2}{2}}A_\mu ^2+{\displaystyle \frac{1}{\sqrt{N}}}A_\mu i\varphi ^I\stackrel{}{_\mu }\varphi ^I+\mathrm{\Delta }.`$
Unlike the fermionic theory of the previous section, this theory does not reduce to an induced QED for $`v=0`$.
To leading order in $`\frac{1}{N}`$, we must renormalize the scalar mass and self-coupling as
$`\mathrm{\Delta }m^2`$ $``$ $`m_r^2m^2={\displaystyle \frac{\lambda }{(4\pi )^2}}\left(\mathrm{\Lambda }^2m_r^2\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_r^2}}\right)`$ (28)
$`{\displaystyle \frac{(4\pi )^2}{\lambda _r}}`$ $`=`$ $`{\displaystyle \frac{(4\pi )^2}{\lambda }}+\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}1`$ (29)
and shift the auxiliary field $`\alpha `$ by
$$\alpha =i\sqrt{\frac{N}{\lambda }}\mathrm{\Delta }m^2+\sqrt{\frac{\lambda _r}{\lambda }}\delta \alpha ,$$
(30)
where the shifted field $`\delta \alpha `$ has a vanishing expectation value $`\delta \alpha =0`$. To leading order in $`\frac{1}{N}`$, the full propagator of the fluctuation $`\delta \alpha `$ is given by
$$\stackrel{~}{\delta \alpha }(k)\delta \alpha =1/\left(1+\frac{\lambda _r}{(4\pi )^2}_0^1𝑑x\mathrm{ln}\frac{\mu ^2}{m_r^2+x(1x)k^2}\right).$$
(31)
Before calculating the vertex functions involving the vector field $`A_\mu `$, we note the implication of the equation of motion for $`A_\mu `$. The action density is quadratic with respect to $`A_\mu `$, and the equation of motion gives
$$A_\mu =B_\mu \frac{i}{\sqrt{N}v^2}\varphi ^I\stackrel{}{_\mu }\varphi ^I$$
(32)
This implies that the composite field $`B_\mu `$ is an interpolating field of the photon. The calculation of the proper vertex of $`A_\mu `$ and $`B_\nu `$ indeed gives
$$\stackrel{~}{A_\mu }(k)B_\nu =\delta _{\mu \nu }\frac{1}{v^2}_p\frac{p^2}{(p^2+m_r^2)^2}\delta _{\mu \nu }\frac{1}{v^2}\frac{\mathrm{\Lambda }^2}{(4\pi )^2}$$
(33)
where we have ignored the terms of order $`\frac{m_r^2}{\mathrm{\Lambda }^2}`$. (Fig. 3) This is $`\delta _{\mu \nu }`$ if we choose
$$v^2\frac{\mathrm{\Lambda }^2}{(4\pi )^2}.$$
(34)
We will see that Eq. (34) is required by the fine-tuning of the mass parameter. Hence, the equation of motion (32) implies that the counterterm in the action density is equivalent to
$$\mathrm{\Delta }=\frac{a}{2N}\left(\varphi ^I\varphi ^I\right)A_\mu ^2+\frac{b}{8N}(A_\mu ^2)^2$$
(35)
to leading order in $`\frac{1}{N}`$. Therefore, $`\mathrm{\Delta }`$ gives rise to the vertices in Fig. 4.
Let us proceed with the two-point vertex of the photon. To leading order in $`\frac{1}{N}`$, we obtain
$`\mathrm{\Pi }_{\alpha \beta }={\displaystyle \frac{1}{e^2}}[m_\gamma ^2\delta _{\alpha \beta }+{\displaystyle \frac{1}{\xi }}k_\alpha k_\beta `$ (36)
$`+(k^2\delta _{\alpha \beta }k_\alpha k_\beta )(1{\displaystyle \frac{e^2}{(4\pi )^2}}{\displaystyle _0^1}dx(12x)^2\mathrm{ln}{\displaystyle \frac{m_r^2+x(1x)k^2}{\mu ^2}})],`$
where the renormalized parameters are defined by
$`{\displaystyle \frac{(4\pi )^2}{e^2}}`$ $``$ $`{\displaystyle \frac{1}{3}}\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{\mu ^2}}{\displaystyle \frac{1}{2}}`$ (37)
$`m_\gamma ^2`$ $``$ $`e^2\left[v^2{\displaystyle \frac{1}{(4\pi )^2}}\left(\mathrm{\Lambda }^2+m_r^22m_r^2\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m_r^2}}\right)\right]`$ (38)
$`{\displaystyle \frac{1}{\xi }}`$ $``$ $`{\displaystyle \frac{e^2}{(4\pi )^2}}{\displaystyle \frac{1}{6}}.`$ (39)
Here, $`\mu `$ is an arbitrary renormalization scale as usual. Eq. (38) implies that we need to fine-tune the squared mass parameter $`v^2`$ such that $`0<m_\gamma ^2\mathrm{\Lambda }^2`$.
Now that we have identified the photon mass $`m_\gamma `$ and the gauge coupling $`e`$, we wish to proceed with verifying the equivalence of our theory with the scalar QED. This will be done by checking three Ward identities. (See Appendix A for a summary of the Ward identities.) The first is the Ward identity for the three-point vertex of the scalar and photon. To leading order in $`\frac{1}{N}`$, the vertex receives no radiative correction (Fig. 5), and the Ward identity is automatically satisfied.
Next we examine the scalar-scalar-photon-photon vertex. To leading order in $`\frac{1}{N}`$, it is given by the four diagrams in Fig. 6, where we denote the propagator of $`\delta \alpha `$ by a broken line.
The third and fourth terms involve the counterterm proportional to $`a`$. The Ward identity demands that at zero external momenta this be given by
$$\mathrm{\Pi }_{\mu \nu }^{\varphi \varphi ^{}}|_{\mathrm{zero}\mathrm{momenta}}=\frac{2}{N}\delta _{\mu \nu }.$$
(40)
To leading order in $`\frac{1}{N}`$ we obtain
$$\mathrm{\Pi }_{\mu \nu }^{\varphi \varphi ^{}}|_{\mathrm{zero}\mathrm{momenta}}=\frac{2}{N}\delta _{\mu \nu }\frac{\mathrm{ln}\frac{\mathrm{\Lambda }^2}{m_r^2}1+\frac{a}{2}\frac{(4\pi )^2}{\lambda }\frac{1}{2}}{\mathrm{ln}\frac{\mathrm{\Lambda }^2}{m_r^2}1+\frac{(4\pi )^2}{\lambda }}.$$
(41)
Therefore, we must choose $`a`$ as
$$a=\frac{\lambda }{(4\pi )^2}+2.$$
(42)
This is an ordinary tuning, since $`\lambda `$ is a quantity of order $`1`$.
Finally we examine the four-photon vertex. Many graphs contribute to leading order in $`\frac{1}{N}`$ as in Fig. 7.
The last term in Fig. 7 is the counterterm proportional to $`b`$. The Ward identity demands that the four-photon vertex vanish for zero external momenta. The calculation is straightforward but somewhat involved. We can use Eq. (40) to simplify the calculation. The final result is given by
$$\mathrm{\Pi }_{\alpha \beta \gamma \delta }|_{\mathrm{zero}\mathrm{momenta}}=\frac{1}{N}\frac{1}{(4\pi )^2}\left(\delta _{\alpha \beta }\delta _{\gamma \delta }+\delta _{\alpha \gamma }\delta _{\beta \delta }+\delta _{\alpha \delta }\delta _{\beta \gamma }\right)\left(\frac{4}{3}+ab\right).$$
(43)
Therefore, the Ward identity demands
$$b=\frac{\lambda }{(4\pi )^2}+\frac{2}{3}$$
(44)
where we have used Eq. (42). Thus, with the choice of the coefficients $`a,b`$ given by Eqs. (42,44), the theory defined by the action density (24) is equivalent to the scalar QED.
Before we close this section, we make a brief comment on the relation of the above model to the $`\mathrm{CP}^{N1}`$ model. (See also Chapter 5 of ref. .) The $`\mathrm{CP}^{N1}`$ model can be obtained formally from the action density (24) by imposing a non-linear constraint
$$\varphi ^I\varphi ^I=\frac{Nv^2}{2}.$$
(45)
Since the $`O(2N)`$ non-linear sigma model is equivalent to the $`O(2N)`$ linear sigma model, the $`\mathrm{CP}^{N1}`$ model is equivalent to the model discussed in this section, and therefore to the scalar QED.
## 5 Conclusion
We have seen that once the mass parameters are fine-tuned, theories which are non-renormalizable by the usual power counting rule reduce to renormalizable theories at energies below the cutoff $`\mathrm{\Lambda }`$ as long as we ignore quantities inversely proportional to $`\mathrm{\Lambda }`$. The renormalized parameters at a finite energy scale $`\mu `$ appear as the marginally irrelevant dependence on the cutoff of the order of $`\frac{1}{\mathrm{ln}\frac{\mathrm{\Lambda }}{\mu }}`$. We have studied two matter-only models without elementary gauge fields whose current-current contact interactions at the cutoff scale give rise to interactions mediated by the abelian gauge fields at low energies. The gauge field was dynamically generated to complete renormalizability of the theory just as the sigma field is dynamically generated in the $`O(N)`$ non-linear sigma model.
In this paper the massive photon was obtained not as a consequence of the Higgs mechanism: it appeared as part of the gauge fixing terms. The mass is simply allowed by the Ward identities in the case of abelian gauge theories. In the next paper we will study a matter-only model which exhibits the Higgs mechanism.
It should be interesting to extend our analysis to non-abelian gauge theories. We expect that the covariant gauge fixing term will arise naturally just as for QED, but in the case of non-abelian gauge theories the covariant gauge requires the Faddeev-Popov (FP) ghosts. Hence, for a matter-only model to become a non-abelian gauge theory, the FP ghosts must be generated dynamically. It will be extremely interesting if this is the case.
In ref. an induced QCD was studied in the $`1/N`$ expansions in the hope of uncovering the non-perturbative dynamics of QCD. We also hope that the reformulation of gauge theories as matter-only theories will find practical applications in understanding, for example, the physics of QED at energies very high compared to the electron mass but still much lower than the cutoff scale.
This work was partially supported by the Grant-In-Aid for Scientific Research (No. 11640279) from the Ministry of Education, Science, and Culture, Japan.
## Appendix A Ward identities
In determining the coefficients of the counterterms, we have imposed Ward identities. We remind the reader of the Ward identities for both QED with electrons and the scalar QED. For QED with electrons, we have three Ward identities to satisfy:
$`k_\mu \mathrm{\Pi }_{\mu \nu }(k)`$ $`=`$ $`{\displaystyle \frac{k_\nu }{e^2}}\left(m_\gamma ^2+{\displaystyle \frac{1}{\xi }}k^2\right)`$ (46)
$`ik_\mu \mathrm{\Pi }_\mu ^{\psi \overline{\psi }}(p,k)`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{N}}}\left(\mathrm{\Pi }^{\psi \overline{\psi }}(p)\mathrm{\Pi }^{\psi \overline{\psi }}(p+k)\right)`$ (47)
$`k_\alpha \mathrm{\Pi }_{\alpha \beta \gamma \delta }`$ $`=`$ $`0,`$ (48)
where $`\mathrm{\Pi }_\mu ^{\psi \overline{\psi }}(p,k)`$ is the electron-photon interaction vertex, and $`\mathrm{\Pi }^{\psi \overline{\psi }}(p)`$ is the inverse electron propagator with momentum $`p`$.
The Ward identities for the scalar QED are similarly given by
$`k_\mu \mathrm{\Pi }_{\mu \nu }(k)`$ $`=`$ $`{\displaystyle \frac{k_\nu }{e^2}}\left(m_\gamma ^2+{\displaystyle \frac{1}{\xi }}k^2\right)`$ (49)
$`ik_\mu \mathrm{\Pi }_\mu ^{\varphi \varphi ^{}}(p,k)`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{N}}}\left(\mathrm{\Pi }^{\varphi \varphi ^{}}(p)\mathrm{\Pi }^{\varphi \varphi ^{}}(p+k)\right)`$ (50)
$`ik_\mu \mathrm{\Pi }_{\mu \nu }^{\varphi \varphi ^{}}(p,k,l)`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{N}}}\left(\mathrm{\Pi }_\nu ^{\varphi \varphi ^{}}(p,l)+\mathrm{\Pi }_\nu ^{\varphi \varphi ^{}}(p+k,l)\right)`$ (51)
$`k_\alpha \mathrm{\Pi }_{\alpha \beta \gamma \delta }`$ $`=`$ $`0`$ (52)
using a similar notation as for QED.
## Appendix B Integrals with a momentum cutoff
In calculating the one-loop integrals with a momentum cutoff, we have used the following formulas where contributions of order $`\frac{m^2}{\mathrm{\Lambda }^2}`$ or less are ignored.
$`{\displaystyle _{p<\mathrm{\Lambda }}}{\displaystyle \frac{1}{p^2+m^2}}`$ $``$ $`{\displaystyle _{p^2<\mathrm{\Lambda }^2}}{\displaystyle \frac{d^4p}{(2\pi )^4}}{\displaystyle \frac{1}{p^2+m^2}}`$ (53)
$`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left(\mathrm{\Lambda }^2m^2\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m^2}}\right)`$
$`{\displaystyle _{p<\mathrm{\Lambda }}}{\displaystyle \frac{p_\mu p_\nu }{(p^2+m^2)^2}}`$ $`=`$ $`{\displaystyle \frac{\delta _{\mu \nu }}{4}}{\displaystyle _p}{\displaystyle \frac{p^2}{(p^2+m^2)^2}}`$ (54)
$`=`$ $`{\displaystyle \frac{\delta _{\mu \nu }}{4}}{\displaystyle \frac{1}{(4\pi )^2}}\left(\mathrm{\Lambda }^22m^2\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m^2}}+m^2\right)`$
$`{\displaystyle _{p<\mathrm{\Lambda }}}{\displaystyle \frac{1}{(p^2+m^2)^2}}`$ $`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m^2}}1\right)`$ (55)
$`{\displaystyle _{p<\mathrm{\Lambda }}}{\displaystyle \frac{p_\mu p_\nu }{(p^2+m^2)^3}}`$ $`=`$ $`{\displaystyle \frac{\delta _{\mu \nu }}{4}}{\displaystyle \frac{1}{(4\pi )^2}}\left(\mathrm{ln}{\displaystyle \frac{\mathrm{\Lambda }^2}{m^2}}{\displaystyle \frac{3}{2}}\right)`$ (56) |
warning/0002/cond-mat0002227.html | ar5iv | text | # Density of States of a d-wave Superconductor in the Presence of Strong Impurity Scatterers: a Non Perturbative Result
## I The model
The generic Hamiltonian for a d-wave superconductor can be written
$$H_0=\underset{𝐤}{}\varphi _𝐤^{}\left[\epsilon _𝐤\sigma _3+\mathrm{\Delta }_𝐤\sigma _1\right]\varphi _𝐤.$$
(1)
It describes BCS quasiparticles with the kinetic energy $`\epsilon _𝐤=W\left(\mathrm{cos}k_x+\mathrm{cos}k_y\right)\mu `$ ($`\mu `$ is the chemical potential) in the presence of the spin singlet superconducting order parameter $`\mathrm{\Delta }_𝐤=\mathrm{\Delta }_0\left(\mathrm{cos}k_x\mathrm{cos}k_y\right)`$. Distances are measured in units of the lattice constant. The $`\sigma _i`$ are the Pauli matrices in the particle-hole space
$`\begin{array}{cc}\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),& \sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)\end{array}.`$The spinor $`\varphi _𝐤^{}=(c_{𝐤,}^{},c_{𝐤,})`$ creates a particle and a hole with momenta $`𝐤`$ and $`𝐤`$, respectively. We shall present our results in terms of a $`d_{x^2y^2}`$ state even though our conclusions apply more generally to any state where $`\mathrm{\Delta }_𝐤`$ vanishes linearly along a direction parallel to the Fermi surface.
Instead of using the Nambu formalism, we work with the diagonalized version of (1) in order to access directly the properties of quasiparticles. The Bogoliubov transformation that diagonalizes $`H_0`$ is given by
$`c_𝐤`$ $`=`$ $`u_𝐤\alpha _𝐤v_𝐤\beta _𝐤,\alpha _𝐤=u_𝐤c_𝐤+v_𝐤c_𝐤^{},`$ (2)
$`c_𝐤^{}`$ $`=`$ $`v_𝐤\alpha _𝐤+u_𝐤\beta _𝐤,\beta _𝐤=v_𝐤c_𝐤+u_𝐤c_𝐤^{},`$ (3)
where $`\alpha _𝐤`$ and $`\beta _𝐤`$ create a particle and a hole with momentum $`k`$ and the coefficients $`u_𝐤`$ and $`v_𝐤`$ satisfy
$$\begin{array}{c}u_𝐤^2=1/2\left(1+\frac{\epsilon _𝐤}{\omega _𝐤}\right),\hfill \\ v_𝐤^2=1/2\left(1\frac{\epsilon _𝐤}{\omega _𝐤}\right),\hfill \\ u_𝐤v_𝐤=\frac{\mathrm{\Delta }_𝐤}{2\omega _𝐤},\hfill \end{array}$$
(4)
with $`\omega _𝐤=\sqrt{\epsilon _𝐤^2+\mathrm{\Delta }_𝐤^2}`$. Given the short-hand notation
$$\begin{array}{c}\psi _{𝐤,0}^{}\alpha _𝐤^{},\hfill \\ \psi _{𝐤,1}^{}\beta _𝐤^{},\hfill \end{array}$$
(5)
the BCS Hamiltonian can now be rewritten
$`H_0={\displaystyle \underset{𝐤}{}}{\displaystyle \underset{\nu =0,1}{}}\omega _𝐤(1)^\nu \psi _{𝐤,\nu }^{}\psi _{𝐤,\nu }.`$The disorder is introduced through $`N`$ repulsive scalar potentials $`V_0`$ located at random positions in the lattice:
$$H_I=V_0\underset{i=1}{\overset{N}{}}\underset{\sigma =,}{}c_{i\sigma }^{}c_{i\sigma }.$$
(6)
The full BCS Hamiltonian $`H=H_0+H_I`$ describes a dirty d-wave superconductor.
## II The T-matrix equation
### A The Hamiltonian
With the help of a Fourier transformation to the reciprocal lattice, the impurity potential becomes
$`H_I`$ $`={\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}{\displaystyle \underset{𝐤,𝐤^{}}{}}e^{i(𝐤𝐤^{})𝐑_i}\left(c_𝐤^{}c_𝐤^{}+c_𝐤^{}c_𝐤^{}\right),`$ (7)
where $`𝒱`$ is the volume of the system. Rewriting the impurity term in terms of quasiparticles gives
$`c_𝐤^{}c_𝐤^{}=u_𝐤u_𝐤^{}\alpha _𝐤^{}\alpha _𝐤^{}v_𝐤u_𝐤^{}\beta _𝐤^{}\alpha _𝐤^{}u_𝐤v_𝐤^{}\alpha _𝐤^{}\beta _𝐤^{}+v_𝐤v_𝐤^{}\beta _𝐤^{}\beta _𝐤^{}.`$Thus
$$c_𝐤^{}c_𝐤^{}=\underset{\nu ,\nu ^{}=0}{\overset{1}{}}(1)^\nu (1)^\nu ^{}t_{𝐤\nu }t_{𝐤^{}\nu ^{}}\psi _{𝐤\nu }^{}\psi _{𝐤^{}\nu ^{}},$$
(8)
where we have introduced the short-hand notation
$$\begin{array}{cc}\begin{array}{c}t_{𝐤,0}u_𝐤,\hfill \\ t_{𝐤,1}v_𝐤.\hfill \end{array}& \end{array}$$
(9)
Similarly,
$`c_𝐤^{}c_𝐤^{}=\left(v_𝐤v_𝐤^{}\alpha _𝐤^{}\alpha _𝐤^{}+v_𝐤u_𝐤^{}\beta _𝐤^{}\alpha _𝐤^{}+u_𝐤v_𝐤^{}\alpha _𝐤^{}\beta _𝐤^{}+u_𝐤u_𝐤^{}\beta _𝐤^{}\beta _𝐤^{}\right)+const.,`$and, neglecting the constant term, we get
$$c_𝐤^{}c_𝐤^{}=\underset{\nu ,\nu ^{}}{}t_{𝐤\nu +1}t_{𝐤^{}\nu ^{}+1}\psi _{𝐤\nu }^{}\psi _{𝐤^{}\nu ^{}}.$$
(10)
In summary, the random BCS Hamiltonian can be written
$`H`$ $`=`$ $`{\displaystyle \underset{𝐤,\nu }{}}\omega _𝐤(1)^\nu \psi _{𝐤,\nu }^{}\psi _{𝐤,\nu }`$ (11)
$`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}{\displaystyle \underset{𝐤,𝐤^{},\nu ,\nu ^{}}{}}e^{i(𝐤𝐤^{})𝐑_i}\left[(1)^\nu (1)^\nu ^{}t_{𝐤\nu }t_{𝐤^{}\nu ^{}}t_{𝐤\nu +1}t_{𝐤^{}\nu ^{}+1}\right]\psi _{𝐤\nu }^{}\psi _{𝐤^{}\nu ^{}}.`$ (12)
### B Equations of motion
As the impurities break translation invariance, the anomalous two-point function
$$G_{\mathrm{𝐤𝐪}}^{\nu \nu ^{}}(\tau )=T_\tau \left[\psi _{𝐤\nu }(\tau )\psi _{𝐪\nu ^{}}^{}(0)\right]$$
(13)
depends on two momenta. The equations of motion are
$$\underset{𝐪,m}{}_{\mathrm{𝐤𝐪}}^{\nu m}G_{\mathrm{𝐪𝐤}^{}}^{m\nu ^{}}=\delta _{\mathrm{𝐤𝐤}^{}}\delta _{\nu \nu ^{}},$$
(14)
where $`_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}`$ is
$`_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}`$ $`=`$ $`[_\tau +(1)^\nu \omega _𝐤]\delta _{\mathrm{𝐤𝐤}^{}}\delta _{\nu \nu ^{}}`$ (15)
$`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}e^{i(𝐤𝐤^{})𝐑_i}(1)^\nu (1)^\nu ^{}t_{𝐤\nu }t_{𝐤^{}\nu ^{}}`$ (16)
$``$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}e^{i(𝐤𝐤^{})𝐑_i}t_{𝐤\nu +1}t_{𝐤^{}\nu ^{}+1}.`$ (17)
We are dealing with a problem of non-interacting particles scattered by the static potential generated by $`N`$ impurities. As we shall see, the anomalous Green function in Eq. (13) can be solved by inverting a $`2N\times 2N`$ matrix. To this end, define the one-point functions
$$G_{𝐤\nu }^0(\tau )\frac{1}{_\tau +(1)^\nu \omega _𝐤},$$
(18)
$$\begin{array}{c}g_{𝐤\nu }^1(𝐑_i)(1)^\nu e^{i𝐤𝐑_i}t_{𝐤\nu }G_{𝐤\nu }^0,\hfill \\ g_{𝐤\nu }^2(𝐑_i)e^{i𝐤𝐑_i}t_{𝐤\nu +1}G_{𝐤\nu }^0,\hfill \end{array}$$
(19)
together with
$`A_{ij}`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}e^{i𝐤(𝐑_i𝐑_j)}(t_{𝐤\nu })^2G_{𝐤\nu }^0},`$ (20)
$`C_{ij}`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}e^{i𝐤(𝐑_i𝐑_j)}(t_{𝐤\nu +1})^2G_{𝐤\nu }^0},`$ (21)
$`B_{ij}`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}e^{i𝐤(𝐑_i𝐑_j)}(1)^nt_{𝐤\nu }t_{𝐤\nu +1}G_{𝐤\nu }^0}.`$ (22)
The integration over $`𝐤`$ is to be performed over the first Brillouin zone. Equation (15) is inserted into the equations of motion (14). We then solve a $`2N\times 2N`$ system of linear equations (see appendix A). This gives
$$G_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}=G_{\mathrm{𝐤𝐤}^{}}^0\frac{V_0}{𝒱}𝐍_{𝐤\nu }^T\widehat{M}^1𝐍_{𝐤^{}\nu ^{}},$$
(23)
with $`𝐍_{𝐤\nu }`$ a vector made of the $`2N`$ components
$$\begin{array}{ccc}𝐍_{𝐤\nu }\left(\begin{array}{c}𝐍_{𝐤\nu }^1\\ 𝐍_{𝐤\nu }^2\end{array}\right),\hfill & 𝐍_{𝐤\nu }^1\left(\begin{array}{c}g_{𝐤\nu }^1(𝐑_1)\\ \mathrm{}\\ g_{𝐤\nu }^1(𝐑_N)\end{array}\right),\hfill & 𝐍_{𝐤\nu }^2\left(\begin{array}{c}g_{𝐤\nu }^2(𝐑_1)\\ \mathrm{}\\ g_{𝐤\nu }^2(𝐑_N)\end{array}\right),\hfill \end{array}$$
(24)
and $`\widehat{M}`$ is a $`2N\times 2N`$ matrix defined by
$$\widehat{M}=\left[\begin{array}{cc}\widehat{I}+V_0\widehat{A}& V_0\widehat{B}\\ V_0\widehat{B}& \widehat{I}+V_0\widehat{C}\end{array}\right],$$
(25)
with $`\widehat{A}`$, $`\widehat{B}`$, $`\widehat{C}`$ are $`N\times N`$ matrices whose matrix elements are respectively $`A_{ij}`$, $`B_{ij}`$ and $`C_{ij}`$. $`\widehat{I}`$ is the identity matrix. With these definitions the T-matrix equation can be written
$$G_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}=G_{\mathrm{𝐤𝐤}^{}}^0+𝐍_{𝐤\nu }^T\widehat{T}𝐍_{𝐤^{}\nu ^{}},$$
(26)
with $`\widehat{T}=\frac{V_0}{𝒱}\widehat{M}^1`$.
## III A Sum Rule for the Density of States
The increment in the density of states induced by the impurities is
$$\delta \rho (\omega )=\frac{1}{\pi }Im\underset{𝐤\nu }{}\delta G_{\mathrm{𝐤𝐤}}^{\nu \nu }(\omega +i0^+),$$
(27)
where
$$\begin{array}{cc}& \delta G_{\mathrm{𝐤𝐤}}^{\nu \nu }\frac{V_0}{𝒱}𝐍_{𝐤\nu }^T\widehat{M}^1𝐍_{𝐤\nu }.\end{array}$$
(28)
Recall that the summation over $`𝐤`$ is restricted to the first Brillouin zone. We can rewrite
$`\delta G_{\mathrm{𝐤𝐤}}^{\nu \nu }`$ $`=`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i,j}{}}N_{𝐤\nu }^iM_{ij}^1N_{𝐤\nu }^j`$ (29)
$`=`$ $`{\displaystyle \underset{i,j}{}}M_{ij}^1\left({\displaystyle \frac{V_0}{𝒱}}N_{𝐤\nu }^jN_{𝐤\nu }^i\right).`$ (30)
If we go to Matsubara frequencies and define
$`G_{𝐤\nu }^0={\displaystyle \frac{1}{i\omega _n(1)^\nu \omega _𝐤}},`$we notice that
$`{\displaystyle \underset{𝐤\nu }{}}\left(V_0N_{𝐤\nu }^jN_{𝐤\nu }^i\right)={\displaystyle \frac{}{i\omega _n}}\widehat{M}_{ij},`$and then
$$\underset{𝐤\nu }{}\delta G_{\mathrm{𝐤𝐤}}^{\nu \nu }(i\omega _n)=\frac{1}{𝒱}Tr\left[\widehat{M}^1\frac{\widehat{M}}{i\omega _n}\right],$$
(31)
where the Trace is assumed to run over $`𝐤`$, $`\nu `$ as well as over the matrix indices. By convention we call
$$\delta G\underset{𝐤\nu }{}\delta G_{\mathrm{𝐤𝐤}}^{\nu \nu }.$$
(32)
The equation (31) can also be written ($`\mathrm{ln}Det=Tr\mathrm{ln}`$)
$$\delta G(i\omega _n)=\frac{1}{𝒱}\frac{\left(\mathrm{ln}Det\widehat{M}\right)}{i\omega _n}.$$
(33)
We used this expression in a previous paper in order to derive the additional density of states. In this paper we will prefer to use the following expression
$$\delta G(i\omega _n)=\frac{1}{𝒱}Tr\left[\widehat{M}^2\frac{\widehat{M}^2}{2i\omega _n}\right].$$
(34)
## IV The matrix elements of $`\widehat{M}`$
As seen previously the evaluation of the density of states relies on calculating the determinant of $`\widehat{M}`$. We begin by evaluating its matrix elements. In this section, we just quote the result; the explicit calculation being given in Appendix B. In order to perform this calculation we make the assumption that the energy $`\omega \mathrm{\Delta }_0`$. This will enable us to linearize the spectrum for small energies.
Second we assume that $`W=\mathrm{\Delta }_0`$ and the chemical potential $`\mu =0`$ in order to simplify the calculation. We will see later what happens when these two conditions are relaxed. Under these conditions there are four nodes in the Brillouin zone (cf. figure 1) located at $`(\pm \frac{\pi }{2},\pm \frac{\pi }{2})`$. With $`𝐤^{}(\frac{\pi }{2},\frac{\pi }{2})+𝐤`$ we get
$`\omega _𝐤^{}^2`$ $`=`$ $`W^2\left[\left(\mathrm{cos}k_x^{}+\mathrm{cos}k_y^{}\mu \right)^2+\left(\mathrm{cos}k_x^{}\mathrm{cos}k_y^{}\right)^2\right]`$ (35)
$``$ $`2W^2k^2.`$ (36)
Thus, $`\omega _𝐤=Dk`$ with $`D=\sqrt{2}W`$. after integrating over the four nodes in the Brillouin zone, we find
$`A_{ij}`$ $`=`$ $`𝒜^0(𝐑_{ij})+𝒜^1(𝐑_{ij}),`$ (37)
$`C_{ij}`$ $`=`$ $`𝒜^0(𝐑_{ij})𝒜^1(𝐑_{ij}),`$ (38)
$`B_{ij}`$ $`=`$ $`^1(𝐑_{ij}),`$ (39)
where
$`𝒜^0(𝐑)`$ $``$ $`i\omega _n{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{1}{(i\omega _n)^2\omega _𝐤^2}}e^{i𝐤𝐑}`$ (40)
$`=`$ $`^0(𝐑){\displaystyle \frac{i\omega _n}{2\pi D^2}}K_0\left(\left|{\displaystyle \frac{R\omega _n}{D}}\right|\right),`$ (41)
whereby
$$\begin{array}{cc}& ^0(𝐑)=2\mathrm{cos}\frac{\pi }{2}(R_x+Ry)+2\mathrm{cos}\frac{\pi }{2}(R_xR_y),\end{array}$$
(42)
and $`K_0`$ is the Bessel function of rank zero. Note that $`|\omega _n|=\sqrt{(i\omega _n)^2}`$. As shown in Appendix B, Eq. (B18),
$`𝒜^1(𝐑)`$ $``$ $`{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{\epsilon _𝐤}{(\omega _n)^2\omega _𝐤^2}}e^{i𝐤𝐑}`$ (43)
$`=`$ $`^1(𝐑){\displaystyle \frac{\omega _n}{2\sqrt{2}\pi D^2}}K_1\left(\left|{\displaystyle \frac{R\omega _n}{D}}\right|\right),`$ (44)
with
$$\begin{array}{cc}& ^1(𝐑)=2\mathrm{sin}\frac{\pi }{2}(R_x+Ry)(\mathrm{cos}\phi +\mathrm{sin}\phi )+2\mathrm{cos}\frac{\pi }{2}(R_xR_y)(\mathrm{cos}\phi \mathrm{sin}\phi ),\end{array}$$
(45)
where $`\phi `$ is the angle between $`𝐑`$ and the x-axis and $`K_1`$ is the Bessel function of rank one. In the same manner
$`^1(𝐑)`$ $``$ $`{\displaystyle \underset{𝐤}{}}{\displaystyle \frac{\mathrm{\Delta }_𝐤}{(\omega _n)^2\omega _𝐤^2}}e^{i𝐤𝐑}`$ (46)
$`=`$ $`^2(𝐑){\displaystyle \frac{\omega _n}{2\sqrt{2}\pi D^2}}K_1\left(\left|{\displaystyle \frac{R\omega _n}{D}}\right|\right),`$ (47)
with
$$\begin{array}{cc}& ^2(𝐑)=2\mathrm{sin}\frac{\pi }{2}(R_x+R_y)(\mathrm{cos}\phi \mathrm{sin}\phi )+2\mathrm{cos}\frac{\pi }{2}(R_xR_y)(\mathrm{cos}\phi +\mathrm{sin}\phi ).\end{array}$$
(48)
Note that the point $`𝐑=0`$ is rather special with $`𝒜^0(\mathrm{𝟎})=\frac{4i\omega _n}{4\pi D^2}\mathrm{ln}\left|\omega _n/D\right|`$ and with $`𝒜^1(\mathrm{𝟎})=^1(\mathrm{𝟎})=0`$.
## V Evaluation of the density of states
### A Unitary limit and low energies
In this section we will evaluate the leading term in the density of states in the limit of low frequencies. When $`0<R|\omega _n|D`$ we have
$$\begin{array}{cc}K_0\left(\frac{R|\omega _n|}{D}\right)\mathrm{ln}\left(\frac{R|\omega _n|}{D}\right),& K_1\left(\frac{R|\omega _n|}{D}\right)\frac{D}{|\omega _n|R},\end{array}$$
(49)
so that in the limit $`|\omega _n|/D0`$ we have
$$\left|𝒜^0(𝐑)\right|\left|𝒜^1(𝐑)\right|.$$
(50)
In the limit of low frequencies and for $`𝐑0`$, this enables us to neglect $`𝒜^0(𝐑)`$ as compared to $`𝒜^1(𝐑)`$ in the evaluation of the matrix $`\widehat{M}`$ in Eq. 25. For $`𝐑=\mathrm{𝟎}`$, $`𝒜^0(\mathrm{𝟎})`$ and $`𝒜^1(\mathrm{𝟎})`$ are negligible as compared to $`𝒜^1(𝐑)`$ for $`𝐑\mathrm{𝟎}`$. In the sequel we can safely avoid the point $`𝐑=\mathrm{𝟎}`$ in summations over $`𝐑`$ that occur during the evaluation of $`\widehat{M}^2`$. Second we notice that the Bessel function $`K_0`$ and $`K_1`$ have an exponential cut-off at $`R_{max}=D/|\omega _n|`$ so that we can safely use the approximation:
$$𝒜^1(𝐑_{ij})\{\begin{array}{cc}\frac{^1(𝐑_{ij})}{2\sqrt{2}\pi DR_{ij}},\text{if}\hfill & R<D/|\omega _n|,\hfill \\ 0,\hfill & \text{elsewhere},\hfill \end{array}$$
(51)
$$^1(𝐑_{ij})\{\begin{array}{cc}\frac{^2(𝐑_{ij})}{2\sqrt{2}\pi DR_{ij}},\text{if}\hfill & R<D/|\omega _n|,\hfill \\ 0,\hfill & \text{elsewhere}.\hfill \end{array}$$
(52)
An important remark to make is that the $`\omega `$-dependence of the matrix elements of $`\widehat{M}`$ appears only through the upper cut-off of the Bessel functions. In what follows, we make the crucial assumption of the unitary limit, i.e., that $`V_0\mathrm{}`$. Recalling the form of $`\widehat{M}`$ in equation (25) we see that in this limit the identity matrix in $`\widehat{M}`$ becomes negligible when compared to $`A_{ij}`$ and $`C_{ij}`$.
### B The divergences appear
Our aim is now to factorize the leading divergences in this problem. In order to make the divergence apparent, it is more convenient to work with
$$\delta G(i\omega _n)=Tr\left[\widehat{M}^2\frac{\widehat{M}^2}{2i\omega _n}\right].$$
(53)
From section V A, we expect logarithmic factors $`\mathrm{ln}\left|D/\omega _n\right|`$ to appear. An important point to stress is that for any configuration of the impurities, we will always find some factors $`\mathrm{ln}\left|D/\omega _n\right|`$ in $`\widehat{M}^2`$. Within the unitary approximation we have
$$\widehat{M}^2=\left[\begin{array}{cc}\widehat{A}^2+\widehat{B}^2& \widehat{A}\widehat{B}\widehat{B}\widehat{A}\\ \widehat{A}\widehat{B}\widehat{B}\widehat{A}& \widehat{A}^2+\widehat{B}^2\end{array}\right].$$
(54)
To see that $`\mathrm{ln}\left|D/\omega _n\right|`$ is necessarily present in the diagonal terms of $`\widehat{M}^2`$, we estimate
$`M_{ii}^2`$ $`=`$ $`{\displaystyle \underset{j}{}}\left(A_{ij}A_{ji}+B_{ij}B_{ji}\right)`$ (55)
$`=`$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{(2\pi D)^2}}{\displaystyle \frac{\left[(^1)^2+(^2)^2\right]}{R_{ij}^2}},`$ (56)
where the summation over $`j`$ is restricted to $`0<R_{ij}<D/|\omega _n|`$. Provided the impurities are rather homogeneously scattered in the system (around each impurity site $`𝐑_i`$ one can find a macroscopic amount of impurities inside a circle of radius $`D/|\omega _n|`$, cf. figure 2 ), we can take the continuous limit,
$$\widehat{M}_{ii}^2\frac{2\pi V_0}{(2\pi D)^2}_1^{D/|\omega _n|}\frac{(^1)^2+(^2)^2}{R}𝑑R.$$
(57)
Now $`(^1)^2+(^2)^2`$ are oscillatory but always positive so for all impurity at site $`𝐑_i`$ we have
$$\widehat{M}_{ii}^2=C\mathrm{ln}\left|\frac{D}{\omega _n}\right|,$$
(58)
where $`C`$ is a constant.
The biggest coefficients in $`\widehat{M}^2`$ are situated on the diagonal. To see it, we distinguish between the off-diagonal elements that are diagonal in the particle-hole grading and those that are not. The magnitude of off-diagonal elements that are diagonal in the particle-hole grading can be written
$$\widehat{M}_{ik}^2=\underset{j}{}\frac{1}{2\pi D}\frac{\left[_{ij}^1_{jk}^1+_{ij}^2_{jk}^2\right]}{R_{ij}R_{jk}}.$$
(59)
As $`ik`$, $`\widehat{M}_{ik}^2`$ picks up an oscillatory prefactor (a combination of $`e^{\pm i\pi /2R_{ik}}e^{\pm i\phi }`$ as in eqns. (LABEL:pref1), (LABEL:pref2), (LABEL:pref3) ). Furthermore, the logarithmic divergence gets rescaled by $`R_{ik}`$:
$$\widehat{M}_{ik}^2=COsc(𝐑_{ik})\mathrm{ln}\left|\frac{DR_{ik}}{\omega _n}\right|,$$
(60)
where $`\left|Osc(𝐑_{ik})\right|1`$. Hence for all sites $`ki`$ we have
$$\left|\frac{\widehat{M}_{ik}^2}{\widehat{M}_{ii}^2}\right|1.$$
(61)
In turn, the terms in the off-diagonal blocks of $`\widehat{M}^2`$ with respect to the particle-hole vanish. Indeed, if we denote the off-diagonal elements of $`\widehat{M}^2`$ with respect to the particle-hole grading by
$$T_{ik}A_{ij}B_{jk}B_{ij}A_{jk},$$
(62)
we have
$$T_{ik}=\underset{j}{}\frac{1}{(2\pi D)^2}\left(\frac{^1(𝐑_{ij})^2(𝐑_{jk})}{R_{ij}R_{jk}}\frac{^2(𝐑_{ij})^1(𝐑_{jk})}{R_{ij}R_{jk}}\right),$$
(63)
where here the sum runs over the points $`𝐑_j`$ such that $`R_{ij}<D/|\omega _n|`$ and $`R_{jk}<|\omega _n|`$. Noticing the symmetry of $`^1`$ and $`^2`$ under the transformation $`\phi \phi +\pi `$, we can show that the two terms on the r. h. s. of Eq. (63) cancel identically. Indeed the integrands $`\frac{^1(R_{ij})}{R_{ij}}`$ and $`\frac{^2(R_{ik})}{R_{ik}}`$ are represented respectively within each circle of figure 3 ( both of these terms have a cut-off at $`R_{max}=D/|\omega _n|`$). The summation zone is the surface of intersection of the two disks. The first term in (63) is represented in the upper drawing whereas the second one is represented in the lower drawing. For each summation point $`𝐑_j`$ in the upper drawing there is a symmetric one $`\stackrel{~}{𝐑}_j`$ in the lower drawing such that $`𝐑_{ij}=\stackrel{~}{𝐑}_{jk}`$ and $`\stackrel{~}{𝐑}_{ij}=𝐑_{jk}`$. Thus $`T_{ik}=0`$.
In conclusion we find it convenient to define a matrix $`\widehat{S}`$ by
$$\widehat{M}^2=C\mathrm{ln}\left|\frac{D}{\omega _n}\right|\widehat{S},$$
(64)
where $`C`$ is the constant defined in (58), independent of disorder. The matrix $`\widehat{S}`$ depends on the particular configuration of the impurities in the system. It satisfies
$$\left|S_{ij}\right|1.$$
(65)
### C Asymptotic value of the density of states
Substituting the value of $`\widehat{M}^2`$ from (64) into the equation (53) we find
$`\delta G(i\omega _n)`$ $`=`$ $`𝒯_{div}(i\omega _n)+(i\omega _n),`$ (66)
where $``$ denotes the average over disorder, $`\widehat{I}`$ is the $`2N\times 2N`$ density matrix, and
$$\begin{array}{cc}𝒯_{div}(i\omega _n)\frac{1}{2𝒱}\frac{}{i\omega _n}\mathrm{ln}\mathrm{ln}(D/|\omega _n|)Tr\widehat{I},|calR(i\omega _n)\frac{1}{2𝒱}Tr\widehat{S}^1\frac{}{i\omega _n}\widehat{S}.& \end{array}$$
(67)
Then the first term in Eq. (66) $`𝒯_{div}`$ is responsible for the singular density of states that we obtain. Indeed it gives rise to
$$𝒯_{div}(i\omega _n)=\frac{N}{𝒱}\frac{1}{i\omega _n\mathrm{ln}|\omega _n/D|},$$
(68)
where we recall that $`N`$ is the number of impurities. After analytic continuation (remember that $`|\omega _n|=\sqrt{(i\omega _n)^2}`$, Eq. (B18) ) and assuming that the reminder $``$ in Eq. (66) is negligible, we get
$$\delta \rho (\omega )\frac{1}{\pi }n_iIm\left[\frac{1}{(\omega +i\delta )\mathrm{ln}\left(\frac{i\omega +\delta }{D}\right)}\right],$$
(69)
and thus
$$\begin{array}{cc}& \delta \rho (\omega )\frac{n_i}{2}\frac{1}{|\omega |\left[\mathrm{ln}^2(|\omega |/D)+(\pi /2)^2\right]},\hfill \end{array}$$
(70)
where $`n_i`$ is the density of impurities in the system. We note that this expression is normalizable:
$$_D^D\delta \rho (\omega )𝑑\omega =2n_i.$$
(71)
In order to prove the result (70) we still have to show that the reminder $``$ in (53) is negligible as compared to $`𝒯_{div}`$. In order to do this we have to give some insight about the form of $`\widehat{S}^1`$.
### D The form of $`\widehat{S}^1`$
The matrix $`\widehat{S}`$ is invertible (since $`\widehat{M}`$ is invertible) and we will find a reasonable candidate to the inverse of $`\widehat{S}`$ in order to give an estimation of the reminder $``$. Inverting $`\widehat{S}`$ means we can find a matrix $`\widehat{S}^1`$ such that for any given pair of sites $`i`$ and $`k`$
$$\underset{j}{}\widehat{S}_{ij}\widehat{S}_{jk}^1=\delta _{ik}.$$
(72)
We introduce a pictorial representation of $`\widehat{S}`$ by drawing a disk ( called $`\widehat{S}`$-disk) of radius $`|D/\omega _n|`$ centered at $`𝐑_i`$. To each location $`𝐑_j`$ of an impurity there corresponds a matrix element $`S_{ij}`$ which depends on the vector $`𝐑_{ij}=𝐑_i𝐑_j`$. From Eqs. (51,52) we recall that
$$\{\begin{array}{ccc}|S_{ij}|1,\hfill & \text{ for}& R_{ij}<|D/\omega _n|,\hfill \\ S_{ij}=0,\hfill & \text{for}& R_{ij}|D/\omega _n|,\hfill \end{array}$$
(73)
inside the disk, $`\widehat{S}`$ has some non vanishing matrix elements $`|S_{ij}|<1`$. Outside the disk, $`S_{ij}=0`$. It’s important to note at this point that the only dependence on $`\omega `$ in the $`\widehat{S}`$ matrix comes from the cut-off. In order to satisfy (72) we must presume that $`\widehat{S}^1`$ has the same cut-off $`|D/\omega _n|`$ as $`\widehat{S}`$. Hence we represent again $`\widehat{S}_{ik}^1`$ by a disk (called $`\widehat{S}^1`$-disk), but centered around $`𝐑_k`$ this time. Now the summation $`_j\widehat{S}_{ij}\widehat{S}_{jk}^1`$ runs over the intersection of the $`\widehat{S}`$-disk and the $`\widehat{S}^1`$-disk. The key difficulty in order to invert $`\widehat{S}`$ is to find a matrix $`\widehat{S}^1`$, such that when the two disks have the same center we have
$`{\displaystyle \underset{j}{}}\widehat{S}_{ij}\widehat{S}_{ji}^1=1,`$whereas when the two centers differ, even by a small amount, we have
$`\begin{array}{cc}_j\widehat{S}_{ij}\widehat{S}_{jk}^1=0,\hfill & ki.\hfill \end{array}`$This is illustrated on figure 4 where on the left side (case (a)) the two circles are centered at the same point and on the right side (case (b)) the two centers differ by a tiny amount. In both cases the intersecting area of the two disks is almost identical, but in case (a) the result has to be $`1`$ whereas in case (b) it has to be $`0`$.
The matrix elements of $`\widehat{S}`$ inside the disk are random, but the condition $`|\widehat{S}_{ij}|1`$ is independent of the realization of disorder. The worst possible situation for differentiating between cases (a) and (b) is when all the matrix elements inside the $`\widehat{S}`$-disk have their maximum value $`1`$.
We define the external boundary of the $`\widehat{S}^1`$-disk. By external boundary we mean the circle exactly adjacent to the disk and external to it. Upon the external boundary the matrix elements of $`\widehat{S}^1`$ are defined as non vanishing and negative, so that the summation over it compensates the summation over the intersection of the $`\widehat{S}`$-disk and the $`\widehat{S}^1`$-disk. In case (a) the external boundary doesn’t touch the $`\widehat{S}`$-disk and thus has no effect upon the summation over the intersection of the $`\widehat{S}`$-disk and the $`\widehat{S}^1`$-disk. Alternatively in case (b) the summations over the intersection of the $`\widehat{S}`$-disk and the $`\widehat{S}^1`$-disk and over the external boundary cancel out.
Let’s take an explicit example where $`S_{ij}=1`$ if $`R_{ij}<|D/\omega _n|`$ and $`S_{ij}=0`$ elsewhere. We define $`\widehat{S}_{ij}^1`$ in the following way. For $`R_{ij}<|D/\omega _n|`$, $`\widehat{S}_{ij}^1`$ is proportional to a random configuration of $`\pm 1`$ such that after integration over a disk of volume $`V=\pi \left|D/\omega _n\right|^2`$ we get
$`{\displaystyle \underset{V}{}}\pm 1`$ $`=`$ $`\sqrt{V}`$ (74)
$`=`$ $`\sqrt{\pi }{\displaystyle \frac{D}{|\omega _n|}}.`$ (75)
By the central limit theorem, there are many random configurations of $`\pm 1`$ that verify this condition. We then take the external boundary to be proportional to $`(1/\sqrt{\pi })`$ with the same proportionality constant. Thus
$$\{\begin{array}{cc}\widehat{S}_{ij}^1=A(\pm 1),\hfill & \text{if}R_{ij}<|D/\omega _n|,\hfill \\ \widehat{S}_{ij}^1=A\frac{1}{\sqrt{\pi }},\hfill & \text{if}R_{ij}=|D/\omega _n|,\hfill \\ \widehat{S}_{ij}^1=0,\hfill & \text{if}R_{ij}>|D/\omega _n|,\hfill \end{array}$$
(76)
where $`A`$ is a constant. We notice that when $`𝐑_i`$ and $`𝐑_k`$ are infinitely close, then
$`{\displaystyle \underset{boundary}{}}\widehat{S}_{ij}^1=A\sqrt{\pi }{\displaystyle \frac{D}{|\omega _n|}},`$and exactly compensates the summation over the volume inside. The proportionality constant $`A`$ is fixed so that
$`{\displaystyle \underset{j}{}}\widehat{S}_{ij}\widehat{S}_{ij}^1=1,`$
$`\begin{array}{cc}\text{thus}\hfill & A={\displaystyle \frac{|\omega _n|}{\sqrt{\pi }D}}.\hfill \end{array}`$
In summary
$$\{\begin{array}{cc}\widehat{S}_{ij}^1=\frac{|\omega _n|}{\sqrt{\pi }D}(\pm 1),\hfill & \text{if}R_{ij}<|D/\omega _n|,\hfill \\ \widehat{S}_{ij}^1=\frac{|\omega _n|}{\pi D},\hfill & \text{if}R_{ij}=|D/\omega _n|,\hfill \\ \widehat{S}_{ij}^1=0,\hfill & \text{if}R_{ij}>|D/\omega _n|.\hfill \end{array}$$
(77)
Now for each intermediate case where $`ki`$ (cf. figure 5) we want to be sure that the intersection of the $`\widehat{S}`$-disk and the $`\widehat{S}^1`$-disk compensates the sum over the external boundary of $`\widehat{S}^1`$ which crosses the $`\widehat{S}`$-disk. This is obviously not the case for any configuration of random $`\pm 1`$ in $`\widehat{S}^1`$ but we believe there is one (and actually only one because there is only one inverse for $`\widehat{S}`$!) configuration which satisfies it for all positions of $`𝐑_i`$ and $`𝐑_k`$.
As represented in figure V D we call $`\theta `$ the angle made by the center $`i`$ with the two points $`A`$ and $`B`$, intersection between the two circles. The area of summation is given by $`A_\theta =\theta R^24R^2\mathrm{sin}\theta `$ and the intersecting arc’s length by $`L_\theta =\theta R`$. They both scale in respectively $`R^2`$ and $`R`$ with varying prefactors. So the desired configuration of random $`\pm 1`$ in the $`\widehat{S}^1`$-disk has to be “denser” towards the center of the circle than towards the boundary. We are not able to write down explicitly this configuration of random $`\pm 1`$ inside the area of the $`\widehat{S}^1`$-disk, but actually it doesn’t matter, because as we will see in the next paragraph the evaluation of the reminder doesn’t depend on it.
### E Evaluation of the reminder $``$
Now we evaluate
$$=\frac{1}{2𝒱}Tr\widehat{S}^1\frac{}{i\omega _n}\widehat{S}.$$
(78)
An important point is that as $`\widehat{S}`$ depends on $`\omega _n`$ only via its boundary, $`\widehat{S}/\omega _n`$ is a matrix with non zero values only on the external boundary of the $`\widehat{S}`$-disk. explicitly, taking our example where $`\widehat{S}_{ij}=1`$ if $`R_{ij}<D/|\omega _n|`$ we get
$$\{\begin{array}{cc}\frac{\widehat{S}}{\omega _n}=1,\hfill & \text{if}R_{ij}=D/|\omega _n|,\hfill \\ \frac{\widehat{S}}{\omega _n}=0,\hfill & \text{elsewhere}.\hfill \end{array}$$
(79)
Since the matrix elements in the $`\widehat{S}^1`$-disk have random (positive and negative) signs whereas upon the external boundary they have constant sign (negative here), the maximum value of $`\left|\widehat{S}^1\widehat{S}/\omega _n\right|`$ is reached when $`i=k`$, that is when the external boundary of the $`\widehat{S}`$-disk (where the matrix elements of $`\widehat{S}/\omega _n`$ are non vanishing) matches the external boundary of the $`\widehat{S}^1`$-disk (cf. Fig. 6).
In our special case we get
$`\left|\widehat{S}^1{\displaystyle \frac{\widehat{S}}{\omega _n}}\right|`$ $`=`$ $`\left({\displaystyle \frac{|\omega _n|}{\pi D}}\right)\left({\displaystyle \frac{2\pi D}{|\omega _n|}}\right),`$ (80)
$`=`$ $`2.`$ (81)
Since $`\left|\widehat{S}/\omega _n\right|`$ has its maximum value when all the matrix elements of $`\widehat{S}`$ inside $`\widehat{S}`$-disk equal $`1`$, we have indeed evaluated an upper bound for the reminder $``$. The reminder $``$ is thus negligible compared to the leading divergence in the density of states in the limit $`|\omega _n|/D1`$.
## VI Discussion
The first question to address is what happens when the different conditions under which our calculation was performed are relaxed. First consider the more realistic situation where the bandwidth $`W`$ and the superconducting gap $`\mathrm{\Delta }_0`$ are not equal (experimentally, we have $`\mathrm{\Delta }_00.10W`$). According to Ref. the power law dependence of the matrix elements of $`\widehat{M}`$ (cf. equations (51) and (52)) is still preserved, but gets an overall prefactor of $`\mathrm{\Delta }_0/W`$. The scale under which our calculation is valid is now $`|\omega |V_0\mathrm{\Delta }_0/W`$. In addition, the hopping matrix $`\widehat{M}`$ will show a strong spatial anisotropy . This anisotropy can be absorbed into the overall prefactors $`^0`$, $`^1`$ and $`^2`$ (resp. eqns. (LABEL:pref1), (LABEL:pref2) and (LABEL:pref3)) entering the definition of the matrix elements $`A_{ij}`$, $`C_{ij}`$ and $`B_{ij}`$. Since only the square of these factors enter the leading divergence (cf. Eq. (55)), we believe the result stays unchanged.
What happens if the bandwidth are not symmetric anymore, that is if $`\mu 0`$ but still $`\mu \mathrm{\Delta }_0`$ ?
First the nodes are moved away from the point $`(\pm \pi /2,\pm \pi /2)`$ so that transversal nodes are now separated by the vectors $`𝐐=(\pi (1\delta ),\pi (1\delta ))`$ and $`𝐐^{}=(\pi (1\delta ),\pi (1\delta ))`$ where $`\delta =\mu /\mathrm{\Delta }_0`$ and $`\mu `$ is the increase in the chemical potential. This leads to a change of the phase factors in $`A_{ij}`$ and $`B_{ij}`$. Namely we get $`^0(R)=2\left[\mathrm{cos}\left(𝐐𝐑/2\right)+\mathrm{cos}\left(𝐐^{}𝐑/2\right)\right]`$;
$`^1(R)`$ $`=`$ $`2[\mathrm{sin}(𝐐𝐑/2)(\mathrm{cos}\phi +\mathrm{sin}\phi ),`$ (82)
$``$ $`\mathrm{sin}(𝐐^{}𝐑/2)(\mathrm{cos}\phi \mathrm{sin}\phi )],`$ (83)
and
$`^2(R)`$ $`=`$ $`2[\mathrm{sin}(𝐐𝐑/2)(\mathrm{cos}\phi \mathrm{sin}\phi ),`$ (84)
$``$ $`\mathrm{sin}(𝐐^{}𝐑/2)(\mathrm{cos}\phi +\mathrm{sin}\phi )].`$ (85)
As $`\delta `$ is a small parameter, this change in the phase won’t affect the existence of the logarithmic divergence in $`\widehat{M}^2`$. Additionally, away from half filling, the bands of quasi particles and quasi holes become asymmetric to account for the removing of particles in the system. The difference induced in $`A_{ij}`$, $`B_{ij}`$ and $`C_{ij}`$ comes from the highest part of the energy spectrum, where $`k1`$ and $`\omega D`$ and shouldn’t affect our result.
Our solution is valid under the assumption of unitary limit, meaning that $`V_0`$ is the largest scale in the problem. As shown in , to non interacting single impurities is associated the creation of a bound states decaying as $`1/(R\mathrm{ln}R)`$. Following Balatsky et al. () this would lead to delocalization due to the formation of an impurity band. The result we find for the density of states is indeed reminiscent of a Dyson-like singularity ( in d=1, a Dyson-like singularity corresponds to a density of state diverging in $`1/\left(|\omega |\mathrm{ln}^3(|\omega |/D)\right)`$ ), associated in one dimension with delocalization .
What happens when the condition of unitarity is relaxed is still very much an open problem. In the case of weak disorder, some $`\sigma `$-model analysis have been performed concluding to a vanishing density of state under an energy scale $`E_2=D_F/\xi ^2`$, where $`D_F`$ is the bare diffusion constant and $`\xi `$ is the localization length. This result is strongly supported by symmetry considerations. Indeed, a disordered d-wave superconductor belongs to class C<sub>I</sub>, according to the classification of ref. , meaning that the Hamiltonian is invariant under time-reversal symmetry as well as spin rotation symmetry. According to random matrix theory a universal behavior is expected, inducing a vanishing density of states on the scale of the level spacing induced by the finite localazation length.
Consider one impurity with scattering potential $`V_0`$. The effective potential at the impurity site can be evaluated exactly and is given by $`\overline{V}=\frac{V_0}{1V_0|\omega |\mathrm{ln}\left|D/\omega \right|}`$. In the unitary limit ( $`V_0\mathrm{}`$) we get $`\overline{V}=\frac{1}{|\omega |\mathrm{ln}\left|D/\omega \right|}`$. This effective potential diverges when $`\omega `$ goes to zero. On the other hand, we notice that the derivation of non linear sigma models for disordered systems requires Gaussian disorder, and especially require that the average effective disorder potential vanishes $`V=0`$ (the second moment is nonzero). If $`V`$ is nonzero but is a constant as a function of energy, it can be absorbed as a redefinition of the chemical potential, but the case where the effective potential would diverge as $`\omega `$ goes to zero belongs to another universality class: the energy of an effective non-linear sigma model would renormalized to $`\omega \mathrm{\Sigma }(\omega )`$, where $`\mathrm{\Sigma }(\omega )`$ diverges when $`omega`$ goes to zero. The energy scale under which the non-linear sigma model describes the diffusive modes then becomes inaccessible since the effective energy never gets close to zero. In our problem the result obtained on the density of states indicates that the self-energy is of the form $`\mathrm{\Sigma }(\omega )=\frac{1}{\left|\omega \mathrm{ln}^2(\omega /D)\right|}`$, diverging as $`\omega `$ goes to zero, but still different from the one impurity case. The feynmann diagrams leading to such this self-energy will be studied in a future publication . We believe the method presented here, using the T-matrix equation takes care in a non perturbative way of the leading divergence in the unitary limit. One possible scenario which would reconciliate the two limits of weak and strong disorder is that the unitarity limit fixed point is unstable (as soon as $`V_0`$ becomes finite, the effective potential saturates), but the cross-over regime close to unitarity is still very much influenced by the strong disorder fixed point.
We would like to thank A.V. Balatsky, P. Coleman, M. Hettler, R. Joynt, C. Mudry, R. Narayanan, A. M. Tsvelik, T. Xiang for useful discussions related to this work. We are especially grateful to J. Chalker and B.D. Simon for discussions concerning symmetries of disordered d-wave superconductors. As this paper was submitted, a recent numerical study confirmed the singular density of states in the unitary limit. However they concluded that the resonant density of states alone was not sufficient to induce delocalisation.
This work is supported by NSF Grant No. DMR 9813764 and by (CP) a Bourse Lavoisier and the research fund from the EPSRC, UK.
## A Derivation of the T-matrix Equation
Starting form equation (14), replacing the lagrangian of (15) into it and multiplicating from the left by $`G^0`$ gives
$`G_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}`$ $`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}e^{i𝐤𝐑_i}(1)^\nu t_{𝐤\nu }G_{𝐤\nu }^0G_{𝐤^{}\nu ^{}}^1(𝐑_i)`$ (A1)
$`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}e^{i𝐤𝐑_i}t_{𝐤\nu +1}G_{𝐤\nu }^0G_{𝐤^{}\nu ^{}}^2=G_{𝐤\nu }^0\delta _{\mathrm{𝐤𝐤}^{}}\delta _{\nu \nu ^{}},`$ (A2)
where
$$\begin{array}{c}G_{𝐤\nu }^1(𝐑_i)_{𝐪m}e^{i𝐪𝐑_i}(1)^mt_{𝐪m}G_{\mathrm{𝐪𝐤}}^{m\nu },\hfill \\ G_{𝐤\nu }^2(𝐑_i)_{𝐪m}e^{i𝐪𝐑_i}t_{𝐪m+1}G_{\mathrm{𝐪𝐤}}^{m\nu }.\hfill \end{array}$$
(A3)
Now we have two unknown functions $`G^1`$ and $`G^2`$ that we evaluate using formula (A1):
$`G_{𝐤\nu }^1(𝐑_j)`$ $`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}{\displaystyle \underset{𝐪m}{}}e^{i𝐪(𝐑_i𝐑_j)}(t_{𝐪m})^2G_{𝐪m}^0G_{𝐤\nu }^1(𝐑_i)`$ (A4)
$`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}{\displaystyle \underset{𝐪m}{}}e^{i𝐪(𝐑_i+𝐑_j)}(1)^mt_{𝐪m}t_{𝐪m+1}G_{𝐪m}^0G_{𝐤\nu }^2(𝐑_i)=(1)^\nu e^{i𝐤𝐑_j}t_{𝐤\nu }G_{𝐤\nu }^0;`$ (A5)
$`G_{𝐤\nu }^2(𝐑_j)`$ $`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}{\displaystyle \underset{𝐪m}{}}e^{i𝐪(𝐑_i𝐑_j)}(1)^mt_{𝐪m}t_{𝐪m+1}G_{𝐪m}^0G_{𝐤\nu }^1(𝐑_i)`$ (A6)
$`+`$ $`{\displaystyle \frac{V_0}{𝒱}}{\displaystyle \underset{i}{}}{\displaystyle \underset{𝐪m}{}}e^{i𝐪(𝐑_i+𝐑_j)}(t_{𝐪m})^2G_{𝐪m}^0G_{𝐤\nu }^2(𝐑_i)=e^{i𝐤𝐑_j}t_{𝐤\nu +1}G_{𝐤\nu }^0.`$ (A7)
These two equations can be rewritten matricially as
$$\begin{array}{c}\left(\delta _{ij}+V_0A_{ij}\right)G_{𝐤\nu }^1(𝐑_j)+V_0B_{ij}G_{𝐤\nu }^2(𝐑_j)=N_{𝐤\nu }^1(𝐑_j),\hfill \\ \left(\delta _{ij}+V_0C_{ij}\right)G_{𝐤\nu }^2(𝐑_j)+V_0B_{ij}G_{𝐤\nu }^1(𝐑_j)=N_{𝐤\nu }^2(𝐑_j).\hfill \end{array}$$
(A8)
Define the $`2N`$ vector $`𝐕`$
$$\begin{array}{ccc}𝐕\left(\begin{array}{c}𝐕^1\\ 𝐕^2\end{array}\right),& 𝐕_{𝐤\nu }^1\left(\begin{array}{c}G_{𝐤\nu }^1(𝐑_1)\\ \mathrm{}\\ G_{𝐤\nu }^1(𝐑_N)\end{array}\right),& 𝐕_{𝐤\nu }^2\left(\begin{array}{c}G_{𝐤\nu }^2(𝐑_1)\\ \mathrm{}\\ G_{𝐤\nu }^2(𝐑_N)\end{array}\right),\end{array}$$
(A9)
and the equation (A8) can be written
$$\widehat{M}𝐕_{𝐤\nu }=𝐍_{𝐤\nu }.$$
(A10)
But then equation (A1) becomes
$$G_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}+\frac{V_0}{𝒱}𝐍_{𝐤\nu }𝐕_{𝐤\nu }=G_{𝐤\nu }^0\delta _{\mathrm{𝐤𝐤}^{}}\delta _{\nu \nu ^{}}.$$
(A11)
Insertion of (A10) yields
$$G_{\mathrm{𝐤𝐤}^{}}^{\nu \nu ^{}}=G_{\mathrm{𝐤𝐤}^{}}^0\frac{V_0}{𝒱}𝐍_{𝐤\nu }^T\widehat{M}^1𝐍_{𝐤^{}\nu ^{}}.$$
(A12)
## B Calculation of the matrix elements $`A_{ij}`$, $`B_{ij}`$ and $`C_{ij}`$
$`A_{ij}`$ $`=`$ $`{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2k}{(2\pi )^2}e^{i𝐤(𝐑_i𝐑_j)}(t_{𝐤n})^2G_{𝐤n}^0}`$ (B1)
$`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}e^{i𝐤𝐑_{ij}}\left(\frac{u_𝐤^2}{i\omega _n\omega _𝐤}+\frac{v_𝐤^2}{i\omega _n+\omega _𝐤}\right)},`$ (B2)
where the integration runs over the first Brillouin zone. Since
$$\begin{array}{c}u_𝐤^2=\frac{1}{2}\left(1+\frac{\epsilon _𝐤}{\omega _𝐤}\right),\hfill \\ v_𝐤^2=\frac{1}{2}\left(1\frac{\epsilon _𝐤}{\omega _𝐤}\right),\hfill \end{array}$$
(B3)
we get
$$A_{ij}=𝒜^0(𝐑_{ij})+𝒜^1(𝐑_{ij}),$$
(B4)
with
$$\begin{array}{c}𝒜^0(𝐑_{ij})i\omega _n\frac{d^2k}{(2\pi )^2}\frac{e^{i𝐤𝐑_{ij}}}{(i\omega _n)^2\omega _𝐤^2},\hfill \\ 𝒜^1(𝐑_{ij})\frac{d^2k}{(2\pi )^2}\frac{\epsilon _𝐤}{(i\omega _n)^2\omega _𝐤^2}e^{i𝐤𝐑_{ij}}.\hfill \end{array}$$
(B5)
Similarly
$$C_{ij}=𝒜^0(𝐑_{ij})𝒜^1(𝐑_{ij}),$$
(B6)
and
$`B_{ij}`$ $`=`$ $`{\displaystyle \frac{d^2k}{(2\pi )^2}\frac{\mathrm{\Delta }_𝐤}{(i\omega _n)^2\omega _𝐤^2}e^{i𝐤𝐑_{ij}}}`$ (B7)
$``$ $`^1(𝐑_{ji}).`$ (B8)
### 1 Evaluation of $`𝒜^0(𝐑_{ji})`$
As we can see on figure 1 the spectrum has four nodes at the points
$$\begin{array}{cccc}P_1=(\frac{\pi }{2},\frac{\pi }{2}),& P_2=(\frac{\pi }{2},\frac{\pi }{2}),& P_3=(\frac{\pi }{2},\frac{\pi }{2}),& P_4=(\frac{\pi }{2},\frac{\pi }{2}).\end{array}$$
(B9)
Under the assumptions $`\mu =0`$ and $`\mathrm{\Delta }_0=W`$, we can linearize the spectrum around each node in the following way:
$$\begin{array}{cc}𝐤^{}(\frac{\pi }{2},\frac{\pi }{2})+𝐤,\hfill & \omega _𝐤^{}^2=W^2\left[(\mathrm{cos}k_x^{}+\mathrm{cos}k_y^{})^2+(\mathrm{cos}k_x^{}\mathrm{cos}k_y^{})^2\right],\hfill \end{array}$$
(B10)
and we get
$$\begin{array}{cc}\omega _𝐤^{}^2D^2k^2,\text{with}\hfill & D=\sqrt{2}W.\hfill \end{array}$$
(B11)
Similarly,
$$\begin{array}{cc}\epsilon _𝐤^{}W(k_x+k_y),\hfill & \\ \mathrm{\Delta }_𝐤^{}(k_xk_y)\mathrm{\Delta }_0.\hfill & \end{array}$$
(B12)
In order to evaluate $`𝒜^0(𝐑)`$ we divide the integral into a sum of four integrals around each node :
$`𝒜^0(𝐑)`$ $`=`$ $`{\displaystyle \underset{k^{}BZ_1}{}}i\omega _n{\displaystyle \frac{e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}+{\displaystyle \underset{k^{}BZ_2}{}}i\omega _n{\displaystyle \frac{e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}`$ (B13)
$`+`$ $`{\displaystyle \underset{k^{}BZ_3}{}}i\omega _n{\displaystyle \frac{e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}+{\displaystyle \underset{k^{}BZ_4}{}}i\omega _n{\displaystyle \frac{e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}.`$ (B14)
In each term, the fact of developing around a particular node gives a specific prefactor, so that we have
$$𝒜^0(𝐑)=^0(𝐑)(i\omega _n)\underset{k}{}\frac{e^{i𝐤𝐑}}{(i\omega _n)^2D^2k^2},$$
(B15)
with
$$\begin{array}{cc}& ^0(𝐑)=e^{i\frac{\pi }{2}(R_x+R_y)}+e^{i\frac{\pi }{2}(R_xR_y)}+e^{i\frac{\pi }{2}(R_x+R_y)}+e^{i\frac{\pi }{2}(R_xR_y)}.\hfill \end{array}$$
(B16)
Thus
$$^0(𝐑)=2\mathrm{cos}\frac{\pi }{2}(R_x+R_y)+2\mathrm{cos}\frac{\pi }{2}(R_xR_y).$$
(B17)
Now calling $`\theta `$ the angle between $`𝐤`$ and $`𝐑`$ (cf. figure 7),
$`𝒜^0(𝐑)`$ $`=`$ $`^0(𝐑)(i\omega _n){\displaystyle _0^1}{\displaystyle \frac{kdk}{(2\pi )^2}}{\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle \frac{e^{ikR\mathrm{cos}\theta }}{(i\omega _n)^2D^2k^2}}`$ (B18)
$`=`$ $`^0(𝐑){\displaystyle \frac{(i\omega _n)}{2\pi D^2}}{\displaystyle _0^1}k𝑑k{\displaystyle \frac{J_0(kR)}{(\omega _n/D)^2+k^2}}`$ (B19)
$`=`$ $`^0(𝐑){\displaystyle \frac{i\omega _n}{2\pi D^2}}K_0\left(R|\omega _n|/D\right),`$ (B20)
where $`K_0`$ is the Bessel function of rank zero. Note that we have defined $`|\omega _n|=\sqrt{(i\omega _n)^2}`$.
### 2 Calculation of $`𝒜^1(𝐑)`$
As previously, we can decompose the summation in the Brillouin zone into four parts:
$`𝒜^1(𝐑)`$ $`=`$ $`{\displaystyle \underset{k^{}BZ_1}{}}{\displaystyle \frac{\epsilon _k^{}e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}+{\displaystyle \underset{k^{}BZ_2}{}}{\displaystyle \frac{\epsilon _k^{}e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}`$ (B21)
$`+`$ $`{\displaystyle \underset{k^{}BZ_3}{}}{\displaystyle \frac{\epsilon _k^{}e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}+{\displaystyle \underset{k^{}BZ_4}{}}{\displaystyle \frac{\epsilon _k^{}e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}.`$ (B22)
If we call $`\theta `$ the angle between $`𝐤`$ and $`𝐑`$ and $`\phi `$ the angle between $`𝐑`$ and the $`x`$-axis as represented on figure 7, the first term in (B21) can be written
$`T_1`$ $``$ $`{\displaystyle \underset{k^{}BZ_1}{}}{\displaystyle \frac{\epsilon _k^{}e^{i𝐤^{}𝐑}}{(i\omega _n)^2\omega _𝐤^{}^2}}`$ (B23)
$`=`$ $`e^{i\frac{\pi }{2}(R_x+R_y)}{\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle _0^1}k𝑑k{\displaystyle \frac{D}{\sqrt{2}(2\pi )^2}}{\displaystyle \frac{(k_x+k_y)e^{ikR\mathrm{cos}\theta }}{(i\omega _n)^2D^2k^2}}.`$ (B24)
Using
$$\begin{array}{c}k_x=k(\mathrm{cos}\theta \mathrm{cos}\phi \mathrm{sin}\theta \mathrm{sin}\phi ),\hfill \\ k_y=k(\mathrm{cos}\theta \mathrm{sin}\phi +\mathrm{sin}\theta \mathrm{sin}\phi ),\hfill \end{array}$$
(B25)
we get
$`T_1`$ $`=`$ $`e^{i\frac{\pi }{2}(R_x+R_y)}(\mathrm{cos}\phi +\mathrm{sin}\phi ){\displaystyle \frac{D}{\sqrt{2}(2\pi ^2)}}{\displaystyle _0^1}k^2𝑑k{\displaystyle _0^{2\pi }}𝑑\theta {\displaystyle \frac{(\mathrm{cos}\theta +\mathrm{sin}\theta )}{(i\omega _n)^2D^2k^2}}e^{ikR\mathrm{cos}\theta }`$ (B26)
$`=`$ $`e^{i\frac{\pi }{2}(R_x+R_y)}(\mathrm{cos}\phi +\mathrm{sin}\phi ){\displaystyle \frac{iD}{\sqrt{2}2\pi }}{\displaystyle _0^1}k^2𝑑k{\displaystyle \frac{J_1(kR)}{(i\omega _n)^2D^2k^2}}e^{ikR\mathrm{cos}\theta }`$ (B27)
$`=`$ $`{\displaystyle \frac{i}{2\sqrt{2}\pi }}e^{i\frac{\pi }{2}(R_x+R_y)}{\displaystyle \frac{\omega _n}{D^2}}K_1\left({\displaystyle \frac{R|\omega _n|}{D}}\right).`$ (B28)
Thus, after summing over the four nodes we get,
$$𝒜^1(𝐑)=\frac{^1(𝐑)}{2\sqrt{2}\pi }\frac{\omega _n}{D^2}K_1\left(\frac{R|\omega _n|}{D}\right),$$
(B29)
with
$$^1=2\left[\mathrm{sin}\frac{\pi }{2}(R_x+R_y)\left(\mathrm{cos}\phi +\mathrm{sin}\phi \right)+\mathrm{sin}\frac{\pi }{2}(R_xR_y)\left(\mathrm{cos}\phi \mathrm{sin}\phi \right)\right].$$
(B30)
The evaluation of $`^1(𝐑)`$ is done in the same way as the one of $`𝒜^1(𝐑)`$. |
warning/0002/astro-ph0002046.html | ar5iv | text | # Effects of massive star formation on the ISM of dwarf galaxies
## 1 Introduction
To construct a comprehensive picture of a galaxy’s history, understanding the distribution of its energy budget is a fundamental step. For this we must consider observations covering several characteristic wavelength regimes, thus, sampling the various components of the interstellar medium (ISM). While the UV to NIR wavelength continua give us relatively direct probes of the stellar populations, this radiation is subject to varying amounts of absorption before we view it. Some of this energy is absorbed by the gas directly in HII regions or transferred to the gas in photodissociation regions (PDRs) and reemitted as molecules, bands and atomic ionic and recombination lines, from wavelengths covering the UV to FIR and beyond. Some of the stellar energy is absorbed by the dust, revealed through extinction, and reradiated in MIR to submillimeter wavelengths as thermal emission. Therefore, models of the ISM in galaxies must consider these interdependent processes and be self-consistent. Our knowledge of the wavelength window from the MIR to the FIR has been limited by the low spatial and spectral resolution provided by IRAS, and has been rather sketchy when it comes to detailed studies of the ISM of individual galaxies. The Infrared Space Observatory (ISO) has been a recent turning point in this effort, providing high spectral and spatial resolution and unprecedented sensitivity in the MIR to FIR. We are incorporating our MIR and FIR observations in a study of the energy redistribution in starburst galaxies to understand the effects of the star formation on the surrounding gas and dust. Here we report the progress to date in our study of star forming low-metallicity dwarf galaxies, which, in the absence of major dynamical complications, allow us to ‘simplify’ model assumptions and the interpretation of observations.
## 2 Far-infrared observations: the \[CII\] cooling line
As an indirect probe of the star formation activity, we have obtained KAO (Kuiper Airborne Observatory) and ISO observations of the 158 $`\mu `$m $`{}_{}{}^{2}P_{3/2}^{}^2P_{1/2}`$ far infrared \[CII\] fine structure line emission in a sample of 15 dwarf galaxies with metallicities ranging from 0.1 to 0.5 solar. As the ionization potential of carbon is 11.3 eV, less than that of HI, photons escaping the HII regions, dissociate CO, and ionize carbon in the photodissociation regions (PDRs) on the surfaces of nearby molecular clouds exposed to the stellar UV radiation. The observed \[CII\] intensity can be traced back to the radiation source due to the fact that the UV photons heat the dust which emits thermal radiation in the MIR to submillimeter wavelengths. Energetic electrons, ejected from the dust through the photoelectric effect, heat the gas. The gas subsequently cools via emission from molecules and atomic fine structure lines, predominantly the 158 $`\mu `$m \[CII\] and the 63 $`\mu `$m \[OI\] transitions in PDRs. There has been a long history of development of PDR models which provide tools to differentiate physical properties, such as density (n), radiation field strength (G<sub>0</sub>) and filling-factors in galaxies (see review and references in ).
### 2.1 \[CII\] Survey of Dwarf Galaxies
The ratio of I\[CII\]/I(CO) is a useful measure of the PDR emission relative to the molecular core emission and is an indicator of the degree of star formation activity in galaxies. Active galaxies have a ratio of I\[CII\]/I(CO) $``$ 6300, which is 3 times greater than that observed in more quiescent galaxies . Our \[CII\] survey shows that for dwarf galaxies, this ratio ranges from 6000 to 70,000, which is up to 10 times greater than those for normal metallicity starburst galaxies (Figure 1) . We also observe an overall enhancement in the I\[CII\]/FIR ratios (where FIR is defined as the sum of the IRAS 60 and 100 $`\mu `$m bands) in these regions compared to those in normal metallicity galaxies, which was also noted in the LMC . The ratio of I\[CII\]/FIR is a direct measure of the fraction of UV energy reemerging in the \[CII\] cooling line, and is usually between 0.1% and 1% for normal metallicity galaxies , while we find up to 2% for dwarf galaxies. Observations of CO in dwarf galaxies have been very challenging and the glaring underabundance of observed CO in dwarf galaxies and relatively high FIR/CO luminosities have often been interpreted as unusually high star formation efficiency. While all these observational effects are a consequence of the lower metal abundance and decreased dust to gas ratio, we do not find an unambiguous direct correlation of the I\[CII\]/I(CO) and I\[CII\]/FIR ratios in our surveys with metallicity.
The reduced dust abundance in these environments allows the UV radiation to penetrate deeper, leaving a smaller CO core surrounded by a larger C<sup>+</sup>\- emitting region, thus enhancing the I\[CII\]/I(CO) ratios . Consequently, as the FUV flux travels further, the intensity becomes geometrically diluted, resulting in a lower beam-averaged FIR flux, accounting for the increased I\[CII\]/FIR ratios .
Using the results of recent PDR models that consider the effects of reduced metallicity , we can find solutions for the dwarf galaxies for clouds in our beam described by 2 different cases. One possible solution (case A) is for clouds with low A<sub>v</sub> ($``$3) and equal densities (n) in the CO and C<sup>+</sup>\- emitting regions with n ranging from 10<sup>3</sup> to 10<sup>4.5</sup> cm<sup>-3</sup> and low to moderate G<sub>o</sub> (normalized to the local interstellar radiation field intensity, 1.3x10<sup>-4</sup> erg s<sup>-1</sup> cm<sup>-2</sup> sr<sup>-1</sup>) ranging from 10<sup>1.5</sup> to 10<sup>3</sup>. Another possible solution (case B) is a higher A<sub>v</sub> ($``$10) with the density of the CO-emitting region (n<sub>CO</sub>) $`>`$ the density of the C<sup>+</sup>\- emitting region (n<sub>CII</sub>) which gives higher ranges of G<sub>o</sub> ($``$ 10<sup>2.5</sup> to 10<sup>3.5</sup>). We can put further constraints on these solutions through stellar population modeling. Based on our modeled SED for IIZw40, for example, case A is a solution (section 4.1). Arguments for molecular cloud stability point toward case B for the LMC . Decreasing the A<sub>v</sub> (case A) or increasing the n<sub>CO</sub> relative to n<sub>CII</sub> (case B) has the similar effect of reducing the CO-emitting core and increasing the C<sup>+</sup>\- emitting zone and increasing the CO-to-H<sub>2</sub> conversion factor . Based on \[CII\] observations in IC10, for example, we speculated that up to 100 times more H<sub>2</sub> may be ‘hidden’ in a C<sup>+</sup>-emitting regions compared to that deduced only from CO observations and using the Galactic CO-to-H<sub>2</sub> conversion factor . The presence of H<sub>2</sub> in the C<sup>+</sup>\- emitting region is due to the self-shielding of H<sub>2</sub> from UV photons or shielding by dust .
## 3 Mid-Infrared Observations
We are studying some of these galaxies in our \[CII\] survey with followup MIR observations. In Figure 2 we show ISOCAM spectra covering 5 to 17 $`\mu `$m for 3 galaxies from our \[CII\] survey, IIZw40, NGC 1140 and NGC 1569 along with that of the notoriously metal-poor SBS0335-052 . The spectra represent the total emission from the galaxies except in the case of the NGC 1569 spectra, which samples the region around the H$`\alpha `$ peak $`\mathrm{\#}`$2 (see ). As often seen in starburst galaxies (e.g. ), the MIR spectra are dominated by steeply rising continua longward of $``$ 10 $`\mu `$m, as evident in NGC 1569, IIZw40 and SBS0335-052 (Figure 2). Thermal emission from hot small grains with mean temperatures of the order of 100’s of K are responsible for the MIR continuum emission. The unidentified infrared bands (UIBs) at 6.2, 7.7, 8.6, 11.3 and 12.6 $`\mu `$m, are proposed to be due to aromatic hydrocarbon particles undergoing stochastic temperature fluctuations (i.e, PAHs ; coal grains ) and are observed to peak on the PDR zones around the HII regions but are destroyed deep within HII regions . While the UIBs are not obvious in the spectra of IIZw40 and SBS0335-052, and are only very weakly present NGC 1569, they can be distinguished in the spectrum of NGC 1140. Several ground state fine-structure nebular lines are present also in 3 of the spectra, the most prominent being 15.6 $`\mu `$m \[NeIII\] (energy potential $``$ 41 eV) and 10.5 $`\mu `$m \[SIV\] (energy potential $``$ 35 eV). Weaker, lower energy lines may also present, such as the 8.9 $`\mu `$m \[ArIII\] line and the \[NeII\] 12.8 $`\mu `$m line, which can be blended with the 12.6 $`\mu `$m UIB. All of these spectra look very different from one another and all differ significantly from those of normal metallicity starburst galaxies. Normal starburst galaxies show prominent UIBs, in contrast to AGNs, which are devoid of UIBs (e.g. ). When compared to spectra characteristic of PDRs and HII regions, ie, M17 , IIZw40 is remarkably similar to that of an HII region. In contrast, NGC 1140, which has a very flat continuum, yet very strong \[NeIII\] line, does have a more obvious contribution from PDR regions in its spectra. The MIR spectra of N66, the most prominent HII region in the SMC, also shows a scarcity of UIBs in the vicinity of the most massive central cluster , as does the low metallicity source NGC 5253 .
In some starburst galaxies, amorphous silicate is seen in absorption centered at 9 and 18 $`\mu `$m (e.g. , , ). We can fit the MIR region of the IIZw40 spectrum with a blackbody of 193 K and and an absorption equivalent to A<sub>v</sub> $``$ 4. We caution interpretation of the dust temperature we derive assuming a blackbody, since the dust emitting in the MIR is expected to be undergoing stochastic heating events, rather than being in thermal equilibrium with the radiation field. The amount of absorption in IIZw40 (A<sub>v</sub> $``$ 4) has yet to be confirmed. In SBS0335-052, a very low metallicity galaxy (1/40 solar), A<sub>v</sub> $``$ 20 deduced from the absorption in the ISOCAM MIR spectra (Figure 2) . The presence of a significant amount of dust in such a low metallicity galaxy is surprising, since star formation in SBS0335-052 began as recently as 100 Myr ago . Such high extinction implies that the current star formation rate, hidden by dust, can be underestimated by at least 50%
### 3.1 Effects of the starburst activity on the dwarf galaxy MIR spectra
As a consequence of the decreased dust abundance in dwarf galaxies, the ISM throughout the galaxies is effected globally by the hard radiation field of the massive stellar clusters. These galaxies contain evidence for Wolf-Rayet stars and super star clusters have been detected in NGC 1140 , NGC 1569 and SBS0335-052 . The harsh radiation field, which more easily permeates the ISM compared to normal metallicity environments, can destroy the UIB carriers, for example, over very extensive spatial areas. The effect of the pervasive radiation field can be witnessed in NGC 1569 (Figure 3). Photodissociation occurs on global scales. Violent activity is revealed by the H$`\alpha `$ distribution and the 15.8 $`\mu `$m \[NeIII\] emission, with giant streamers suspected to originate from the energetic winds of the super star clusters A & B, (shown in the figure as white stars). The UIBs, \[SIV\] and \[NeIII\] emission seem to avoid the super star clusters, which blow out much of the gas and dust on relatively short time scales. This effect is also seen in the CO , HI and the H$`\alpha `$ distribution. Likewise we see the destruction of the UIBs in the beam- averaged spectrum of the entire galaxy of IIZw40 and SBS0335-052 (due to our lack of spatial resolution we do not see the details within these galaxies in the MIR).
## 4 Spectral Energy Distribution
We compile broad-band data from the literature for IIZw40, NGC 1569 and NGC 1140, and together with our MIR data, construct stellar spectral energy distributions (SEDs). In doing so, we fit the observed optical and NIR data with stellar evolution models of PEGASE , taking into account the results of photoionization modeling of the MIR line emission using CLOUDY . This is an attempt to reconstruct the input stellar spectra consistent from the viewpoints of both the stellar evolution and photoionization.
### 4.1 Combined stellar evolution and photoionization model results
Using PEGASE with an instantaneous star formation rate, metallicity 0.2 solar and a Salpeter IMF (with upper and lower mass cut offs of .1 and 120 solar masses), we find solutions to observed broad band stellar light for various ages and ionization parameters. Diagnostic optical and NIR lines in the literature exist for all of these sources for a variety of apertures. The ISOCAM MIR observations also provide important diagnostic lines of neon, sulphur and argon, and has been recently addressed by others, including . For example, the \[NeIII\]/\[NeII\] ratio, is a measure of T<sub>eff</sub>, the hardness of the radiation field, and therefore traces the massive stellar population. For the dwarf galaxies, we find \[NeIII\]/\[NeII\] ratios in the range of 5 to 10 - much higher than those for normal metallicity galaxies ($``$1) . The extreme values of the \[NeIII\]/\[NeII\] ratios are due to effects of the low-metallicities of the systems: the T<sub>eff</sub> of the stars increases as the metallicity decreases for a specific stellar age. High ratios of \[NeIII\]/\[NeII\] and the prominent \[SIV\] in these spectra limit the age of the present star formation to $`<`$ 5 Myr. Beyond this age, the massive stars have died and the \[NeIII\]/\[NeII\] ratio drops dramatically. The high excitation 24.9 $`\mu `$m \[OIV\] line, covered by the ISO SWS data, is observed in some dwarf galaxies and has been proposed to be due to the presence of Wolf-Rayet stars . For NGC 1569, NGC 1140 and IIZw40, we construct composite stellar SEDs that require a 75% to 95 % mass fraction of an ’older’ population ranging in age from about 10 Myr to 30 Myr along with 5% to 30% of a very young population, $`<`$ 5 Myr. The broad band optical and NIR data alone reveal predominantly the older population in our apertures. Figure 4 shows an example of the resultant composite SED obtained for IIZw 40, and the extreme-ultraviolet (EUV) radiation which the young, massive stellar population traces. Observational evidence for the presence of Wolf-Rayet stars also indicates a very young stellar population .
## 5 Dust modeling
Having modeled the radiation field above, we next use the stellar spectra of IIZw40, NGC 1569 and NGC 1140 as input to a dust model to deduce the nature of the various dust components emitting in the MIR and the FIR. This is an important step since dust plays a major role in influencing the chemical and physical state of the ISM. We use the Désert et al. model , which calculates the IR emission from large silicate grains (BGs), very small amorphous carbon grains (VSGs), and stochastically-heated polycyclic aromatic hydrocarbons (PAHs), for various grain size distributions. This model is rather empirical in its approach and thus does not give an exact fit to the details of the observed spectrum. For example, the 8.6 $`\mu `$m UIB is not well-matched and no emission from bands at wavelengths longer than 11.3 $`\mu `$m are included. The model is currently in the process of modification using up-to-date laboratory-measured optical constants for a wide range of likely interstellar grain materials.
### 5.1 Dust in low-metallicity galaxies
In Figure 5 we show, as an example, the ISOCAM MIR spectrum and the IRAS data points for IIZw40 and NGC 1569, where we have plotted the emission from the PAH (dashed line), VSG (dotted line) and BG (dashed-dotted line) components. In these galaxies the MIR spectrum is clearly dominated by emission from VSGs with very little PAH emission. The BG component dominantes the overall dust emission with mass fractions ranging from 93% to 99% for the 3 galaxies, while the PAH mass fraction is relatively insignificant - 5 orders of magnitude lower. This model gives a PAH/VSG mass ratio for NGC 1569 and IIZw 40 of 2 to 3x10<sup>-4</sup> and 10 times this for NGC 1140. The Désert et al. model applied to the Galactic cirrus gives PAH/VSG mass ratio $`1`$. Thus, even compared to the VSG population, we find an insignificant mass fraction of PAHs, reflecting the fact that the PAHs are destroyed in the hard radiation fields in these galaxies. PAHs are thought to be the primary particles responsible for the photoelectric heating process and are incorporated in PDR models . Our preliminary results, while not statistically robust at this stage, suggest that even in the absence of PAHs, the photoelectric effect is efficient, as both IIZw40 and NGC 1569 are relatively prominent \[CII\] sources from our survey. On the contrary, in NGC 1140, where PAHs are more obvious in the MIR spectra (Figure 2), we do not detect \[CII\]. VSGs (sizes determined from model $``$40 to 300A), which are very abundant relative to the PAHs in NGC 1569 and IIZw40, and less so in NGC 1140, may therefore, be the more efficient sources of photoelectric gas heating in these environments, rather than PAHs. More detailed studies of these galaxies will be carried out using the analytical dust model of Városi and Dwek , which takes into account radiative transfer in a two-phase clumpy environment and considers various geometries.
## 6 Summary
Tracers of various components of the ISM show evidence of effects of the hard stellar radiation field in dwarf galaxies on the surrounding ISM due to the decreased dust abundace, allowing photoionization over large galactic scales to occur. From our survey of the 158 $`\mu `$m \[CII\] PDR cooling line in dwarf galaxies, we observe an increased penetration of the FUV radiation field which enhances the I\[CII\]/I(CO) emission in dwarf galaxies up to a factor of 10 more than in normal metallicity star burst galaxies. We also find a small enhancement in the I\[CII\]/FIR ratio in dwarf galaxies. Followup MIR ISOCAM spectroscopy provides details of ionic lines, UIBs and small hot small grain emission distribution in dwarf galaxies. The strong MIR \[NeIII\]/\[NeII\] ratios are signatures of the hard radiation fields and indicate the presence of young massive stellar populations in dwarf galaxies. Because of the increase in T<sub>eff</sub> in low metallicity environments, this ratio is enhanced at least 5 to 10 times more in dwarf galaxies than in normal metallicity galaxies. The penetrating radiation field also effects the dust components, destroying the UIBs in some dwarf galaxies on global scales, as is evident in the MIR spectra and in the dust modeling.
## 7 Acknowledgements
This work, still in progress, results from a series of observations from ISO and the KAO and includes a number of collaborators such as S. Colgan, N. Geis, M. Haas, D. Hollenbach, A. Jones, P. Maloney, A. Poglitsch, D. Ragaigne, B. Smith and M. Wolfire. I have benefited from invaluable discussions with E. Dwek and A. Jones on dust modeling. I thank W. Waller for his H$`\alpha `$ image of NGC 1569. |
warning/0002/hep-th0002034.html | ar5iv | text | # Exact Renormalization Group Equations. An Introductory Review.
## 1 Introduction
The formal discussion of consequences of the renormalization group works best if one has a differential form of the renormalization group transformation. Also, a differential form is useful for the investigation of properties of the $`\epsilon `$ expansion to all orders (…) A longer range possibility is that one will be able to develop approximate forms of the transformation which can be integrated numerically; if so one might be able to solve problems which cannot be solved any other way.
By “exact renormalization group equation (ERGE)”, we mean the continuous (i.e. not discrete) realization of the Wilson renormalization group (RG) transformation of the action in which no approximation is made and also no expansion is involved with respect to some small parameter of the action <sup>1</sup><sup>1</sup>1On the meaning of the word “exact”, see a discussion following the talk by Halperin in .. Its formulation — under a differential form — is known since the early seventies . However, due to its complexity (an integro-differential equation), its study calls for the use of approximation (and/or truncation) methods. For a long time it was natural to use a perturbative approach (based on the existence of a small parameter like the famous $`\epsilon `$-expansion for example). But, the standard perturbative field theory (e.g. see ) turned out to be more efficient and, in addition the defenders of the nonperturbative approach have turned towards the discrete formulation of the RG due to the problem of the ”stiff” differential equation (see a discussion following a talk given by Wegner ). This is why it is only since the middle of the eighties that substantial studies have been carried out via:
* the truncation procedures in the scaling field method (extended studies of )
* the explicit consideration of the local potential approximation and of the derivative expansion
* an appealing use, for field theoreticians, of the ERGE
In the nineties there has been a rapid growth of studies in all directions, accounting for scalar (or vector) fields, spinor, gauge fields, finite temperature, supersymmetry, gravity, etc…
In this paper we report on progress in the handling of the ERGE. Due to the abundance of the literature on the subject and because this is an introductory report, we have considered in detail the various versions of the ERGE only in the scalar (or vector) case<sup>2</sup><sup>2</sup>2However, because of the review by Wetterich and collaborators in this volume, we have not reported on their work as it deserves. The reader is invited to refer to their review and also to from which he could realize the rich variety of calculations that may be done within the nonperturbative framework of the ERGE.. This critical review must also be seen as an incitement to look at the original papers of which we give a list as complete as possible.
Let us mention that the ERGE is almost ignored in most of the textbooks on the renormalization group except notably in (see also ).
## 2 Exact Renormalization Group Equations
### 2.1 Introduction
There are four representations of the ERGE: the functional differential equation, the functional integral, the infinite set of partial differential equations for the couplings $`u_n(𝐩_1,\mathrm{},𝐩_n;t)`$ (eq. 11.19 of ) which are popularly known because they introduce the famous “beta” functions $`\beta _n\left(\left\{u_n\right\}\right)`$:
$$\frac{\text{d}u_n}{\text{d}t}=\beta _n\left(\left\{u_n\right\}\right)$$
and the infinite hierarchy of the ordinary differential equations for the scaling fields $`\mu _i(t)`$ (the RG-scale parameter $`t`$ is defined in section 2.2 and the scaling fields $`\mu _i(t)`$ in section 2.2.4).
In this review we shall only consider the functional differential representation of the ERGE.
There is not a unique form of the ERGE, each form of the equation is characterized by the way the momentum cutoff $`\mathrm{\Lambda }`$ is introduced. (In perturbative RG, this kind of dependence is known as the regularization-scheme dependence or “scheme dependence” in short .) The important point is that the various forms of the ERGE embody a unique physical content in the sense that they all preserve the same physics at large distances and, via the recourse to a process of limit, yield the same physics at small distances (continuum limits). The object of this part is to present the main equations used in the literature and connections between their formally (but not necessarily practically) equivalent forms. We do not derive them in detail here since one may find the derivations in several articles or reviews (that we indicate below).
Although the (Wilson) renormalization group theory owes much to the statistical physics as recently stressed by M. E. Fisher we adopt here the notations and the language of field theory.
Before considering explicitly the various forms of the ERGE (in sections 2.4-2.6), we find it essential to fix the notations and to remind some fundamental aspects of the RG.
### 2.2 Notations, reminders and useful definitions
We denote by $`\mathrm{\Lambda }`$ the momentum cutoff and the RG-“time” $`t`$ is defined by $`\frac{\mathrm{\Lambda }}{\mathrm{\Lambda }_0}=`$e<sup>-t</sup> in which $`\mathrm{\Lambda }_0`$ stands for some initial value of $`\mathrm{\Lambda }.`$
We consider a scalar field $`\varphi (𝐱)`$ with $`𝐱`$ the coordinate vector in an Euclidean space of dimension $`d`$. The Fourier transformation of $`\varphi (𝐱)`$ is defined as:
$$\varphi (𝐱)=_p\varphi _p\text{e}^{i𝐩x}$$
in which
$$_p\frac{d^dp}{(2\pi )^d}$$
(1)
and $`\varphi _p`$ stands for the function $`\varphi (𝐩)`$ where $`𝐩`$ is the momentum vector (wave vector).
The norms of $`𝐱`$ and $`𝐩`$ are noted respectively $`x`$ and $`p`$ ($`x\sqrt{𝐱𝐱}`$). However, when no confusion may arise we shall denote the vectors $`𝐱`$ and $`𝐩`$ by simply $`x`$ and $`p`$ as in (1) for example. Sometimes the letters $`k`$ and $`q`$ (or $`𝐤`$ and $`𝐪`$) will refer also to momentum variables.
It is useful to define $`K_d`$ as the surface of the $`d`$-dimensional unit sphere divided by $`\left(2\pi \right)^d`$, i.e.:
$$K_d=\frac{2\pi ^{\frac{d}{2}}}{\left(2\pi \right)^d\mathrm{\Gamma }\left(\frac{d}{2}\right)}$$
(2)
We shall also consider the case where the field has $`N`$ components $`\varphi =(\varphi _1,\mathrm{},\varphi _N)`$ that we shall also denote generically by $`\varphi _\alpha `$.
The action $`S[\varphi ]`$ (the Hamiltonian divided by $`k_\text{B}T`$ for statistical physics) is a general semi-local functional of $`\varphi `$. Semi-local means that, when it is expanded<sup>3</sup><sup>3</sup>3We do not need to assume this expansion as the unique form of $`S[\varphi ]`$ in general. in powers of $`\varphi `$, $`S[\varphi ]`$ involves only powers of $`\varphi (𝐱)`$ and of its derivatives with respect to $`x^\mu `$ (that we denote $`^\mu \varphi `$ or even $`\varphi `$ instead of $`\varphi /x^\mu `$). This characteristics is better expressed in the momentum-space (or wave-vector-space). So we write:
$$S[\varphi ]=\underset{n}{\overset{\mathrm{}}{}}_{p_1\mathrm{}p_n}u_n(𝐩_1,\mathrm{},𝐩_n)\varphi _{p_1}\mathrm{}\varphi _{p_n}\widehat{\delta }\left(𝐩_1+\mathrm{}+𝐩_n\right)$$
(3)
in which $`\widehat{\delta }\left(𝐩\right)\left(2\pi \right)^d\delta ^d\left(𝐩\right)`$ is the $`d`$-dimensional delta-function:
$$\delta ^d\left(𝐱\right)=_p\text{e}^{i𝐩x}$$
Notice that the O(1) symmetry $`\varphi \varphi `$, also called the $`Z_2`$ symmetry, is not assumed neither here nor in the following sections except when it is explicitly mentioned.
The $`u_n`$’s are invariant under permutations of their arguments.
For the functional derivative with respect to $`\varphi `$, we have the relation:
$$\frac{\delta }{\delta \varphi _p}=\text{d}^dx\text{e}^{i𝐩x}\frac{\delta }{\delta \varphi (𝐱)}$$
so that in performing the functional derivative with respect to $`\varphi _p`$ we get rid of the $`\pi `$-factors involved in the definition (1), e.g..:
$`{\displaystyle \frac{\delta }{\delta \varphi _p}}`$ $`\left[{\displaystyle _{p_1\mathrm{}p_n}}u_n(𝐩_1,\mathrm{},𝐩_n)\varphi _{p_1}\mathrm{}\varphi _{p_n}\widehat{\delta }\left(𝐩_1+\mathrm{}+𝐩_n\right)\right]=`$
$`n{\displaystyle _{p_1\mathrm{}p_{n1}}}u_n(𝐩_1,\mathrm{},𝐩_{n1},𝐩)\varphi _{p_1}\mathrm{}\varphi _{p_{n1}}\widehat{\delta }\left(𝐩_1+\mathrm{}+𝐩_{n1}+𝐩\right)`$
in order to lighten the notations we sometimes will write $`\frac{\delta }{\delta \varphi }`$ instead of $`\frac{\delta }{\delta \varphi (𝐱)}`$ when no confusion may arise.
Let us also introduce:
* The generating functional $`Z[J]`$ of Green’s functions:
$$Z[J]=𝒵^1𝒟\varphi \mathrm{exp}\left\{S[\varphi ]+J\varphi \right\}$$
(4)
in which
$$J\varphi d^dxJ(𝐱)\varphi (𝐱)$$
$`J(𝐱)`$ is an external source, and $`𝒵`$ is a normalization such that $`Z[0]=1`$. Indeed $`𝒵`$ is the partition function:
$$𝒵=𝒟\varphi \mathrm{exp}\left\{S[\varphi ]\right\}$$
(5)
* The generating functional $`W[J]`$, of connected Green functions, is related to $`Z[J]`$ as follows:
$$W[J]=\mathrm{ln}\left(Z[J]\right)$$
(6)
Notice that if one defines $`𝒲`$
$$\text{e}^𝒲=𝒵$$
then $`𝒲`$ is minus the free energy.
* The Legendre transformation which defines the generating functional $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ of the one-particle-irreductible (1PI) Green functions (or simply vertex functions):
$`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]+W\left[J\right]J\mathrm{\Phi }`$ $`=`$ $`0`$
$`{\displaystyle \frac{\delta \mathrm{\Gamma }\left[\mathrm{\Phi }\right]}{\delta \mathrm{\Phi }(x)}}|_J`$ $`=`$ $`J(x)`$ (7)
in which we have introduced the notation $`\mathrm{\Phi }`$ to make a distinction between this and the (dummy) field variable $`\varphi `$ in (4). In the following we shall not necessarily make this distinction.
#### 2.2.1 Dimensions
First let us define our conventions relative to the usual (i.e. engineering) dimensions.
In the following, we refer to a system of units in which the dimension of a length scale $`L`$ is $`1`$:
$$\left[L\right]=1$$
and a momentum scale like $`\mathrm{\Lambda }`$ has the dimension $`1`$:
$$\left[\mathrm{\Lambda }\right]=1$$
The usual classical dimension of the field (in momentum unit):
$$d_\varphi ^c=\left[\varphi \right]=\frac{1}{2}\left(d2\right)$$
(8)
is obtained by imposing that the coefficient of the kinetic term $``$d$`{}_{}{}^{d}x\left(\varphi (x)\right)^2`$ in $`S`$ is dimensionless \[it is usually set to $`\frac{1}{2}`$\].
The dimension of the field is not always given by (8). Indeed one knows that the field may have an anomalous dimension:
$$d_\varphi ^a=\frac{1}{2}\left(d2+\eta \right)$$
(9)
with $`\eta `$ a non-zero constant defined with respect to a non-trivial fixed point.
In RG theory, the dimension of the field depends on the fixed point in the vicinity of which the (field) theory is considered. Hence we introduce an adjustable dimension of the field, $`d_\varphi `$ , which controls the scaling transformation of the field:
$$\varphi \left(s𝐱\right)=s^{d_\varphi }\varphi \left(𝐱\right)$$
(10)
For the Fourier transformation, we have:
$$\varphi \left(s𝐩\right)=s^{d_\varphi d}\varphi \left(𝐩\right)$$
(11)
Since the dimension of any dimensioned (in the classical meaning for $`\varphi `$) quantity is expressed in term of a momentum scale, we use $`\mathrm{\Lambda }`$ to reduce all dimensioned quantities into dimensionless quantities. In the following we deal with dimensionless quantities and, in particular, the notation $`p`$ will refer to a dimensionless momentum variable. However sometimes, for the sake of clarity, we will need to reintroduce the explicit $`\mathrm{\Lambda }`$-dependence, e.g. via the ratio $`p/\mathrm{\Lambda }`$.
It is also useful to notice that, with a dimensionless $`p`$, the following derivatives are equivalent (the derivative is taken at constant dimensioned momentum):
$$\frac{}{t}=\mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}=p\frac{}{p}$$
#### 2.2.2 Transformations of the field variable
In order to discuss invariances in RG theory, it is useful to consider a general transformation of the field which leaves invariant the partition function. Following , we replace $`\varphi _p`$ by $`\varphi _p^{}`$ such that
$$\varphi _p^{}=\varphi _p+\sigma \mathrm{\Psi }_p[\varphi ]$$
(12)
where $`\sigma `$ is infinitesimally small and $`\mathrm{\Psi }_p`$ a function which may depend on all Fourier components of $`\varphi `$. Then one has
$$S[\varphi ^{}]=S[\varphi ]+\sigma _p\mathrm{\Psi }_p[\varphi ]\frac{\delta S[\varphi ]}{\delta \varphi _p}$$
Moreover we have:
$$𝒟\varphi ^{}=𝒟\varphi \frac{\left\{\varphi ^{}\right\}}{\left\{\varphi \right\}}=𝒟\varphi \left(1+\sigma _p\frac{\delta \mathrm{\Psi }_p[\varphi ]}{\delta \varphi _p}\right)$$
The transformation must leave the partition function $`𝒵`$ \[eq. (5)\] invariant. Therefore one obtains
$`𝒵`$ $`=`$ $`{\displaystyle 𝒟\varphi ^{}\mathrm{exp}\left\{S[\varphi ^{}]\right\}}`$
$`=`$ $`{\displaystyle 𝒟\varphi \mathrm{exp}\left\{S[\varphi ]\sigma 𝒢_{\text{tra}}\left\{\mathrm{\Psi }\right\}S[\varphi ]\right\}}`$
with
$$𝒢_{\text{tra}}\left\{\mathrm{\Psi }\right\}S[\varphi ]=_p\left(\mathrm{\Psi }_p\frac{\delta S}{\delta \varphi _p}\frac{\delta \mathrm{\Psi }_p}{\delta \varphi _p}\right)$$
(13)
which indicates how the action transforms under the infinitesimal change (12):
$$\frac{\text{d}S}{\text{d}\sigma }=𝒢_{\text{tra}}\left\{\mathrm{\Psi }\right\}S$$
In the case of $`N`$ components, the expression (13) generalizes obviously:
$$𝒢_{\text{tra}}\left\{𝚿\right\}S[\varphi ]=\underset{\alpha =1}{\overset{N}{}}_p\left(\mathrm{\Psi }_p^\alpha \frac{\delta S}{\delta \varphi _p^\alpha }\frac{\delta \mathrm{\Psi }_p^\alpha }{\delta \varphi _p^\alpha }\right)$$
(14)
#### 2.2.3 Rescaling
We consider an infinitesimal change of (momentum) scale:
$$pp^{}=sp=(1+\sigma )p$$
(15)
with $`\sigma `$ infinitesimally small. Introducing the rescaling operator of , the consequence on $`S[\varphi ]`$ is written as
$$SS^{}=S+\sigma 𝒢_{\text{dil}}S$$
(16)
and thus:
$$\frac{\text{d}S}{\text{d}\sigma }=𝒢_{\text{dil}}S$$
(17)
Considering $`S`$ as given by (3), then $`\mathrm{\Delta }S=\sigma 𝒢_{\text{dil}}S`$ may be expressed by gathering the changes induced by (15) on the various factors in the sum, namely:
the differential volume $`_{i=1}^n`$d$`{}_{}{}^{d}p_{i}^{}=(1+\sigma )^{nd}_{i=1}^nd^dp_i^{}`$ induces a change $`\mathrm{\Delta }S_1`$ which may be written as
$$\mathrm{\Delta }S_1=\sigma \left(d_p\varphi _p\frac{\delta }{\delta \varphi _p}\right)S$$
the couplings $`u_n(𝐩_1,\mathrm{},𝐩_n)=u_n(s^1𝐩_1^{},\mathrm{},s^1𝐩_n^{})`$ induce a change $`\mathrm{\Delta }S_2`$ which may be written as
$$\mathrm{\Delta }S_2=\sigma \left(_p\varphi _p𝐩_p^{}\frac{\delta }{\delta \varphi _p}\right)S$$
where the prime on the derivative symbol ($`_p^{}`$) indicates that the momentum derivative does not act on the delta-functions.
the delta-functions $`\widehat{\delta }\left(𝐩_1+\mathrm{}+𝐩_n\right)=\widehat{\delta }\left(s^1𝐩_1^{}+\mathrm{}+s^1𝐩_n^{}\right)`$ induce a change $`\mathrm{\Delta }S_3`$ which may be written as
$$\mathrm{\Delta }S_3=\sigma \left(dS\right)$$
the field itself $`\varphi \left(𝐩\right)=\varphi \left(s^1𝐩^{}\right)`$ induces a change $`\mathrm{\Delta }S_4`$ which, according to (11), is dictated by
$$\varphi \left(s^1𝐩^{}\right)=(1+\sigma )^{dd_\varphi }\varphi \left(𝐩\right)$$
Hence, to the first order in $`\sigma `$, we have
$$\mathrm{\Delta }S_4=\sigma \left[\left(dd_\varphi \right)_p\varphi _p\frac{\delta }{\delta \varphi _p}\right]S$$
Summing the four contributions, $`\mathrm{\Delta }S=_{i=1}^4\mathrm{\Delta }S_i`$ we obtain:
$$\mathrm{\Delta }S=\sigma \left(dS_p\varphi _p𝐩_p^{}\frac{\delta }{\delta \varphi _p}d_\varphi _p\varphi _p\frac{\delta }{\delta \varphi _p}\right)S$$
This expression may be further simplified by allowing the momentum derivative $`_p`$ to act also on the delta-functions (this eliminates the prime and absorbs the term $`dS`$). We thus may write
$$\mathrm{\Delta }S=\sigma \left(_p\varphi _p𝐩_p\frac{\delta }{\delta \varphi _p}+d_\varphi _p\varphi _p\frac{\delta }{\delta \varphi _p}\right)S$$
hence:
$$𝒢_{\text{dil}}S=\left(_p\varphi _p𝐩_p\frac{\delta }{\delta \varphi _p}+d_\varphi _p\varphi _p\frac{\delta }{\delta \varphi _p}\right)S$$
(18)
As in the case of $`𝒢_{\text{tra}}`$, the generalization to $`N`$ components is obvious \[see eq. (14)\].
The writing of the action of $`𝒢_{\text{dil}}`$ \[eq. (18)\] may take on two other forms in the literature:
1. due to the possible integration by parts of the first term:
$$𝒢_{\text{dil}}S=\left(_p\left(𝐩_p\varphi _p\right)\frac{\delta }{\delta \varphi _p}+\left(dd_\varphi \right)_p\varphi _p\frac{\delta }{\delta \varphi _p}\right)S$$
(19)
2. if one explicitly performs the derivative with respect to $`𝐩`$ acting on the $`\delta `$-functions:
$$𝒢_{\text{dil}}S=dS\left(_p\varphi _p𝐩_p^{}\frac{\delta }{\delta \varphi _p}+d_\varphi _p\varphi _p\frac{\delta }{\delta \varphi _p}\right)S$$
(20)
in which now $`_p^{}`$ does not act on the $`\delta `$-functions.
An other expression of the operator $`𝒢_{\text{dil}}`$ may be found in the literature , it is:
$$𝒢_{\text{dil}}=d\mathrm{\Delta }_{}d_\varphi \mathrm{\Delta }_\varphi $$
(21)
where $`\mathrm{\Delta }_\varphi =\varphi .\frac{\delta }{\delta \varphi }`$ is the ‘phi-ness’ counting operator: it counts the number of occurrences of the field $`\varphi `$ in a given vertex and $`\mathrm{\Delta }_{}`$ may be expressed as
$$\mathrm{\Delta }_{}=d+_p\varphi _p𝐩_p\frac{\delta }{\delta \varphi _p}$$
(22)
i.e. the momentum scale counting operator $`+d`$. Operating on a given vertex it counts the total number of derivatives acting on the fields $`\varphi `$ .
Notice that we have not introduced the anomalous dimension $`\eta `$ of the field. This is because it naturally arises at the level of searching for a fixed point of the ERGE. As we indicate in the following section, the introduction of $`\eta `$ is related to an invariance.
#### 2.2.4 Linearized RG theory
Following Wegner , we write the ERGE under the following formal form
$$\frac{\text{d}S}{\text{d}t}=𝒢S$$
Near a fixed point $`S^{}`$ (such that $`𝒢S^{}=0`$) we have:
$$\frac{\text{d}\left(S^{}+\mathrm{\Delta }S\right)}{\text{d}t}=\mathrm{\Delta }S+𝒬\mathrm{\Delta }S$$
in which the RG operator has been separated into a linear $``$ and a quadratic $`𝒬`$ parts.
The eigenvalue equation:
$$𝒪_i^{}=\lambda _i𝒪_i^{}$$
defines scaling exponents $`\lambda _i`$ and a set (assumed to be complete) of eigenoperators $`𝒪_i^{}`$. Hence we have for any $`S(t)`$:
$$S(t)=S^{}+\underset{i}{}\mu _i(t)𝒪_i^{}$$
In which $`\mu _i`$ are the “scaling fields” which in the linear approximation satisfy:
$$\frac{\text{d}\mu _i(t)}{\text{d}t}=\lambda _i\mu _i(t)$$
Which yields:
$$S(t)=S^{}+\underset{i}{}\mu _i(0)t^{\lambda _i}𝒪_i^{}$$
(23)
##### Scaling operators
There are three kinds of operators associated to well defined eigenvalues and called “scaling operators” in which are classified as follows<sup>4</sup><sup>4</sup>4In perturbative field theory, where the implicit fixed point is Gaussian, the scaling fields are called superrenormalizable, nonrenormalizable and strictly renormalizable respectively.:
* $`\lambda _i>0`$, the associated scaling field $`\mu _i`$ (or the operator $`𝒪_i^{}`$) is relevant because it brings the action away from the fixed point.
* $`\lambda _i<0`$, the associated scaling field $`\mu _i`$ is irrelevant because it decays to zero when $`t\mathrm{}`$ and $`S(t)`$ finally approaches $`S^{}`$.
* $`\lambda _i=0,`$ the associated scaling field $`\mu _i`$ is marginal and $`S^{}+\mu _i𝒪_i^{}`$ is a fixed point for any $`\mu _i`$. This latter property may be destroyed beyond the linear order<sup>5</sup><sup>5</sup>5This is the case of the renormalized $`\varphi ^4`$-coupling constant for $`d=4`$ with respect to the Gaussian fixed point: it is marginal in the linear approximation and irrelevant beyond. It is marginally irrelevant..
In critical phenomena, the relevant scaling fields alone are responsible for the scaling form of the physical quantities: e.g., in its scaling form, the free energy depends only on the scaling fields . The irrelevant scaling fields induce corrections to the scaling form .
To be specific, the positive eigenvalue of a critical fixed point (once unstable), say $`\lambda _1`$, is the inverse of the correlation length critical exponent $`\nu `$ and the less negative eigenvalue is equal to (minus) the subcritical exponent $`\omega `$.
In modern field theory, only the relevant (or marginally relevant) scaling fields are of interest in the continuum limit : they correspond to the renormalized couplings (or masses) of field theory the continuum limit of which being defined “at” the considered fixed point (see section 2.10).
##### Redundant operators and reparametrization invariance
In addition to scaling operators, there are redundant operators . They come out due to invariances of the RG (see also ). Thus they can be expressed in the form $`𝒢_{\text{tra}}\left\{\mathrm{\Phi }\right\}S^{}`$ and the associated exponents $`\lambda _i`$ (which, in general, have nonuniversal values) are spurious since the free energy does not depend on the corresponding redundant fields $`\mu _i`$ (by construction of the transformation generator $`𝒢_{\text{tra}}`$ which leaves the partition function invariant, see section 2.2.2).
Although unphysical, the redundant fields cannot be neglected. For example, a well known redundant operator is<sup>6</sup><sup>6</sup>6A clear explanation of this may be found in p.101–102. $`\frac{\delta S^{}}{\delta \varphi _0}`$ which may be written under the form $`𝒢_{\text{tra}}\left\{\mathrm{\Phi }\right\}S^{}`$ with $`\mathrm{\Phi }_q=\widehat{\delta }(q)`$. and, most often, has the eigenvalue $`\lambda =\frac{1}{2}\left(d2+\eta \right)`$. Since $`\lambda >0`$ for $`d=3`$, this operator is relevant with respect to the Wilson-Fisher (i.e. Ising-like for $`d=3`$) fixed point although it is not physical. Indeed, as pointed out by Hubbard and Schofield in , the fixed point becomes unstable in presence of a $`\varphi ^3(x)`$ term which, however, may be eliminated by the substitution $`\varphi _0\varphi _0+\mu `$, which is controlled by the operator $`\frac{\delta S^{}}{\delta \varphi _0}`$ since:
$$S^{}\left(\varphi _0+\mu \right)=S^{}\left(\varphi _0\right)+\mu \frac{\delta S^{}}{\delta \varphi _0}+O\left(\mu ^2\right)$$
This redundant operator is not really annoying because it is sufficient to consider actions that are even functional of $`\varphi `$ ($`Z_2`$-symmetric) to get rid of $`\frac{\delta S^{}}{\delta \varphi _0}`$.
Less obvious and more interesting for field theory is the following redundant operator:
$$𝒪_1=_q\left[\frac{\delta ^2S}{\delta \varphi _q\delta \varphi _q}\frac{\delta S}{\delta \varphi _q}\frac{\delta S}{\delta \varphi _q}+\varphi _q\frac{\delta S}{\delta \varphi _q}\right]$$
(24)
which has been studied in detail by Riedel et al . When the RG transforms the field variable linearly (as in the present review), $`𝒪_1`$ has, once and for all, the eigenvalue $`\lambda _1=0`$; it is absolutely marginal .
$`𝒪_1`$ is redundant because it may be written under the form $`𝒢_{\text{tra}}\left\{\mathrm{\Phi }\right\}S`$ with:
$$\mathrm{\Phi }_q=\varphi _q\frac{\delta S}{\delta \varphi _q}$$
The redundant character of $`𝒪_1`$ is related to the invariance of the RG transformation under a change of the overall normalization of $`\varphi `$ . This invariance is also called the “reparametrization invariance” .
The most general realization of this symmetry is not linear, this explains why $`𝒪_1`$ is so complicated (otherwise, in the case of a linear realization of the invariance, $`𝒪_1`$ would reduce to simply $`_q\varphi _q\frac{\delta S}{\delta \varphi _q}`$).
As consequences of the reparametrization invariance :
* a line of equivalent fixed points exists which is parametrized by the normalization of the field,
* a field-rescaling parameter that enters in the ERGE must be properly adjusted in order for the RG transformation to have a fixed point (the exponent $`\eta `$ takes on a specific value).
We illustrate these two aspects with the Gaussian fixed point in section 2.7 after having written down the ERGE.
* Due to the complexity of $`𝒪_1`$, truncations of the ERGE (in particular the derivative expansion, see section 4.1) may easily violate the reparametrization invariance , in which case the line of equivalent fixed points becomes a line of inequivalent fixed points yielding different nonuniversal values for the exponent $`\eta `$. However the search for a vestige of the invariance may be used to determine the best approximation for $`\eta `$ . In the case where the invariance is manifestly linearly realized and momentum independent (as when a regularization with a sharp cutoff is utilized for example, see below), then the derivative expansion may preserve the invariance and as a consequence, $`\eta `$ is uniquely defined (see section 4.1).
To understand why $`\eta `$ must take on a specific value, it is helpful to think of a linear eigenvalue problem. “The latter may have apparently a solution for each arbitrary eigenvalue, but the fact that we can choose the normalization of the eigenvector at will over-determines the system, making that only a discrete set of eigenvalues are allowed.” (see section 4.1)
### 2.3 Principles of derivation of the ERGE
The Wilson RG procedure is carried out in two steps (see also for example):
1. an integration of the fluctuations $`\varphi (𝐩)`$ over the range e$`{}_{}{}^{t}<\left|p\right|1`$ which leaves the partition function (5) invariant,
2. a change of the length scale by a factor e<sup>-t</sup> in all linear dimensions to restore the original scale $`\mathrm{\Lambda }`$ of the system, i. e. $`𝐩p^{}=`$e$`{}_{}{}^{t}𝐩`$
For infinitesimal value of $`t`$, step 2 corresponds to a change in the effective action \[see eqs (15, 16 and 17 with $`\sigma =t`$)\]
$$SS^{}=S+t𝒢_{\text{dil}}S$$
(25)
inducing a contribution to
$$\dot{S}\frac{S}{t}$$
(26)
which is equal to $`𝒢_{\text{dil}}S`$.
The step of reducing the number of degrees of freedom (step 1) is the main part of the RG theory. It is sometimes called “coarse grain decimation” by reference to a discrete realization of the RG transformation or Kadanoff’s transformation , it is also called sometimes “blocking” or “coarsening”. It carries its own arbitrariness due to the vague notion of “block”, i.e. in the present review the unprecised way of separating the high from the low momentum frequencies. Step 1 may be roughly introduced as follows.
We assume that the partition function may be symbolically written as
$$𝒵=\underset{p1}{}𝒟\varphi _p\mathrm{exp}\left\{S[\varphi ]\right\}$$
(27)
then after performing the integrations of step 1, we have:
$$𝒵=\underset{pe^t}{}𝒟\varphi _p\mathrm{exp}\left\{S^{}[\varphi ]\right\}$$
with
$$\mathrm{exp}\left\{S^{}[\varphi ;t]\right\}=\underset{e^t<p1}{}𝒟\varphi _p\mathrm{exp}\left\{S[\varphi ]\right\}$$
(28)
$`S^{}[\varphi ;t]`$ is named the Wilson effective action. By considering an infinitesimal value of $`t`$, one obtains an evolution equation for $`S`$ under a differential form, i.e. an explicit expression for $`\dot{S}`$.
As indicated by Wegner , the infinitesimal “blocking” transformation of $`S`$ (step 1) may sometimes<sup>7</sup><sup>7</sup>7When the cutoff is smooth. be expressed as a transformation of the field of the form introduced in section 2.2.2. Hence the general expression of the ERGE may formally be written as follows :
$$\dot{S}=𝒢_{\text{dil}}S+𝒢_{\text{tra}}\left\{𝚿\right\}S$$
in which $`𝚿`$ has different expressions depending on the way one introduces the cutoff $`\mathrm{\Lambda }`$. For example, in the case of the Wilson ERGE \[see eq. (30) below\], $`𝚿`$ has the following form :
$$\mathrm{\Psi }_p=\left(c+2p^2\right)\left(\varphi _p\frac{\delta S}{\delta \varphi _p}\right)$$
As it is introduced just above in (27-28), the cutoff is said sharp or hard<sup>8</sup><sup>8</sup>8It corresponds to a well defined boundary between low and high momentum frequencies, to be opposed to a smooth cutoff which corresponds to a blurred boundary.. It is known that a sharp boundary in momentum space introduces non-local interactions in position space which one would like to avoid. Nevertheless, a differential ERGE has been derived which has been used several times with success under an approximate form. Indeed, in the leading approximation of the derivative expansion (local potential approximation), most of the differences between a sharp and smooth cutoff disappear, and moreover, as stressed by Morris , the difficulties induced by the sharp cutoff may be circumvented by considering the Legendre transformation (7) (see section 2.6).
### 2.4 Wegner-Houghton’s sharp cutoff version of the ERGE
The equation has been derived in , one may also find an interesting detailed presentation in , it reads:
$$\dot{S}=\underset{t0}{lim}\frac{1}{2t}\left[_p^{}\mathrm{ln}\left(\frac{\delta ^2S}{\delta \varphi _p\delta \varphi _p}\right)_p^{}\frac{\delta S}{\delta \varphi _p}\frac{\delta S}{\delta \varphi _p}\left(\frac{\delta ^2S}{\delta \varphi _p\delta \varphi _p}\right)^1\right]+𝒢_{\text{dil}}S+\text{const}$$
(29)
in which we use (26) and the prime on the integral symbol indicates that the momenta are restricted to the shell $`e^t<|p|1`$, and $`𝒢_{\text{dil}}S`$ corresponds to any of the eqs (1820) with $`d_\varphi (t)`$ set to a constant in .
The explicit terms correspond to the step 1 (decimation or coarsening) of section 2.3, while $`𝒢_{\text{dil}}S`$ refers to step 2 (rescaling) as indicated in section 2.2.3. The additive constant may be neglected in field theory \[due to the normalization of (4)\].
### 2.5 Smooth cutoff versions of the ERGE
#### 2.5.1 Wilson’s incomplete integration
The first expression of the exact renormalization group equation under a differential form has been presented as far back as 1970 before publication in the famous Wilson and Kogut review (see chapt. 11). The step 1 (decimation) of this version (referred to below as the Wilson ERGE) consists in an “incomplete” integration in which large momenta are more completely integrated than small momenta (see chapt. 11 of and also p. 70 for the details).
The Wilson RG equation in our notations reads (with the change $`S`$ compared to ):
$$\dot{S}=𝒢_{\text{dil}}S+_p\left(c+2p^2\right)\left(\frac{\delta ^2S}{\delta \varphi _p\delta \varphi _p}\frac{\delta S}{\delta \varphi _p}\frac{\delta S}{\delta \varphi _p}+\varphi _p\frac{\delta S}{\delta \varphi _p}\right)$$
(30)
Some short comments relative to (30):
The term $`𝒢_{\text{dil}}S`$ (which comes out of the rescaling step 2) is given by one of the eqs. (18-21) but, in , the choice $`d_\varphi =d/2`$ has been made. The somewhat mysterious function $`c`$ (denoted $`\frac{\text{d}\rho }{\text{d}t}`$ in ) must be adjusted in such a way as to obtain a useful fixed point . This adjustment is related to the reparametrization invariance (see section 2.7 for an example). Notice that $`c`$ is precisely introduced in (30) in front of the operator $`𝒪_1`$ of eq. (24) which controls the change of normalization of the field (see section 2.2.4). Indeed, in the vicinity of the fixed point, we have :
$$c=1\frac{\eta }{2}$$
(31)
and most often $`c`$ is considered as a constant (as in for example). Notice that the unusual (for field theory) choice<sup>9</sup><sup>9</sup>9See ref. for a further explanation of this choice. $`d_\varphi =d/2`$ in , leads to the same anomalous dimension (9) at the fixed point.
#### 2.5.2 Polchinski’s equation
With a view to study field theory, Polchinski has derived his own smooth cutoff version of the ERGE (see also section 2.10.2). A general ultraviolet (UV) cutoff function $`K(p^2/\mathrm{\Lambda }^2)`$ is introduced (we momentaneously restore the dimensions) with the property that it vanishes rapidly when $`p>\mathrm{\Lambda }`$. (Several kinds of explicit functions $`K`$ may be chosen, the sharp cutoff would be introduced with the Heaviside step function $`K(x)=\mathrm{\Theta }(1x)`$.) The Euclidean action reads:
$$S[\varphi ]\frac{1}{2}_p\varphi _p\varphi _pp^2K^1(p^2/\mathrm{\Lambda }^2)+S_{\text{int}}[\varphi ]$$
(32)
Compared to , the “mass” term has been incorporated into $`S_{\text{int}}[\varphi ]`$, this does not corrupt in any way the eventual analysis of the massive theory because the RG naturally generates quadratic terms in $`\varphi `$ and the massive or massless character of the (field) theory is not defined at the level of eq. (32) but a posteriori in the process of defining the continuum limit (modern conception of the renormalization of field theory, see sections 2.10.1 and 3.4.1).
Polchinski’s ERGE is obtained from the requirement that the coarsening step (step 1 of section 2.3) leaves the generating functional $`Z[J]`$ \[eq. (4)\] invariant. The difficulty of dealing with an external source is circumvented by imposing that $`J(p)=0`$ for $`p>\mathrm{\Lambda }`$. As in , the derivation relies upon the writing of an (ad hoc) expression under the form of a complete derivative with respect to the field in such a way as to impose d$`Z[J]/`$d$`\mathrm{\Lambda }=0`$. The original form of Polchinski’s equation accounts only for the step 1 and reads (for more details on this derivation, see for example ):
$$\mathrm{\Lambda }\frac{\text{d}S_{\text{int}}}{\text{d}\mathrm{\Lambda }}=\frac{1}{2}_pp^2\mathrm{\Lambda }\frac{\text{d}K}{\text{d}\mathrm{\Lambda }}\left(\frac{\delta S_{\text{int}}}{\delta \varphi _p}\frac{\delta S_{\text{int}}}{\delta \varphi _p}\frac{\delta ^2S_{\text{int}}}{\delta \varphi _p\delta \varphi _p}\right)$$
(33)
then, considering the rescaling (step 2) and the complete action, the Polchinski ERGE is (see for example ):
$$\dot{S}=𝒢_{\text{dil}}S_pK^{}(p^2)\left(\frac{\delta ^2S}{\delta \varphi _p\delta \varphi _p}\frac{\delta S}{\delta \varphi _p}\frac{\delta S}{\delta \varphi _p}+\frac{2p^2}{K(p^2)}\varphi _p\frac{\delta S}{\delta \varphi _p}\right)$$
(34)
in which all quantities are dimensionless and $`K^{}(p^2)`$ stands for d$`K(p^2)/`$d$`p^2`$.
Let us mention that one easily arrives at eq. (33) using the observation that the two following functionals:
$$Z[J]=𝒟\varphi \mathrm{exp}\left\{\frac{1}{2}\varphi \mathrm{\Delta }^1\varphi S[\varphi ]+J\varphi \right\}$$
(35)
and
$$Z^{}[J]=𝒟\varphi \mathrm{exp}\left\{\frac{1}{2}\varphi _1\mathrm{\Delta }_1^1\varphi _1\frac{1}{2}\varphi _2\mathrm{\Delta }_2^1\varphi _2S[\varphi _1+\varphi _2]+J\left(\varphi _1+\varphi _2\right)\right\}$$
(36)
are equivalent (up to a multiplicative factor) provided that $`\mathrm{\Delta }=\mathrm{\Delta }_1+\mathrm{\Delta }_2`$ and $`\varphi =\varphi _1+\varphi _2`$ (see appendix<sup>10</sup><sup>10</sup>10Or appendix 24 of the first edition in 1989. 10 of and also ).
#### 2.5.3 Redundant free ERGE
In it is proposed to make the effective dimension of the field $`d_\varphi `$ \[defined in (10)\], which enters Polchinski’s ERGE (34) via the rescaling part $`𝒢_{\text{dil}}S`$ \[see eqs (1821)\], depend on the momentum $`𝐩`$ in such a way as to keep unchanged, along the RG flows, the initial quadratic part $`S_0[\varphi ]\frac{1}{2}_p\varphi _p\varphi _pp^2K^1(p^2)`$. This additional condition imposed on the ERGE would completely eliminate the ambiguities in the definition (the invariances) of the RG transformation leaving no room for any redundant operators (see also chapt 5 of ).
If we understand correctly the procedure, it is similar to (but perhaps more general than) that proposed in to eliminate the redundant operator (24). One fixes the arbitrariness associated to the invariance (reparametrization invariance of section 2.2.4) in order that any RG flow remains orthogonal to the redundant direction(s).
The authors of do not specify what happens in the case where the unavoidable truncation (or approximation) used to study the ERGE breaks some invariances<sup>11</sup><sup>11</sup>11In particular it is known that the derivative expansion developped within ERGE with smooth cutoffs breaks the reparametrization invariance.. In fact, as already mentioned at the end of section 2.2.4, the freedom associated with the redundant directions allows us to search for the region of minimal break of invariance in the space of interactions . If this freedom is suppressed one may not obtain, for universal quantities such as the critical exponents, the optimal values compatible with the approximation (or truncation) used.
### 2.6 ERGE for the Legendre Effective Action
The first derivation<sup>12</sup><sup>12</sup>12The first mention to the Legendre transformation in the expression of the ERGE has been formulated by E. Brézin in a discussion following a talk by Halperin . of an ERGE for the Legendre effective action $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ \[defined in eq. (7)\] has been carried out with a sharp cutoff by Nicoll and Chang (see also ). Their aim was to simplify the obtention of the $`\epsilon `$-expansion from the ERGE. More recent obtentions of this equation with a smooth cutoff are due to Bonini, D’Attanasio and Marchesini , Wetterich , Morris and Ellwanger .
A striking fact arises with the Legendre transformation: the running cutoff $`\mathrm{\Lambda }`$ (or intermediate-scale momentum cutoff, i.e. associated to the “time” $`t`$) acts as an IR cutoff and physical Green functions are obtained in the limit $`\mathrm{\Lambda }0`$ . The reason behind this result is simple to understand. The generating functional $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ is obtained by integrating out all modes (from $`\mathrm{\Lambda }_0=\mathrm{}`$ to $`0`$). If an intermediate cutoff $`\mathrm{\Lambda }\mathrm{\Lambda }_0`$ is introduced and integration is performed only in the range $`[\mathrm{\Lambda },\mathrm{\Lambda }_0]`$, then for the integrated modes (thus for the effective $`\mathrm{\Gamma }_\mathrm{\Lambda }\left[\mathrm{\Phi }\right]`$), $`\mathrm{\Lambda }`$ is an IR cutoff while for the unintegrated modes (for the effective $`S_\mathrm{\Lambda }\left[\varphi \right]`$), $`\mathrm{\Lambda }`$ is an UV cutoff. There is an apparent second consequence: contrary to the ERGE for $`S_\mathrm{\Lambda }\left[\varphi \right]`$, the ERGE for $`\mathrm{\Gamma }_\mathrm{\Lambda }\left[\mathrm{\Phi }\right]`$ will depend on both an IR ($`\mathrm{\Lambda }`$) and an UV ($`\mathrm{\Lambda }_0`$) cutoffs<sup>13</sup><sup>13</sup>13One could think a priori that as a simple differential equation, the ERGE would be instantaneous (would depend only on $`\mathrm{\Lambda }`$ and not on some initial scale $`\mathrm{\Lambda }_0`$), but it is an integro-differential equation.. However it is possible to send $`\mathrm{\Lambda }_0`$ to infinity, see section 2.6.2 for a short discussion of this point.
#### 2.6.1 Sharp cutoff version
In an ERGE for the Legendre transformation $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ \[eq. 7\] is derived (See also ). It reads:
$$\dot{\mathrm{\Gamma }}=𝒢_{\text{dil}}\mathrm{\Gamma }+\frac{1}{2}\frac{d\mathrm{\Omega }}{\left(2\pi \right)^d}\mathrm{ln}\left\{\mathrm{\Gamma }_{q,q}\underset{p}{}\underset{p^{}}{}\mathrm{\Gamma }_{qp}\left(\mathrm{\Gamma }^1\right)_{pp^{^{}}}\mathrm{\Gamma }_{p^{^{}},q}\right\}$$
(37)
in which:
$$\mathrm{\Gamma }_{kk^{}}=\frac{\delta ^2\mathrm{\Gamma }}{\delta \mathrm{\Phi }_k\delta \mathrm{\Phi }_k^{}}$$
the momentum $`q`$ lies on the shell $`q=|𝐪|=1`$ while the integrations on $`p`$ and $`p^{}`$ are performed inside the shell $`]1,\mathrm{\Lambda }_0/\mathrm{\Lambda }]`$ where $`\mathrm{\Lambda }_0`$ is some initial cutoff ($`\mathrm{\Lambda }_0>\mathrm{\Lambda }`$) and $`\mathrm{\Omega }`$ is the surface of the $`d`$-dimensional unit sphere \[$`\mathrm{\Omega }=\left(2\pi \right)^dK_d`$\].
One sees that $`\mathrm{\Lambda }`$ is like an IR cutoff and that (37) depends on the initial cutoff $`\mathrm{\Lambda }_0`$.
#### 2.6.2 Smooth cutoff version
We adopt notations which are close to the writing of (32) and we consider the Wilson effective action with an “additive” IR cutoff $`\mathrm{\Lambda }`$ such that:
$$S_\mathrm{\Lambda }[\varphi ]\frac{1}{2}_p\varphi _p\varphi _pC^1(p,\mathrm{\Lambda })+S_{\mathrm{\Lambda }_0}[\varphi ]$$
(38)
in which $`C(p,\mathrm{\Lambda })`$ is an additive infrared cutoff function which is small for $`p<\mathrm{\Lambda }`$ (tending to zero as $`p0`$) and $`p^2C(p,\mathrm{\Lambda })`$ should be large for $`p>\mathrm{\Lambda }`$ . Due to the additive character of the cutoff function, $`S_{\mathrm{\Lambda }_0}[\varphi ]`$ is the entire action (involving the kinetic term contrary to eq. (32) and to where the cutoff function was chosen multiplicative). In this section, because $`C`$ is naturally dimensioned<sup>14</sup><sup>14</sup>14It is not harmless that $`C`$ is dimensionful because the anomalous dimension $`\eta `$ may be a part of its dimension (see and section 3.1.3) \[contrary to $`K`$ in (32)\], all the dimensions are implicitly restored in order to keep the same writing as in the original papers. The ultra-violet regularization is provided by $`\mathrm{\Lambda }_0`$ and needs not to be introduced explicitly (see and below). The Legendre transformation is defined as:
$$\mathrm{\Gamma }[\mathrm{\Phi }]+\frac{1}{2}_p\mathrm{\Phi }p\mathrm{\Phi }_pC^1(p,\mathrm{\Lambda })=W[J]+J.\mathrm{\Phi }$$
in which $`W[J]`$ and $`\mathrm{\Phi }`$ are defined as usual \[see eq (7)\] from (38).
Then the ERGE reads:
$$\dot{\mathrm{\Gamma }}=𝒢_{\text{dil}}\mathrm{\Gamma }+\frac{1}{2}_p\frac{1}{C}\mathrm{\Lambda }\frac{C}{\mathrm{\Lambda }}\left(1+C\mathrm{\Gamma }_{p,p}\right)^1$$
(39)
When the field is no longer a pure scalar but carries some supplementary internal degrees of freedom and becomes a vector, a spinor or a gauge field etc…, a more compact expression using the trace of operators is often used:
$$\dot{\mathrm{\Gamma }}=𝒢_{\text{dil}}\mathrm{\Gamma }+\frac{1}{2}\text{tr}\left[\frac{1}{C}\mathrm{\Lambda }\frac{C}{\mathrm{\Lambda }}\left(1+C\frac{\delta ^2\mathrm{\Gamma }}{\delta \mathrm{\Phi }\delta \mathrm{\Phi }}\right)^1\right]$$
(40)
which, for example, allows us to include the generalization to $`N`$ components in a unique writing ($`\frac{\delta ^2\mathrm{\Gamma }}{\delta \mathrm{\Phi }\delta \mathrm{\Phi }}`$ has then two supplementary indices $`\alpha `$ and $`\beta `$ corresponding to the derivatives with respect to the fields $`\mathrm{\Phi }_\alpha `$ and $`\mathrm{\Phi }_\beta `$, the trace is relative to both the momenta and the indices).
The equations (39,40) may be obtained, as in , using the trick of eqs.(35 and 36), but see also .
For practical computations it is actually often quite convenient (for example, see section 4.2) to write the flow equation (39) as follows :
$$\dot{\mathrm{\Gamma }}=𝒢_{\text{dil}}\mathrm{\Gamma }+\frac{1}{2}\text{tr}\stackrel{~}{}_t\mathrm{ln}\left(C^1+\frac{\delta ^2\mathrm{\Gamma }}{\delta \mathrm{\Phi }\delta \mathrm{\Phi }}\right)$$
(41)
with $`\stackrel{~}{}_t\mathrm{\Lambda }\frac{}{\mathrm{\Lambda }}`$ acting only on $`C`$ and not on $`\mathrm{\Gamma }`$, i.e. $`\stackrel{~}{}_t=\left(C^1/t\right)\left(/C^1\right)`$.
Wetterich’s expression of the ERGE is identical<sup>15</sup><sup>15</sup>15The correspondence between Wetterich’s notations and ours is as follows: $`\mathrm{\Lambda }k`$; $`C1/R_k`$; $`tt`$. to eqs. (39, 40, 41). Its originality is in the choice of the cutoff function $`C^1`$: to make the momentum integration in (39) converge, a cutoff function is introduced such that only a small momentum interval $`p^2\mathrm{\Lambda }^2`$ effectively contributes (see also the review in ). This feature, which avoids an explicit UV regularization, allows calculations in models where the UV regularization is a delicate matter (e.g. non-Abelian gauge theories). In fact, as noticed in , the ERGE only requires momenta $`p\mathrm{\Lambda }`$ and should not depend on $`\mathrm{\Lambda }_0\mathrm{\Lambda }`$ at all. Indeed, once a finite ERGE is obtained, the flow equation for $`\mathrm{\Gamma }_\mathrm{\Lambda }\left[\mathrm{\Phi }\right]`$ is finite and provides us with an “ERGE”-regularization scheme which is specified by the flow equation, the choice of the infrared cutoff function $`C`$ and the “initial condition” $`\mathrm{\Gamma }_\mathrm{\Lambda }`$ . Most often, there is no need for any UV regularization and the limit $`\mathrm{\Lambda }_0\mathrm{}`$ may be taken safely. In this case, the cutoff function chosen by Wetterich has the following form (up to some factor):
$`C(p,\mathrm{\Lambda })`$ $`=`$ $`{\displaystyle \frac{1f(\frac{p^2}{\mathrm{\Lambda }^2})}{p^2f(\frac{p^2}{\mathrm{\Lambda }^2})}}`$ (42)
$`f(x)`$ $`=`$ $`\text{e}^{2ax^b}`$ (43)
in which the two parameters $`a`$ and $`b`$ may be adjusted to vary the smoothness of the cutoff function.
Although it is almost an anticipation on the expansions (local potential approximation and derivative expansion) considered in the parts to come, we find it worthwhile indicating here the exact equation satisfied by the effective potential which is often used by Wetterich and co-workers (see their review in this volume). Thus following Wetterich , we write the effective (Legendre) action $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ for $`O(N)`$-symmetric systems as follows<sup>16</sup><sup>16</sup>16It is customary to introduce the variable $`\rho `$ for $`O(N)`$ systems because this allows better convergences in some cases (see section 3.3) but in the particular case $`N=1`$ the symmetry assumption is not required.:
$`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ $`=`$ $`{\displaystyle \text{d}^dx\left[U(\rho )+\frac{1}{2}^\mu \mathrm{\Phi }_\alpha 𝒵(\rho ,\mathrm{})_\mu \mathrm{\Phi }_\alpha +\frac{1}{4}^\mu \rho 𝒴(\rho ,\mathrm{})_\mu \rho \right]}`$
$`\rho `$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Phi }^2`$
where the symbol $`\mathrm{}`$ stands for $`_\mu ^\mu `$ and acts only on the right (summation over repeated indices is assumed). Then the exact evolution equation for the effective potential $`U(\rho )`$ follows straightforwardly from (41):
$`\dot{U}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _p}{\displaystyle \frac{1}{C^2}}\mathrm{\Lambda }{\displaystyle \frac{C}{\mathrm{\Lambda }}}\left[{\displaystyle \frac{N1}{M_0}}+{\displaystyle \frac{1}{M_1}}\right]+dUd_\varphi U^{}`$ (44)
$`M_0`$ $`=`$ $`𝒵(\rho ,p^2)p^2+C^1+U^{}`$ (45)
$`M_1`$ $`=`$ $`\left[𝒵(\rho ,p^2)+\rho 𝒴(\rho ,p^2)\right]p^2+C^1+U^{}+2\rho U^{\prime \prime }`$ (46)
in which $`\dot{U}`$ stands for $`U(\rho ,t)/t`$ and $`U^{}`$ and $`U^{\prime \prime }`$ refer to the first and second (respectively) derivatives with respect to $`\rho `$. The two last terms in (44) come from $`𝒢_{\text{dil}}\mathrm{\Gamma }`$ which was not explicitly considered in .
The interest of dealing with an additive cutoff function is that one may easily look for the classes of $`C`$ that allow a linear realization of the reparametrization invariance . It is found that $`C`$ must be chosen as:
$$C(p,\mathrm{\Lambda })=\mathrm{\Lambda }^2\left(\frac{p^2}{\mathrm{\Lambda }^2}\right)^k$$
with $`k`$ an integer such that $`k>d/21`$ to have UV convergence. With this choice, the derivative expansion preserves the reparametrization invariance and $`\eta `$ is uniquely defined (see part 4.1.2). On the contrary, because the cutoff function corresponding to eqs. (42, 43) has an exponential form, the derivative expansion does not provide us with a uniquely defined value of $`\eta `$ (see section 4.3).
##### Sharp cutoff limit
It is possible to obtain the sharp cutoff limit from eqs.(39, 40) provided one is cautious in dealing with the value at the origin of the Heaviside function $`\theta (0)`$ (which is not equal to $`\frac{1}{2}`$) . One obtains the sharp cutoff limit of the flow equation :
$$\dot{\mathrm{\Gamma }}=𝒢_{\text{dil}}\mathrm{\Gamma }+\frac{1}{2}_p\frac{\delta (p1)}{\gamma (p)}\left[\widehat{\mathrm{\Gamma }}\left(1+G\widehat{\mathrm{\Gamma }}\right)^1\right](𝐩,𝐩)$$
(47)
in which the field independent full inverse propagator $`\gamma (p)`$ has been separated from the two-point function :
$$\frac{\delta ^2\mathrm{\Gamma }[\mathrm{\Phi }]}{\delta \mathrm{\Phi }\delta \mathrm{\Phi }}(𝐩,𝐩^{})=\gamma (p)\widehat{\delta }(𝐩+𝐩^{})+\widehat{\mathrm{\Gamma }}[\mathrm{\Phi }](𝐩,𝐩^{})$$
so that $`\widehat{\mathrm{\Gamma }}[0]=0`$, and $`G(p)=\theta (p1)/\gamma (p)`$.
### 2.7 Equivalent fixed points and reparametrization invariance.
To illustrate the line of equivalent fixed points which arises when the reparametrization invariance is satisfied (see the end of section 2.2.4), we consider here the pure Gaussian case of the Wilson ERGE (see also the appendix of and ). No truncation is needed to study the Gaussian case, the analysis below is thus exact.
The action is assumed to have the following pure quadratic form:
$$S_G[\varphi ]=\frac{1}{2}_p\varphi _p\varphi _pR\left(p^2\right)$$
The effect of $`𝒢_{\text{dil}}`$ in (30) yields (the prime denotes the derivative with respect to $`p^2`$):
$$𝒢_{\text{dil}}S_G=\left(\frac{d}{2}d_\varphi \right)_p\varphi _p\varphi _pR_p\varphi _p\varphi _pp^2R^{}$$
while the remaining part (coarsening) gives (up to neglected constant terms):
$$𝒢_{\text{tra}}S_G=_p\varphi _p\varphi _p\left[\left(c+2p^2\right)\left(RR^2\right)p^2R^{}\right]$$
Adding the two contributions to $`\dot{S}_G`$, choosing $`d_\varphi =\frac{d}{2}`$ as in and imposing that the fixed point is reached ($`\dot{S}_G=0`$), one obtains:
$`c`$ $`=`$ $`1`$
$`R^{}(p^2)`$ $`=`$ $`{\displaystyle \frac{zp^2}{\text{e}^{2p^2}+zp^2}}`$
This is a line of (Gaussian) fixed points parametrized by $`z`$. To reach this line, the parameter $`c`$ must be adjusted to $`1`$ \[i.e., following eq. (31), $`\eta =0`$\], the fixed points on the line are equivalent.
The same analysis may be done with the same kind of conclusions with the Polchinski ERGE (see ). For the Wegner-Houghton version (29), the situation is very different in nature . The same kind of considerations yields:
$$R^{}(p^2)=p^{2\eta }$$
in which $`\eta `$ is undetermined. “This phenomenon is of quite different nature as the similar one described above. There the whole line shares the same critical properties, here it does not; there we have well-behaved actions throughout the fixed line, here nearly all of them are terribly non-local (in the sense that we cannot expand the action integrand in a power series of $`p^2`$). What happens is that we have one physical FP (the $`\eta =0`$ case) and a line of spurious ones.
### 2.8 Approximations and truncations
How to deal with integro-differential equations is not known in general. In the case of the RG equations, one often had recourse to perturbative expansions such as the usual perturbation in powers of the coupling but also the famous $`\epsilon `$–expansion (where $`\epsilon =4d`$), the $`1/N`$–expansion or expansions in the dimensionality $`2d`$, let us mention also an expansion exploiting the smallness of the critical exponent<sup>17</sup><sup>17</sup>17The expansion requires also a truncation in powers of the field and some re-expansion of $`\epsilon /4=1d/4`$ in powers of $`\eta ^{1/2}`$ . $`\eta `$ .
When no small parameter can be identified or when one does not want to consider a perturbative approach, one must truncate the number of degrees of freedom involved in order to reduce the infinite system of coupled differential equations to a finite system. The way truncations are introduced is of utmost importance as one may learn from the development of the scaling field method (see ). The ERGE is a useful starting point to develop approximate approach. If from the beginning it has been exclusively seen as a useful tool for the investigation of the $`\epsilon `$-expansion , two complementary approaches to nonperturbative truncations have been proposed:
* an expansion of the effective action in powers of the derivatives of the field :
$$S[\varphi ]=\text{d}^dx\left\{V(\varphi )+\frac{1}{2}Z(\varphi )\left(_\mu \varphi \right)^2+O(^4)\right\}$$
(48)
which is explicitly considered in this review (see sections 3 and 4.1)
* an expansion in powers of the field for the functional $`\mathrm{\Gamma }[\mathrm{\Phi }]`$ (see also ):
$$\mathrm{\Gamma }[\mathrm{\Phi }]=\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{n!}\left(\underset{k=0}{\overset{n}{}}\text{d}^dx_k\mathrm{\Phi }(x_k)\right)\mathrm{\Gamma }^{(n)}(x_1,\mathrm{},x_n)$$
(49)
The flow equations for the $`1`$PI $`n`$–point functions $`\mathrm{\Gamma }^{(n)}`$ are obtained by functional differentiation of the ERGE. The distinction between $`S`$ in (48) and $`\mathrm{\Gamma }`$ in (49) is not essential, we can introduce the two approximations for both $`S`$ and $`\mathrm{\Gamma }`$.
Because the derivative expansion corresponds to small values of $`p`$, it is naturally (quantitatively) adapted to the study of the large distance or low energy physics like critical phenomena or phase transition. We will see, however, that at a qualitative level it is suitable to a general discussion of many aspects of field theory (like the continuum limit, see part 3). Obviously, when bound state formation or nonperturbative momentum dependences are studied, the expansion (49) seems better adapted (see, for examples, ).
Only a few terms of such series will be calculable in practice, since the number of invariants increases rapidly for higher orders (see section 4.1).
### 2.9 Scaling field representation of the ERGE
To date the most expanded approximate method for the solution of the ERGE has been developed by Golner and Riedel in (see also and especially ) from the scaling-field representation. The idea is to introduce the expansion (23) into the ERGE which is thus transformed into an infinite hierarchy of nonlinear ordinary differential equations for scaling fields. For evident reasons<sup>18</sup><sup>18</sup>18The eigenoperators $`𝒪_i^{}`$ and eigenvalues $`\lambda _i`$ of the Gaussian fixed point can be determined exactly., the fixed point chosen for $`S^{}`$ is the Gaussian fixed point. Approximate solutions may be obtained by using truncations and iterations. The approximations appear to be effective also in calculations of properties, like phase diagrams, scaling functions, crossover phenomena and so on (see ).
Because they are unjustly not often mentioned in the literature, we find it fair to extract from the following estimates for $`N=1`$ (see the paper for estimates corresponding to other values of $`N`$):
$`\nu `$ $`=`$ $`0.626\pm 0.009`$
$`\eta `$ $`=`$ $`0.040\pm 0.007`$
$`\omega `$ $`=`$ $`0.855\pm 0.07`$
$`\omega _2`$ $`=`$ $`1.67\pm 0.11`$
$`\omega _5`$ $`=`$ $`2.4\pm 0.4`$
in which $`\omega _2`$ is the second correction-to-scaling exponent and the subscript “5” in $`\omega _5`$ refers to a $`\varphi ^5`$ interaction present in the action and which would be responsible for correction-to-scaling terms specific to the critical behavior of fluids, as opposed to the Ising model which satisfies the symmetry $`Z_2`$, see and also section 3.2.3.
### 2.10 Renormalizability, continuum limits and the Wilson theory
#### 2.10.1 Wilson’s continuum limit
In the section 12.2 of , a nonperturbative realization of the renormalization of field theory is schematically presented. The illustration is done with the example of a fixed point with one relevant direction (i.e. a fixed point which controls the criticality of magnetic systems in zero magnetic field). The resulting renormalized field theory is purely massive (involving only one parameter: a mass).
We find satisfactory<sup>19</sup><sup>19</sup>19Except the expression “critical manifold, which consists of all bare actions yielding a given massless continuum limit”, see section 3.4.1. the presentation of this continuum limit by Morris in (relative to the discussion of his fig. 3, see fig. 3 of the present paper for an illustration with actual RG trajectories), and we reproduce it here just as it is.
In the infinite dimensional space of bare actions, there is the so-called critical manifold, which consists of all bare actions yielding a given massless continuum limit. Any point on this manifold – i.e. any such bare action – flows under a given RG towards its fixed point; local to the fixed point, the critical manifold is spanned by the infinite set of irrelevant operators. The other directions emanating out of the critical manifold at the fixed point, are spanned by relevant and marginally relevant perturbations (with RG eigenvalues $`\lambda _i>0`$ and $`\lambda _i=0`$, respectively). \[In the example of and in fig. 3, there is only one relevant perturbation.\] Choosing an appropriate parametrization of the bare action, we move a little bit away from the critical manifold. The trajectory of the RG will to begin with, move towards the fixed point, but then shoot away along one of the relevant directions towards the so-called high temperature fixed point which represents an infinitely massive quantum field theory.
To obtain the continuum limit, and thus finite masses, one must now tune the bare action back towards the critical manifold and at the same time, reexpress physical quantities in renormalised terms appropriate for the diverging correlation length. In the limit that the bare action touches the critical manifold, the RG trajectory splits into two: a part that goes right into the fixed point, and a second part that emanates out from the fixed point along the relevant directions. This path is known as a Renormalised Trajectory (RT). The effective actions on this path are ‘perfect actions’ .
The continuum limit so defined “at” a critical fixed point has been used by Wilson to show that the $`\varphi ^6`$-field-theory in three dimensions has a nontrivial continuum limit involving no coupling constant renormalization (i.e. the continuum limit involves only a mass as renormalized parameter, but see also ). The fixed point utilized in the circumstances is the Wilson-Fisher (critical) fixed point. In fact, exactly the same limit would have been obtained starting with a $`\varphi ^4`$\- or a $`\varphi ^8`$-bare-theory, since it is the symmetry of the bare action which is important and not the specific form chosen for the initial (bare) interaction ($`\varphi ^4`$, $`\varphi ^6`$, or $`\varphi ^8`$ are all elements of the same $`Z_2`$-symmetric scalar theory, see section 3.4.1 for more details).
It is noteworthy that one (mainly) exclusively presents the Wilson continuum limit as it is illustrated in , i.e. relatively to a critical point which possesses only one relevant direction. Obviously one may choose any fixed point with several relevant directions (as suggested in the Morris presentation reproduced above). The relevant parameters provide the renormalized couplings of the continuum limit. For example the Gaussian fixed point in three dimensions for the scalar theory (a tricritical fixed point with two relevant directions) yields a continuum limit which involves two renormalized parameters (see an interesting discussion with Brézin following a talk given by Wilson ). Indeed, that continuum limit is nothing but the so-called $`\varphi ^4`$-field theory used successfully in the investigation, by perturbative means, of the critical properties of statistical systems below four dimensions and which is better known as the “Field theoretical approach to critical phenomena (see for a review). The scalar field theory below four dimensions defined “at” the Gaussian fixed point involves a mass and a (renormalized) $`\varphi ^4`$-coupling “constant”.
“At” the Gaussian fixed one may also define a massless renormalized theory. To reach this massless theory, one must inhibit the direction of instability<sup>20</sup><sup>20</sup>20The relevant direction that points towards the most stable high-temperature fixed point. of the Gaussian fixed point toward the massive sector. As a consequence, the useful space of bare interactions is limited to the critical manifold alluded to above by Morris (there is one parameter to be adjusted in the bare action). The discussion is as previously but with one dimension less for the space of bare interactions: the original whole space is replaced by the critical submanifold and this latter by the tricritical submanifold. Notice that the massless continuum limit so defined really involves a scale dependent parameter: the remaining relevant direction of the Gaussian fixed point which corresponds to the $`\varphi ^4`$-renormalized coupling (see section 3.4.1 for more details). It differs however from the massless theories sketchily defined by Morris as fixed point theories . Most certainly no mass can be defined right at a fixed point (there is scale invariance there and a mass would set a scale) but at the same time the theory would also have no useful parameter at all (no scale dependent parameter) since, by definition, right at a fixed point nothing changes, nothing happens, there is nothing to describe.
An important aspect of the Wilson continuum limit is the resulting self-similarity emphasized rightly in several occasions by Morris and collaborators (see section 3.4.1). This notion expresses the fact that in a properly defined continuum limit, the effects of the infinite number of degrees of freedom involved in a field theory are completely represented by a (very) small number of flowing (scale dependent) parameters (the relevant parameters of a fixed point): the system is self-similar in the sense that it is exactly (completely) described by the same finite set of parameters seen at different scales (see section 3.4.1). This is exactly what one usually means by renormalizability in perturbative field theory. However the question is nonperturbative in essence. For example, the $`\varphi ^4`$-field theory in four dimensions is perturbatively renormalizable, but it is not self-similar at any scale and especially in the short distance regime due to the UV “renormalon” problems (for a review see ) which prevent the perturbatively renormalized coupling constant to carry exactly all the effects of the other (an infinite number) degrees of freedom: it is not a relevant parameter for a fixed point (see section 2.10.2, 3.4.1 and ).
#### 2.10.2 Polchinski’s effective field theories
There has been a renewed interest of field theoreticians in the ERGE since Polchinski’s paper in 1984. From the properties of the RG flows generated by an ERGE (see section 2.5.2) and by using only “very simple bounds on momentum integrals”, Polchinski presented an easy proof of the perturbative renormalizability of the $`\varphi ^4`$ field theory in four dimensions (see also ). This paper had a considerable success. One may understand the reasons of the resulting incipient interest of field theoreticians in the ERGE, let us cite for example:
Proofs of renormalizability to all orders in perturbation theory were notoriously long and complicated (involving notions of graph topologies, skeleton expansions, overlapping divergences, the forest theorem, Weinberg’s theorem, etc.), …” .
The enthusiasm of some was so great that one has sometimes referred to Polchinski’s presentation, really based on the Wilson RG theory, as the Wilson-Polchinski theory. However the strategy relative to the construction of the continuum limit (modern expression for “renormalizability”) is rather opposite to the ideas of Wilson because they are perturbative in essence. Indeed, while the reference to a fixed point is essential in the Wilson construction of the continuum limit, Polchinski does not need any explicit reference to a fixed point, but in fact refers implicitly to the Gaussian fixed point<sup>21</sup><sup>21</sup>21As in perturbation theory.. A classification of parameters as relevant, irrelevant and marginal is given using a purely classical dimensional analysis (referring to the Gaussian fixed point, see for example). In the Polchinski view, the marginal parameters are then considered as being relevant although in some cases (as the scalar case, see footnote 5) they may actually be (marginally) irrelevant.
The arguments may then lead to confusions. The notion of relevant parameter (the natural candidate for the renormalized parameter), which, in the Wilson theory represents an unstable direction of a fixed point (one goes away from the fixed point which thus in the occasion displays an ultraviolet stability or attractivity) has been replaced in the Polchinski point of view (for the $`\varphi _4^4`$ field) by the least irrelevant parameter which controls the final approach to a fixed point (one goes toward the fixed point which thus presents an infrared stability or attractivity). Thus the renormalized coupling resulting from the “proof” controls only the infrared (large distances or low energy) regime of the scalar theory. The field theory so constructed is actually an “effective” field theory (valid in the infrared regime) and not a field theory well defined in the short distance regime (e.g. even after the “proof” the $`\varphi ^4`$ field theory in four dimensions remains trivial due to the lack of ultraviolet stable fixed point).
Of course, if by chance the Gaussian fixed point is ultraviolet stable (asymptotically free field theories), then the marginal coupling is truly a relevant parameter for the Gaussian fixed point and the perturbatively constructed field theory exists beyond perturbation (in the Wilson sense of section 2.10.1). In that case, one may use the Polchinski approach to prove the existence of a continuum limit (see some references in ).
It is fair to specify that, in several occasions in , Polchinski has emphasized the perturbative character of his proof which is only equivalent to (but simpler than) the usual perturbative proof. In order to be clear, let us precisely indicate the weak point of Polchinski’s arguments which is clearly expressed in the discussion of the fig. 2 of , p. 274 one may read:
We can proceed in this way, thus defining the bare coupling $`\lambda _4^0`$ as a function of $`\lambda _4^\text{R}`$, $`\mathrm{\Lambda }_\text{R}`$, and $`\mathrm{\Lambda }_0`$. Now take $`\mathrm{\Lambda }_0\mathrm{}`$ holding $`\mathrm{\Lambda }_\text{R}`$ and $`\lambda _4^\text{R}`$ fixed.
The objection is that, in Wilson’s theory (i.e. nonperturbatively) it is impossible to make sense to the second sentence without an explicit reference to an (eventually nontrivial) ultra-violet stable fixed point. In perturbation theory, however, no explicit reference to a fixed point is needed since, order by order, terms proportional to, say $`(\mathrm{\Lambda }_\text{R}/\mathrm{\Lambda }_0)^2`$, give exactly zero in the limit “$`\mathrm{\Lambda }_0\mathrm{}`$ holding $`\mathrm{\Lambda }_\text{R}`$ and $`\lambda _4^\text{R}`$ fixed”. However this limit introduces singularities in the perturbative expansion: the famous ultraviolet renormalons which make it ambiguous to resum the perturbative series for the scalar field in four dimensions (for a review on the renormalons see ). See section 3.4.1 for more details.
Actually, using the nonperturbative framework of the ERGE to present a proof of the perturbative renormalizability might be seen as a misunderstanding of the Wilson theory.
## 3 Local potential approximation: A textbook example
### 3.1 Introduction
The local potential approximation (LPA) of the ERGE (the momentum-independent limit of the ERGE) allows to consider all powers of $`\varphi `$. The approximation still involves an infinite number of degrees of freedom which are treated on the same footing within a nonlinear partial differential equation for a simple function $`V(\varphi )`$ \[$`V`$ is the (local) potential, $`\varphi `$ is assumed to be a constant field and thus, except the kinetic term, the derivatives $`\varphi `$ of $`\varphi (x)`$ are all neglected in the ERGE\].
LPA of the ERGE is the continuous version of the Wilson approximate recursion formula (see also , p. 117) which is a discrete (and approximate) realization of the RG (the momentum scale of reference is reduced by a factor two). As shown by Felder , LPA is also similar to a continuous version of the hierarchical model .
This approximation has been first considered in (see also ) from the sharp cutoff version of the ERGE of Wegner and Houghton , it has been rederived by Tokar by using approximate functional integrations and rediscovered by Hasenfratz and Hasenfratz .
LPA amounts to assuming that the action $`S[\varphi ]`$ reduces to the following form:
$$S[\varphi ]=\text{d}^dx\frac{z}{2}\left(_\mu \varphi \right)^2+V(\varphi )$$
(50)
in which $`z`$ is a pure number (a constant usually set equal to 1) and , to set the ideas, $`V(\varphi )`$ is a simple function of $`\varphi _0`$, it has the form (symbolically)
$$V(\varphi )=\underset{n}{}u_n\left(\varphi _0\right)^n\delta (0)$$
(51)
The infinity carried by the delta function would be absent in a treatment at finite volume, it reflects the difficulties of selecting one mode out of a continuum set, such ill-defined factors may be removed within a rescaling of $`\varphi `$ . In the derivation of the approximate equation, instead of using (51) we find it convenient to deal with
$$V(\varphi )=\underset{n}{\overset{\mathrm{}}{}}u_n_{p_1\mathrm{}p_n}\varphi _{p_1}\mathrm{}\varphi _{p_n}\widehat{\delta }\left(𝐩_1+\mathrm{}+𝐩_n\right)$$
(52)
in which the $`u_n`$’s do not depend on the momenta \[see eq. (3)\] and to project onto the zero modes $`\varphi _0`$ of $`\varphi `$ at the end of the calculation.
Eq. (50) is identical to eq. (48) in which $`Z(\varphi )`$ is set equal to the constant $`z`$. LPA may also be considered as the zeroth order of a systematic expansion in powers of the (spatial) derivative of the field (derivative expansion) (see part 4).
In the following, $`\phi `$ stands for $`\varphi _0`$ and primes denote derivatives with respect to $`\phi `$ (at fixed $`t`$):
$`\phi `$ $``$ $`\varphi _0`$
$`V^{}(\phi ,t)`$ $`=`$ $`{\displaystyle \frac{V}{\phi }}|_t`$ (53)
Frequently, as in the present review, $`V^{}(\phi ,t)`$ is replaced by $`f(\phi ,t)`$.
Let us consider the expressions of LPA for the various ERGE’s introduced in part 2.
#### 3.1.1 Sharp cutoff version
The derivation of the local potential approximation for eq. (29) (sharp cutoff version of the LPA of the ERGE) is well known, we only give the result (for more details see ). It reads:
$$\dot{V}=\frac{K_d}{2}\mathrm{ln}\left[z+V^{\prime \prime }\right]+dVd_\varphi \phi V^{}$$
(54)
in which $`K_d`$ is given by (2). The non logarithmic terms come from the contribution $`𝒢_{\text{dil}}S`$ in (29). As usual in field theory, we neglect the field independent contributions \[“const” in (29)\]) to the effective potential $`V`$.
The $`t`$-dependence is entirely carried by the coefficient $`u_n(t)`$ in (52) while $`z`$ is considered as being independent of $`t`$. This condition is required for consistency of the approximation (it prevents the ERGE from generating contribution to the kinetic term: there is no wave function renormalization). Writing down explicitly this condition (namely $`\dot{z}=0`$) provides us with the relation:
$$d_\varphi =\frac{d2}{2}$$
in other word $`\eta =0`$, i.e. the anomalous part of the dimension of the field is zero. This is a characteristic feature of the LPA.
The dependence on $`z`$ in (54) may be removed (up to an additive constant) by the simplest (or naive) change of normalization of the field $`\phi \phi \sqrt{z}`$. \[The exact version (29) is invariant under the same change.\] In order to avoid the useless additive constant terms generated in (54), it is frequent to write down the evolution equation for the derivative $`f=V^{}`$, it comes :
$$\dot{f}=\frac{K_d}{2}\frac{f^{\prime \prime }}{z+f^{}}+\left(1+\frac{d}{2}\right)f+\left(1\frac{d}{2}\right)\phi f^{}$$
(55)
Notice that one could eliminate the factor $`K_d`$ by the change $`f(\phi ,t)\lambda f(\phi /\lambda ,t)`$ with $`K_d\lambda ^2=1`$.
It is interesting also to write down the ERGE in the same approximation when the number of components $`N`$ of the field is variable. With a view to eventually consider large values of $`N`$, it is convenient to redefine the action and the field as follows:
$$SNS\left[\frac{\varphi }{\sqrt{N}}\right]$$
then in the case of O$`(N)`$ symmetric potential, the LPA of (29) yields:
$$\dot{V}=\frac{K_d}{2N}\left\{\left(N1\right)\mathrm{ln}\left[z+\frac{V^{}}{\phi }\right]+\mathrm{ln}\left[z+V^{\prime \prime }\right]\right\}+dVd_\varphi \phi V^{}$$
(56)
or
$$\dot{f}=\frac{K_d}{2N}\left\{\left(N1\right)\frac{\phi f^{}f}{z\phi ^2+\phi f}+\frac{f^{\prime \prime }}{z+f^{}}\right\}+\left(1+\frac{d}{2}\right)f+\left(1\frac{d}{2}\right)\phi f^{}$$
(57)
It may be also useful to express that, in the O$`(N)`$ symmetric case, $`V`$ is a function of $`\phi ^2`$. By setting $`s=\phi ^2`$ and $`u=2`$d$`V/`$d$`s`$, one obtains :
$$\dot{u}=\frac{K_d}{N}\left[\frac{3u^{}+2su^{\prime \prime }}{1+u+2su^{}}+(N1)\frac{u^{}}{1+u}\right]+2u+(2d)su^{}$$
(58)
#### 3.1.2 The Wilson (or Polchinski) version
Due to the originality in introducing the arbitrary scaling parameter $`c`$, it is worthwhile writing down explicitly the LPA of eq. (30). This equation has first been derived in .
From the same lines as previously, it comes:
$$\dot{V}=c\left[V^{\prime \prime }\left(V^{}\right)^2+\phi V^{}\right]+dVd_\varphi \phi V^{}$$
in which $`c`$ is determined by the condition implying no wave function renormalization ($`\dot{z}=0`$) which reads:
$$d22d_\varphi +2c=0$$
For the Wilson choice $`d_\varphi =d/2`$ , it comes $`c=1`$ and from (31), $`\eta =0`$ (as it must). Consequently the LPA of (30) is :
$$\dot{V}=V^{\prime \prime }\left(V^{}\right)^2+\left(1\frac{d}{2}\right)\phi V^{}+dV$$
(59)
which, for the derivative $`f=V^{}`$ yields :
$$\dot{f}=f^{\prime \prime }2ff^{}+\left(1+\frac{d}{2}\right)f+\left(1\frac{d}{2}\right)\phi f^{}$$
Notice that, contrary to the sharp cutoff version, this equation (as the exact version) is not invariant under the simplest rescaling of the field $`\phi \phi \sqrt{z}`$.
In the O$`(N)`$-case, it comes:
$$\dot{V}=\frac{1}{N}V^{\prime \prime }\left(V^{}\right)^2+\frac{N1}{N}\frac{V^{}}{\phi }+\left(1\frac{d}{2}\right)\phi V^{}+dV$$
###### Polchinski’s version
In the LPA, eq. (34) yields exactly the same partial differential equation as previously \[eq. (59)\]. For general $`N`$, it has been studied by Comellas and Travesset under the following form:
$$\dot{u}=\frac{2s}{N}u^{\prime \prime }+\left[1+\frac{2}{N}+(2d)s2su\right]u^{}+(2u)u$$
in which $`s=\frac{1}{2}\phi ^2`$ and $`u=2dV/ds`$ \[the definition of the variables is different from (58)\].
#### 3.1.3 The Legendre transformed version
##### Sharp cutoff version
LPA for the Legendre transformed ERGE has been first written down by Nicoll, Chang and Stanley (see also ) with a sharp cutoff.
From eq. (37), it is easy to verify that one obtains the same equation as in the Wegner-Houghton case (also for general $`N`$). This is because at this level of approximation the effective potential coincides with its Legendre transformation (the Helmholtz potential coincides with the free energy). Also, the sharp cutoff limit leading to eq. (47) yields the correct LPA (54) .
##### Smooth cutoff version
With a smooth cutoff, the LPA of the Legendre version of the ERGE \[eq. (39)\] keeps the integro-differential form except for particular choices of the cutoff function and of the dimension $`d`$ \[e.g. see eq. (65) below\]. This is clearly an inconvenience.
From (44) but without introducing the substitution $`\phi \rho =\frac{1}{2}\phi ^2`$, we may easily write down the LPA for general $`N`$, it reads:
$`\dot{V}`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle _p}{\displaystyle \frac{1}{\stackrel{~}{C}^2}}\left(\left[p^2\stackrel{~}{C}^{}+\stackrel{~}{C}\right]\right)\left({\displaystyle \frac{N1}{M_0^{}}}+{\displaystyle \frac{1}{M_1^{}}}\right)+dVd_\varphi \phi V^{}`$ (60)
$`M_0^{}`$ $`=`$ $`p^2+\stackrel{~}{C}^1+V^{}/\phi `$ (61)
$`M_1^{}`$ $`=`$ $`p^2+\stackrel{~}{C}^1+V^{\prime \prime }`$ (62)
in which $`\stackrel{~}{C}`$ is the dimensionless version of the cutoff function $`C`$ of section 2.6.2 with $`C=\mathrm{\Lambda }^{2+\eta }\stackrel{~}{C}(p^2)`$ and $`\stackrel{~}{C}^{}=`$d$`\stackrel{~}{C}(p^2)/`$d$`p^2`$ ($`p`$ is there also dimensionless). Notice that, following , we have introduced the anomalous dimension $`\eta `$ by anticipation of the anomalous scaling behavior satisfied by the field in the close vicinity of a non trivial fixed point. In the approximation presently considered, $`\eta `$ vanishes and does not appear in (60) but it would have an effect at higher orders of the derivative expansion (see section 4.1).
With the particular choice of cutoff function given by (42, 43), eq. (60) may be written as follows:
$$\dot{V}=\frac{K_d}{4N}\left[\left(N1\right)L_0^d\left(V^{}/\phi \right)+L_0^d\left(V^{\prime \prime }\right)\right]+dVd_\varphi \phi V^{}$$
(63)
in which:
$$L_0^d\left(w\right)=2(2a)^{\frac{2d}{2b}}_0^{\mathrm{}}\text{d}yy^{\frac{d2}{2b}}\frac{\text{e}^y}{\left(1\text{e}^y\right)}\frac{1}{\left[1+\left(\frac{2a}{y}\right)^{\frac{1}{b}}\text{e}^yw\right]}$$
In the sharp cutoff limit $`b\mathrm{}`$ one has:
$$L_0^d\left(w\right)=2\mathrm{ln}\left(1+w\right)+\text{ const}$$
in which “const” is infinite, neglecting this usual infinity, one sees that (63) gives back the expression (54) for $`z=1`$. In order to avoid the infinite “const”, it is preferable to consider the flow equation for the derivative $`f=V^{}`$, in which case the function $`L_1^d\left(w\right)=\frac{}{w}L_0^d\left(w\right)`$ appears in the equation:
$$\dot{f}=\frac{K_d}{4N}\left[\left(N1\right)\left(f^{}/\phi f/\phi ^2\right)L_1^d\left(f/\phi \right)+f^{\prime \prime }L_1^d\left(f^{}\right)\right]$$
See for more details on the function $`L_n^d\left(w\right)`$.
An interesting expression of the flow equation for the Legendre transformed action $`\mathrm{\Gamma }`$ is obtained from the smooth cutoff version of Morris \[see eqs. (39, 40)\] with a pure power law (dimensionless) cutoff function of the form:
$$\stackrel{~}{C}(p^2)=p^{2k}$$
(64)
For $`d=3`$, $`k=1`$ and $`N=1`$, the LPA reads :
$$\dot{V}=\frac{1}{\sqrt{2+V^{\prime \prime }}}+3V\frac{1}{2}\phi V^{}$$
(65)
and for general $`N`$ :
$$\dot{V}=\frac{1}{\sqrt{2+V^{\prime \prime }}}\frac{N1}{\sqrt{2+V^{}/\phi }}+3V\frac{1}{2}\phi V^{}$$
(66)
The choice of the power law cutoff function (64) is dictated by the will to linearly realize the reparametrization invariance . With the sharp cutoff, the power law cutoff is the only known cutoff that satisfies the conditions required to preserve this invariance.
### 3.2 The quest for fixed points
Fixed points are essential in the RG theory. In field theory, they determine the nature of the continuum limits; in statistical physics they control the large distance physics of a critical system.
A fixed point is a solution of the equation
$$\dot{V}^{}=0$$
(67)
From the forms of the equations involved (see the preceding section), it is easy to see that $`V=0`$ (or $`V=`$const) is always a solution of the fixed point equation. This is the Gaussian fixed point. There are two other trivial fixed points which are only accounted for with the Wilson (or Polchinski) version. Following Zumbach , let us write the eq. (59) for the quantity $`\mu (\phi ,t)=\mathrm{exp}\left(V(\phi ,t)\right)`$:
$$\dot{\mu }=\mu ^{\prime \prime }+\left(1\frac{d}{2}\right)\phi \mu ^{}+d\mu \mathrm{ln}\mu $$
(68)
from which the following trivial fixed point solutions are evident (the notations differ from those used in ):
* $`\mu _G^{}=1`$, the Gaussian fixed point mentioned above
* $`\mu _{HT}^{}=\mathrm{exp}\left(\frac{1}{2}\phi ^2+\frac{1}{d}\right)`$, the high-temperature (or infinitely massive) fixed point.
* $`\mu _{LT}^{}=0`$, the low-temperature fixed point.
We are more interested in nontrivial fixed points. But notice that in general, there are two generic ways fixed points can appear as $`N`$ or $`d`$ is varied :
splitting off from existing fixed points (bifurcation)
appearing in pairs in any region.
In the case (a), the signature is the approach to marginality of some operator representing a perturbation on an existing fixed point. The classic example is the Wilson-Fisher fixed point which bifurcates from the Gaussian as $`d`$ goes below four. The study of LPA (in the scalar case) yields no other kind of fixed point, this is why we consider first the vicinity of the Gaussian fixed point.
#### 3.2.1 The Gaussian fixed point
A study of the properties of the Gaussian fixed point may easily be realized by linearization of the flow equations in the vicinity of the origin.
In this linearization, all the LPA equations mentioned in section 3.1 reduce to a unique equation. Considering the derivative $`f(\phi )`$ of the potential $`V(\phi )`$ and a small deviation $`g(\phi )`$ to a fixed point solution $`f^{}(\phi )`$:
$$f(\phi )=f^{}(\phi )+g(\phi )$$
and choosing $`f^{}(\phi )0`$ (the Gaussian fixed point) the equations linearized in $`g`$ yields<sup>22</sup><sup>22</sup>22Up to some change of normalization for eq. (54) and eq. (65) what is authorized in the cases of the sharp cutoff and of the power law cutoff due to (evident, see section 3.1) reparametrization invariance. the unique equation :
$$\dot{g}=g^{\prime \prime }+(1\frac{d}{2})\phi g^{}+(1+\frac{d}{2})g$$
(69)
If one sets :
$$g(\phi ,t)=e^{\lambda t}\alpha h(\beta \phi )$$
with
$$\alpha =\frac{4}{d2}\text{,}\beta =\left(\frac{d2}{4}\right)^{\frac{1}{2}}$$
then (69) reads:
$$h^{\prime \prime }2\phi h^{}+2\frac{2+d2\lambda }{d2}h=0$$
(70)
##### Polynomial form of the potential
If we request the effective potential to be bounded by polynomials<sup>23</sup><sup>23</sup>23There are other possibilities, see below “Nonpolynomial…”. then eq. (70) identifies with the differential equation of Hermite’s polynomials of degree $`n=2k1`$ for the set of discrete values of $`\lambda `$ satisfying:
$$\frac{2+d2\lambda _k}{d2}=2k1\text{ }k=1,2,3,\mathrm{}$$
(71)
since $`f(\phi ,t)`$ is an odd function of $`\phi `$.
The same kind of considerations may be done for general $`N`$, in which case the Hermite polynomials are replaced by the Laguerre polynomials . Since the discussion is similar for all $`N`$, we limit ourselves here to a discussion of the simple case $`N=1`$.
###### Eigenvalues:
From (71), the eigenvalues are defined by:
$$\lambda _k=dk\left(d2\right)\text{ }k=1,2,3,\mathrm{}$$
(72)
then it follows that
* for $`d=4`$: $`\lambda _k=42k`$ $`k=1,2,3,\mathrm{}`$, there are two non-negative eigenvalues: $`\lambda _1=2`$ and $`\lambda _2=0`$
* for $`d=3`$: $`\lambda _k=3k`$ $`k=1,2,3,\mathrm{}`$, there are three non-negative eigenvalues: $`\lambda _1=2`$, $`\lambda _2=1`$ et $`\lambda _3=0`$
###### Eigenfunctions:
If we denote by $`\chi _k(\phi )`$ the eigenfunctions associated to the eigenvalue $`\lambda _k`$, it comes:
* $`\chi _1^+=\phi `$, $`\chi _2^+=\phi ^3\frac{3}{2}\phi `$, $`\chi _3^+=\phi ^55\phi ^3+\frac{15}{4}\phi `$, $`\mathrm{}`$, whatever the spatial dimensionality $`d`$. The upperscript “$`+`$” is just a reminder of the fact that the eigenfunctions are defined up to a global factor and thus the functions $`\chi _k^{}(\phi )=\chi _k^+(\phi )`$ are also eigenfunctions with the same eigenvalue $`\lambda _k`$. This seemingly harmless remark gains in importance after the following considerations.
To decide whether the marginal operator (associated with the eigenvalue equal to zero, i.e. $`\lambda _2`$ in four dimensions, or $`\lambda _3`$ in three dimensions) is relevant or irrelevant, one must go beyond the linear approximation. The analysis is presented in for $`d=4`$. If one considers a RG flow along $`\chi _2^+`$ such that $`g_2(\phi ,t)=c(t)\chi _2^+(\phi )`$, then one obtains, for small $`c`$: $`c(t)=c(0)\left[1Ac(0)t\right]`$ with $`A>0`$. Hence the marginal parameter decreases as $`t`$ grows. As is well known, in four dimensions the marginal parameter is irrelevant. However, if one considers the direction opposite to $`\chi _2^+`$ (i.e. $`\chi _2^{}`$) then the evolution corresponds to changing $`cc`$. This gives, for small values of $`c`$: $`c(t)=c(0)\left[1+A\left|c(0)\right|t\right]`$ and the parameter becomes relevant. The parameter $`c`$ is the renormalized $`\varphi ^4`$ coupling constant $`u_R`$ and it is known that in four dimensions the Gaussian fixed point is IR stable for $`u_R>0`$ but IR unstable for $`u_R<0`$ (if the corresponding action was positive for all $`\phi `$, one could say that the $`\varphi _4^4`$-field theory with a negative coupling is asymptotically free, see section 3.4.1).
For $`d`$ slightly smaller than four, $`\lambda _2`$ is positive and the Gaussian fixed point becomes IR unstable in the direction $`\chi _2^+`$ (and remains IR unstable along $`\chi _2^{}`$). The instability in the direction $`\chi _2^+`$ is responsible for the appearance of the famous Wilson-Fisher fixed point which remains the only known nontrivial fixed point until $`d`$ becomes smaller than 3 in which case it appears a second non trivial fixed point which bifurcates from the Gaussian fixed point. Any dimension $`d_k`$ corresponding to $`\lambda _k=0`$, is a border dimension below which a new fixed point appears by splitting off from the Gaussian fixed point. Eq. (72) gives<sup>24</sup><sup>24</sup>24This is a result already known from the $`ϵ`$-expansion framework .:
$$d_k=\frac{2k}{k1},k=2,3,\mathrm{},\mathrm{}$$
##### Nonpolynomial form of the potential
As pointed out by Halpern and Huang (see also ), there exist nonpolynomial eigenfunctions for the Gaussian fixed point. In four dimensions these nonpolynomial eigenpotentials have the asymptotic form $`\mathrm{exp}\left(c\phi ^2\right)`$ for large $`\phi `$ and provides the Gaussian fixed point with new relevant directions (with positive eigenvalues). From trivial, the scalar field theory in four dimensions would become physically non trivial due to asymptotic freedom and some effort have been made with a view to understand the physical implications of that finding .
Unfortunately, as stressed by Morris (see also ), the finding of Halpern and Huang implies a continuum of eigenvalues and this is opposite to the usual formulation of the RG theory as it is applied to field theory where the eigenvalues take on quantized values. Indeed, usually there are a finite number (preferably small) of relevant (renormalized) parameters, and it is precisely that property which is essential in the renormalization of field theory: if the number of relevant parameters is finite the theory is said renormalizable otherwise it is not. Let us emphasize that, there is no mathematical error in the work of Halpern and Huang (see the reply of Halpern and Huang which maintain their position except for the “line of fixed points”<sup>25</sup><sup>25</sup>25About the infinity of nontrivial fixed points, see also .), the key point is that the theory of “renormalization” for nonpolynomial potentials does not exist. We come back to this discussion in section 3.4.1 where we illustrate, among other notions, the notion of self-similarity which is rightly so dear to Morris .
#### 3.2.2 Non trivial fixed points
As one may see from the equations presented in section 3.1, the fixed point equation (67) is a second order non linear differential equation. Hence a solution would be parametrized by two arbitrary constants. One of these two constants may easily be determined if $`V^{}(\phi )`$ is expected to be an even function of $`\phi `$ \[O(1) symmetry\] then $`V^{}(0)=0`$ may be imposed<sup>26</sup><sup>26</sup>26Or if it is an odd function of $`\phi `$ then $`V^{\prime \prime }(0)=0`$ may be chosen as condition.. It remains one free parameter: a one-parameter family of (nontrivial) fixed points are solutions to the differential equation. But there is not an infinity of physically acceptable fixed points.
As first<sup>27</sup><sup>27</sup>27A discussion on the singular fixed point solutions in the case $`N=\mathrm{}`$ similar to that mentioned here for LPA, may be found in . indicated by Hasenfratz and Hasenfratz (private communication of H. Leutwyler) , studied in detail by Felder<sup>28</sup><sup>28</sup>28Who demonstrates that, for $`d=3`$, there is only one nontrivial fixed point. then by Filippov and Breus and by Morris , all but a finite number of the solutions in the family are singular at some $`\phi _c`$. By requiring the physical fixed point to be defined for all $`\phi `$ then the acceptable fixed points (if they exist) may be all found by adjusting one parameter in $`V(\phi )`$ (see fig. 1). For even fixed points, this parameter is generally chosen to be $`V^{\prime \prime }(0)`$ ($`=\sigma ^{}`$ in the following). For $`N=1`$ the situation is as follows:
* $`d4`$, no fixed point is found except for $`\sigma ^{}=0`$ (Gaussian fixed point).
* $`3d<4`$, one fixed point (the Wilson-Fisher fixed point ) is found for a nonzero value of $`\sigma ^{}`$ which depends on the equation considered. For $`d=3`$ one has:
$`\sigma ^{}=0.461533\mathrm{}`$ ($`0.4615413727\mathrm{}`$ in ) with eq. (54),
$`\sigma ^{}=0.228601293102\mathrm{}`$ in (or at $`V^{}(0)=0.0762099\mathrm{}`$ ) with eq. (59),
$`\sigma ^{}=0.5346\mathrm{}`$ with eq. (65).
* As indicated previously, a new nontrivial fixed point emanates from the origin (the Gaussian fixed point) below each dimensional threshold $`d_k=2k/(k1)`$, $`k=2,3,\mathrm{},\mathrm{}`$ .
We show in fig. 1, in the case $`d=3`$, how the physical fixed point is progressively discovered by adjusting $`\sigma =V^{\prime \prime }(0)`$ to $`\sigma ^{}`$ after several tries (shooting method). The knowledge of the behavior of the solution for large $`\phi `$ (obtained from the flow equation studied) greatly facilitates the numerical determination of $`\sigma ^{}`$ and of the fixed point solution $`V^{}(\phi )`$ (for example see ). For an indication on the numerical methods one can use, references are interesting.
#### 3.2.3 Critical exponents in three dimensions
Once the fixed point has been located, the first idea that generally occurs to someone is to calculate the critical exponents. There is only one exponent to calculate (e.g. $`\nu `$) since $`\eta =0`$. The other exponents are deduced from $`\nu `$ by the scaling relations (e.g. $`\gamma =2\nu `$)<sup>29</sup><sup>29</sup>29The right relation is $`\gamma =\nu (2\eta )`$.. The best way to calculate the exponents is to linearize the flow equation in the vicinity of the fixed point and to look at the eigenvalue problem. One obtains as in the case of the Gaussian fixed point a linear second order differential equation. For example with the Wilson (or Polchinski) version (59), setting $`V(\phi ,t)=V^{}+`$e$`{}_{}{}^{\lambda t}v(\phi )`$, one obtains the eigenvalue equation:
$$v^{\prime \prime }+\left[\left(1\frac{d}{2}\right)\phi 2V^{}\right]v^{}+\left(d\lambda \right)v=0$$
(73)
As Morris explains in the case of eq. (65) , “($`\mathrm{}`$) again one expects solutions to (73) labelled by two parameters, however by linearity one can choose $`v(0)=1`$ (arbitrary normalization of the eigenvectors) and by symmetry $`v^{}(0)=0`$ (or by asymmetry and linearity: $`v(0)=0`$ and $`v^{}(0)=1`$). Thus the solutions are unique, given $`\lambda `$. Now for large $`\phi `$, $`v(\phi )`$ is generically a superposition<sup>30</sup><sup>30</sup>30To obtain this behavior, use the large $`\phi `$ behavior of the fixed point potential $`V^{}(\phi )\frac{1}{2}\phi ^2`$ coming from eq.(59), see . of $`v_1\phi ^{2(d\lambda )/(d+2)}`$ and of $`v_2\mathrm{exp}\left(\frac{d+2}{4}\phi ^2\right)`$. Requiring zero coefficient for the latter restricts the allowed values of $`\lambda `$ to a discret set”.
The reason for which the exponential must be eliminated is the same as previously mentioned in section 3.2.1 to discard nonpolynomial forms of the potential.
For $`d=3`$, the Wilson-Fisher fixed point possesses just one positive eigenvalue $`\lambda _1`$ corresponding to the correlation length exponent ($`\nu =1/\lambda _1`$) and infinitely many negative eigenvalues. In the symmetric case, the less negative $`\lambda _2`$ corresponds to the first correction-to-scaling exponent $`\omega =\lambda _2`$ while $`\lambda _3`$ provides us with the second $`\omega _2=\lambda _3`$ and so on. In the asymmetric case, which is generally not considered (see however ), one may also associate the first negative eigenvalue $`\lambda _1^{\text{as}}`$ to the first non-symmetric correction-to-scaling exponent $`\omega _5=\lambda _1^{\text{as}}`$ (the subscript “5” refers to the $`\varphi ^5`$ interaction term in the action responsible for this kind of correction, see ).
Another way of numerically determining the (leading) eigenvalues is the shooting method. One chooses an initial (simple) action and tries to approach the fixed point (one parameter of the initial action must be finely adjusted). When the flow approaches very close to the fixed point, its rate of approach is controlled by $`\omega `$ (i.e. $`\lambda _2`$). The adjustment cannot be perfect and the flow ends up going away from the fixed point along the relevant direction with a rate controlled by $`\lambda _1=1/\nu `$.
The first determination of $`\nu `$ and $`\omega `$ has been made by Parola and Reatto from eq (54) \[in the course of constructing a unified theory of fluids, for a review see \]. They found<sup>31</sup><sup>31</sup>31In order to appreciate the quality of the estimates the reader may refer to the so-called best values given, for example, in . (for $`d=3`$ and $`N=1`$):
$$\nu =0.689,\omega =0.581$$
estimates which are compatible with the well known results of Hasenfratz and Hasenfratz obtained from eq. (55) using the shooting method:
$$\nu =0.687(1),\omega =0.595(1)$$
The error is only indicative of the numerical inaccuracy of solving the differential equation . Since then, the above results have been obtained several times. The first estimate of $`\omega _2`$ has been given in from the same equation and using again the shooting method:
$$\omega _22.8$$
No error was given due to the difficulty of approaching the fixed point along the second irrelevant direction (two parameters of the initial action must be adjusted ). More accurate estimations of this exponent may be found in .
We also mention an estimate of $`\omega _5`$ from the same eq. (54) :
$$\omega _51.69$$
In themselves the estimates of critical exponents in the local potential approximation do not present a great interest except as first order estimates in a systematic expansion see section 4.1. It is however interesting to notice that the LPA estimates are not unique but depend on the equation studied. Hence with the Legendre version (65) it comes :
$$\nu =0.6604,\omega =0.6285$$
which is closer to the “best” values. And the closest to the “best” are obtained from the Wilson (or Polchinski) version :
$$\nu =0.6496,\omega =0.6557$$
LPA estimates of exponents at various values of $`N`$ and $`d`$ have been published (see for example ). It is also worth mentioning that LPA gives the exact exponents up to $`O(\epsilon )`$ .
#### 3.2.4 Other dimensions $`2<d<3`$
If for $`2<d<3`$ multicritical fixed points appear at the dimensional thresholds $`d_k=2k/(k1)`$, $`k=2,3,\mathrm{},\mathrm{}`$, the Wilson-Fisher (critical) fixed point (once unstable) still exists and the same analysis as above for $`d=3`$ could have been done to estimate the critical exponents. However this kind of calculations have not been performed in the LPA despite some considerations relative to $`2<d<3`$ .
The special case of $`d=2`$ does not yield the infinite set of nonperturbative and multicritical fixed points expected following the conformal field theories but only periodic solutions corresponding to critical sine-Gordon models . This is due to having discarded the nonlocal contributions which are not small for $`d=2`$ ($`\eta =1/4`$ is not small) .
### 3.3 Truncations
One may try to find solutions to the ERGE within LPA by expanding the potential in powers of the constant field variable $`\phi `$. Although it is obviously not a convenient way of studying a non linear partial differential equation, this programme is interesting because the study of the ERGE necessarily requires some sort of truncation (or approximation). It is an opportunity to study the simplest truncation scheme in the simple configuration of the LPA in as much as there are complex systems (e.g. gauge theories) for which the truncation in the powers of the field seems inevitable .
Margaritis, Ódor and Patkós for arbitrary $`N`$ and Haagensen et al for $`2<d<4`$ have tried this kind of truncation on eq. (56). The idea is as follows. One expands $`V(\phi ,t)`$ in powers of $`\phi `$:
$$V(\phi ,t)=\underset{m=1}{\overset{\mathrm{}}{}}c_m(t)\phi ^m$$
and one reports within the flow equation to obtain, e.g. for the fixed point equation, an infinite system of equations for the coefficients $`c_i`$. That system may be truncated at order $`M`$ (i.e. $`c_i0`$ for $`i>M`$) to get a finite easily solvable set of equations (from which solutions may be obtained analytically ). By considering larger and larger values of $`M`$, one may expect to observe some convergence.
As one could think, the method does not generate very good results: an apparent new (but spurious) fixed point is found . Moreover, the estimation of critical exponents (associated to the Wilson-Fisher fixed point which, nevertheless, is identified) shows a poor oscillatory convergence . Indeed Morris has shown that this poor convergence is due to the proximity of a singularity in the complex plane of $`\phi `$.
Surprisingly, the truncation procedure considered just above works very well in the case of $`O(N)`$-symmetric systems. It appears that if one considers the variable $`\rho =\frac{1}{2}_\alpha \phi _\alpha ^2`$ and expands $`V(\rho ,t)`$ about the location $`\rho _0`$ of the minimum of the potential
$$\frac{\text{d}V(\rho ,t)}{\text{d}\rho }|_{\rho =\rho _0}=0$$
then one obtains an impressive apparent convergence toward the correct LPA value of the exponents (the method works also within the derivative expansion ). The convergence of the method has been studied in \[with eq. (54)\] and further in \[with the Legendre effective action (65)\] where it is shown that the truncation scheme associated to the expansion around the minimum of the potential (called co-moving scheme in ) actually does not converge but finally, at a certain large order, leads also to an oscillatory behavior.
We have seen that the estimates of the critical exponents in the LPA depend on the ERGE chosen (see section 3.2.3). It is thought however that for a given ERGE, the estimates should not depend on the form of the cutoff function (the dependence would occur only at next-to-leading order in the derivative expansion ). Nevertheless, the truncation in powers of the field may violate this “scheme independence” and affect the convergence of the truncation. It is the issue studied in with three smooth cutoff functions: hyperbolic tangent, exponential and power-law. An improvement of the convergence is proposed by adjustment of the smoothness of the cutoff function.
### 3.4 A textbook example
Despite its relative defect in precise quantitative predictions, the local potential approximation of the ERGE is more than simply the zeroth order of a systematic expansion in powers of derivatives (see section 4.1). In the first place it is a pedagogical example of the way infinitely many degrees of freedom are accounted for in RG theory. Almost all the characteristics of the RG theory are involved in the LPA. The only lacking features are related to phenomena highly correlated to the non local parts neglected in the approximation. For example in two dimensions, where $`\eta =\frac{1}{4}`$ is not particularly small, LPA is unable to display the expected fixed point structure . But, when $`\eta `$ is small (especially for $`d=4`$ and $`d=3`$), one expects the approximation to be qualitatively correct on all aspects of the RG theory. One may thus trust the results presented in section 3.2 on the search for nontrivial fixed points in four and three dimensions.
It is a matter of fact that much of studies on RG theory are limited to the vicinity of a fixed point. This is easily understood due to the universality of many quantities (exponents, amplitude ratios, scaled equation of state, etc$`\mathrm{}`$) that occurs there. However this limitation greatly curtails the possibilities that RG theory offers. The fixed points and their local properties (relevant and irrelevant directions) are not the only interesting aspects of the RG theory. Let us simply quote, as an example, the crossover phenomenon which reflects the competition between two fixed points. But what is worse than a simple limitation in the use of the theory, is the resulting misinterpretation of the theory. This is particularly true with respect to the definition of the continuum limit of field theory and its relation to the study of critical phenomena. Let us specify a bit this point (more details may be found in )
It is often expressed that the continuum limit of field theory is defined “at” a fixed point and that it is sufficient to look at its relevant directions to get the renormalized parameters, i.e. a simple linear study of the RG theory in the vicinity of the fixed point would be sufficient to define the continuum limit (see for example in ). This is not wrong but incomplete and, actually void of practical meaning. Indeed it is not enough emphasized (or understood) in the literature that, for example, although defined “at” the Gaussian fixed point, the field theoretic approach to critical phenomena is finally applied “at” the Wilson-Fisher fixed point which, in three dimensions, lies far away from the Gaussian fixed point. Then if the renormalized coupling of the $`\varphi _3^4`$-field theory was only defined by the linear properties of the RG theory in the vicinity of the Gaussian fixed point, one would certainly not be able to discover the nontrivial Wilson-Fisher fixed point by perturbative means.
Actually, the relevant directions of a fixed point provide us with exclusively the number and the nature of the renormalized parameters involved in the continuum limit. But the most important step of the continuum limit is the determination of the actual scale dependence (say, the beta functions) of those renormalized parameters. It is at this step that the recourse to RG theory actually makes sense: the attractive RG flow “that emanates out from the fixed point along the relevant direction ” results from the effect of infinitely many degrees of freedom and the problem of determining this flow is nonperturbative in essence. In the Wilson space $`𝒮`$ of the couplings $`\left\{u_n\right\}`$, the flow in the continuum runs along a submanifold (of dimension one if the fixed point has only one relevant direction) which is entirely plunged in $`𝒮`$. The writing of the corresponding action (the “perfect action” ) would require the specification of infinitely many conditions on the action. Because it is an hopeless task to find the “perfect action”<sup>32</sup><sup>32</sup>32However a truncated expression of the perfect action could be useful in actual studies of field theory defined on a lattice., one gives up all idea of a determination of the initial absolute scale dependence<sup>33</sup><sup>33</sup>33Called the functional form of the scale dependence in . and one has recourse to the process of fine tuning some parameter of the (bare) action in the vicinity of the (Gaussian) fixed point (or in the vicinity of a critical or tricritical etc.. surface in $`𝒮`$). This allows us to approach the ideal (or “perfect”) flow one is looking for and which runs along a trajectory tangent to a relevant direction. This flow is determined under a differential form (beta function) since the initial condition is not specified (the “perfect action” is not known) but it results from the effect of all degrees of freedom involved in the theory (by virtue of the relevant direction).
The simplicity of the ERGE in the LPA allows us to visualize the evolution of RG trajectories in the space $`𝒮`$ of the couplings $`\left\{u_n\right\}`$ and thus to illustrate and qualitatively discuss many aspects directly related to the nonperturbative character of the RG theory. In particular the approach to the continuum limit and the various domains of attraction of a fixed point.
#### 3.4.1 Renormalization group trajectories
Following , one considers an initial simple potential \[rather its derivative with respect to $`\phi `$, see (53)\], say<sup>34</sup><sup>34</sup>34The normalization of the couplings $`u_n`$ are here modified ($`u_nnu_n`$) compared to (52).:
$$f(\phi ,0)=u_2(0)\phi +u_4(0)\phi ^3+u_6(0)\phi ^5$$
(74)
corresponding to a point of coordinates $`(u_2(0),u_4(0),u_6(0),0,0,\mathrm{})`$ in $`𝒮`$, and after having numerically determined the associated solution $`f(\phi ,t)`$ of Eq. (55) at a varying “time” $`t`$, one concretely represents the RG trajectories (entirely plunged in $`𝒮`$) by numerically evaluating the derivatives of $`f`$ at the origin ($`\phi =0`$) corresponding to: $`u_2(t)`$, $`u_4(t)`$, $`u_6(t)`$ etc$`\mathrm{}`$
We then are able to visualize the actual RG trajectories by means of projections onto the planes $`\{u_2,u_4\}`$ or $`\{u_4,u_6\}`$ (for example) of the space $`𝒮`$.
For the sake of shortness we limit ourselves to a rapid presentation of figures. The reader is invited to read the original papers.
To approach the Wilson-Fisher fixed point in three dimensions (or $`3d<4`$) starting with (74), it is necessary to adjust the initial value of one coupling (e.g. $`u_2(0)`$) to a critical value $`u_2^c[u_4(0),u_6(0)]`$ (see fig. 2). This is because the Wilson-Fisher fixed point has one relevant direction (which must be thwarted). In such a case the fixed point controls the large distance properties of a critical system and the potential corresponding to (74) with $`u_2(0)=u_2^c[u_4(0),u_6(0)]`$ represents some physical system at criticality. The initial $`f(\phi ,0)`$ lies in the critical submanifold $`𝒮_\text{c}`$ which is locally orthogonal to the relevant eigendirection of the fixed point. As already mentioned in section 2.10.1, the renormalized trajectory (RT) T<sub>0</sub> that emerges from the Wilson-Fisher fixed point tangentially to the relevant direction allows to define a massive continuum limit<sup>35</sup><sup>35</sup>35The nontrivial continuum limit proposed by Wilson in for the so-called $`\varphi _3^6`$-field-theory which involves no coupling constant renormalization but a mass (and a wave function renormalization which cannot be evidenced in the present approximation). (see fig. 3).
In three dimensions, the Gaussian fixed point has two relevant directions. There a field theory involving two (renormalized) parameters (a mass-like and a $`\varphi ^4`$-like coupling) may be constructed. A purely massless field theory may also be constructed by choosing the relevant direction lying in the critical surface (corresponding to the eigenfunction $`\chi _2^+`$ of section 3.2.1). One obtains a one-parameter theory which interpolates between the Gaussian and the Wilson-Fisher fixed points (the renormalized submanifold T<sub>1</sub> of fig. 4). As already mentioned in section 2.10.1, this scale dependent massless theory contradicts the Morris view of massless theories as fixed point theories (thus scale invariant) . In fact a sensible massless (renormalized) scalar theory involves only a $`\varphi ^4`$-like coupling “constant”, say $`g`$, as parameter. But $`g`$ is not at all constant it is scale dependent and this is usually expressed via the beta function
$$\beta (g)=\mu \frac{\text{d}g}{\text{d}\mu }$$
in which $`\mu `$ is some momentum scale of reference and the function $`\beta (g)`$ is defined relatively to the flow running along the attractive submanifold T<sub>1</sub> (the slowest flow in the critical submanifold $`𝒮_\text{c}`$).
The fact that we get a unique scale dependent parameter (renormalizability) illustrates well the notion of “self-similarity” which means that the cooperative effect of the infinite number of degrees of freedom may be reproduced by means of a finite set of (effective or renormalized) scale dependent parameters (the system “looks like” the same when probed at different scales). In the case of the massless scalar theory, the number of renormalized parameters is equal to one: a coupling “constant”<sup>36</sup><sup>36</sup>36It is, perhaps, useful to mention here that the perturbative expression of the $`\beta `$-function (up to one loop order in and up to two loops in ) has been re-obtained from the ERGE once expanded with respect to the number of loops (perturbative study)..
Fig. 4 shows other RG trajectories (different to those attracted to T<sub>1</sub>), see the caption and for more details.
Less known are the RG trajectories drawn on fig. 5 in the sector $`u_4<0`$ . They correspond to:
1. The attractive submanifold T$`{}_{}{}^{\prime \prime }{}_{1}{}^{}`$ which emerges from the Gaussian fixed point and is tangent to the eigenfunction $`\chi _2^{}`$ of section 3.2.1 (symmetric to T<sub>1</sub> which emerges from the Gaussian fixed point tangentially to $`\chi _2^+`$ ). Along this submanifold the flows run toward larger and larger negative values of $`u_4`$. It is interesting to know that this submanifold still exists in four dimensions , it still corresponds to $`\chi _2^{}`$ and the associated eigenvalue is zero but the nonlinear analysis shows that the associated operator ($`\varphi ^4`$-like) is marginally relevant. Hence the $`\varphi _4^4`$ field theory with $`u_4<0`$ would be asymptotically free if the action had not the wrong sign for large values of $`\phi `$ (the negative sign of the $`\varphi _4^4`$ term which is the only dominant term of the “perfect action” in the vicinity of the Gaussian fixed point). It is worthwhile mentioning that this relevant direction of the Gaussian fixed point in four dimensions is different from those discovered by Halpern and Huang precisely because of the self-similarity displayed in the present case and which is a consequence of the discretization of the eigenvalues presented in section 3.2.1 in the case of polynomial interactions .
The attractive submanifold T$`{}_{}{}^{\prime \prime }{}_{1}{}^{}`$ is endless in the infrared direction (no nontrivial fixed point lies in the sector $`u_4<0`$, see section 3.2.2). It is customary to say in that case that T$`{}_{}{}^{\prime \prime }{}_{1}{}^{}`$ is associated to a first order transition. This is because without a fixed point the correlation length $`\xi `$ remains finite. However, although finite, the existence and the length of T$`{}_{}{}^{\prime \prime }{}_{1}{}^{}`$ suggest that $`\xi `$ may be very large. Following Zumbach in a study of the Stiefel nonlinear sigma model<sup>37</sup><sup>37</sup>37The Stiefel nonlinear sigma model is a generalization of the Heisenberg model with the field a real $`N\times P`$ matrix and the action is $`O(N)\times O(P)`$ invariant . within LPA , one may refer in the circumstances to “almost second order phase transition” or equivalently to “weakly first order phase transition” (see for supplementary details).
2. The “tricritical” attractive submanifold which approaches the Gaussian fixed point asymptotically tangentially to $`\chi _3^+`$ (the submanifold tangent to $`\chi _3^{}`$ also exists but is not drawn) necessitates the adjustment of two parameters in the initial action (because in three dimensions the Gaussian fixed point is twice unstable). See the figure caption for more comments and .
The same configuration displayed by fig. 5 has been obtained also by Tetradis and Litim while studying analytical solutions of an ERGE in the LPA for the $`O(N)`$-symmetric scalar theory in the large $`N`$ limit. But they were not able to determine “the region in parameter space which results in first order transitions. Fig. 5 shows this region for the scalar theory in three dimensions.
Fig. 6 shows RG trajectories in the critical surface for $`d=4`$ and $`u_4>0`$. The verification of this configuration was the main aim of the paper by Hasenfratz and Hasenfratz . As emphasized by Polchinski , in the infrared regime all the trajectories approach a submanifold of dimension one (on which runs the slowest RG flow ) before plunging into the Gaussian fixed point. This pseudo renormalized trajectory allows to make sense to the notion of effective field theory (a theory which only makes sense below some finite momentum scale). Because there is no fixed point other than Gaussian, it is not possible to adjust the initial action in such a way that the deviations between the actual RG trajectories and the ideal one dimensional submanifold be reduced to zero at arbitrary large momentum scales. Those irreducible deviations are responsible for the presence of the so-called UV renormalon singularities (for a review see ) in the perturbative construction of the $`\varphi _4^4`$ field theory. Indeed as explained in , the perturbative approach selects the ideal flow (the slowest) and sets a priori equal to zero all possible deviation (this is possible order by order in perturbation) without caring about the genuine physical momentum scale dependence (that requires the explicit reference to a relevant parameter relative to another fixed point to be well accounted for). Especially, it is argued in that the UV renormalon singularities would be absent in the $`\beta `$-function calculated within the minimal subtraction scheme of perturbation theory simply because in that case the scale of reference $`\mu `$ is completely artificial (has no relation with a genuine momentum scale except the dimension). See also , for a discussion of the UV renormalon singularities with the help of an ERGE.
A study of the attraction of RG flows to infrared stable submanifolds are also presented in .
#### 3.4.2 Absence of infrared divergences
It is known that the perturbation expansion of the massless scalar-field theory with $`d<4`$ involves infrared divergences. However, it has been shown that theory “develops by itself” an infinite number of non-perturbative terms that are adapted to make it well defined beyond perturbation. One may show that these purely nonperturbative terms are related to the critical parameter $`u_2^c`$ mentioned above . In the nonperturbative framework of the ERGE there is never infrared divergences and it is thus particularly well adapted to treat problems which are known to develop infrared singularities in the perturbative approach (e.g. Goldstone modes in the broken symmetry phase of $`O(N)`$ scalar theory and in general super renormalizable massless theories).
#### 3.4.3 Other illustrations
##### Limit $`N=\mathrm{}`$, large $`N`$
Exact results are accessible with LPA. The limit $`N=\mathrm{}`$ corresponds to the model of Berlin and Kac which can be exactly solved. Wegner and Houghton have shown that, in this limit their equation (29) is identical to the limit $`N=\mathrm{}`$ of the LPA. A first order non-linear differential equation is then obtained and studied in (see also ). Comparison with the exact model is completely satisfactory. More recently Breus and Filippov and then Comellas and Travesset have studied the large $`N`$ limit of the LPA for the Wilson (or Polchinski) ERGE (and for the Wegner-Houghton in ) and find also agreement with the exact results. However D’Attanasio and Morris have pointed out that the Large $`N`$ limit of the LPA for the Polchinski (or Wilson) ERGE is not exact although it may give correct results. See also in which analytical solutions in the Large $`N`$ limit of the LPA for the Legendre ERGE are obtained and discussed..
##### The RG flow is gradient flow
LPA allows to easily illustrate the property that the RG flow is gradient flow. This property is important because if the RG flow is gradient flow then only fixed points are allowed (limit cycles or more complicated behavior are excluded) and the eigenvalues of the linearized RG in the vicinity of a fixed point are real . The conditions for gradient flow are that the beta functions $`\beta _i\left(\left\{u\right\}\right)`$ \[The infinite set of differential renormalization group equations: $`\dot{u}_i=\beta _i\left(\left\{u\right\}\right)`$\] may be written in terms of a non-singular metric $`g_{ij}\left(\left\{u\right\}\right)`$ and a scalar function $`c\left(\left\{u\right\}\right)`$ :
$$\beta _i\left(\left\{u\right\}\right)=\underset{j}{}g_{ij}\left(\left\{u\right\}\right)\frac{c\left(\left\{u\right\}\right)}{u_j}$$
If $`g_{ij}\left(\left\{u\right\}\right)`$ is a positive-definite metric, the function $`c\left(\left\{u\right\}\right)`$ is monotonically decreasing along the RG flows:
$$\dot{c}=\underset{i}{}\beta _i\frac{c}{u_i}=\underset{i,j}{}g_{ij}\frac{c}{u_i}\frac{c}{u_j}0$$
Following Zumbach one may easily verify that the local potential approximation of the Wilson (or Polchinski) ERGE, written in terms of $`\mu (\phi ,t)=\mathrm{exp}\left(V(\phi ,t)\right)`$ \[eq. (68)\] may be expressed as a gradient flow:
$$g\left(\phi \right)\dot{\mu }=\frac{\delta \left[\mu \right]}{\delta \mu }$$
where
$`g\left(\phi \right)`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{4}}\left(d2\right)\phi ^2\right]`$
$`\left[\mu \right]`$ $`=`$ $`{\displaystyle \text{d}\phi g\left(\phi \right)\left\{\frac{1}{2}\mu ^{\prime \prime }+\frac{d}{4}\mu ^2\left(12\mathrm{ln}\mu \right)\right\}}`$
It has then been shown (see also ) that a c-function may be defined as ($`A`$ is a normalization factor)
$$c=\frac{1}{A}\mathrm{ln}\left(\frac{4}{d}\right)$$
which satisfies in any $`d`$ the two first properties of Zamolodchikov’s c-function and has a counting property which generalizes the third property.
Let us mention also that a c-function has been obtained in the framework of the truncation in powers of the field of section 3.3 by Haagensen et al and that Myers and Periwal have proposed a new form of the ERGE which is similar to a gradient flow.
##### Triviality bounds
It has been argued that an upper bound on the Higgs mass may be estimated from the only trivial character of the scalar field theory in four dimensions. The idea may be roughly illustrated by the following relation:
$$\frac{\mathrm{\Lambda }_{max}}{m}=_g^{\mathrm{}}\frac{dx}{\beta (x)}$$
(75)
in which $`g`$ is the (usual renormalized) $`\varphi ^4`$-coupling of the massive theory, $`m`$ is the mass parameter and $`\mathrm{\Lambda }_{max}`$ the maximum value that the momentum-scale of reference of the scalar theory can take on (it is associated with an infinite value of the coupling $`g`$ since no nontrivial fixed point exists —this is the consequence of triviality). Hence there is a finite relation between $`m`$ and $`\mathrm{\Lambda }_{max}`$.
In order to determine the (triviality) upper bound for the Higgs mass $`m_H`$ (which then replaces $`m`$), one usually refers to the ratio $`R=\frac{m_H}{m_W}`$ in which $`m_W`$ is the mass of the vector boson $`W`$ of the standard model of the electro-weak interaction . The ratio $`R`$ expresses as a function of both $`g`$ (the scalar coupling) and $`G`$ (the gauge coupling):
$$R=\frac{m_H}{m_W}=f(g,G)$$
(76)
It appears that, at fixed $`G`$, $`R`$ is an increasing function of $`g`$ . For example, at tree level it comes:
$$R^2=8\frac{g}{G^2}$$
(77)
Knowing $`G`$ and $`m_W`$ from experiments (usually one considers $`G^20.4`$ and $`m_W80`$ Gev, see for example), the calculation of $`f(g,G)`$ would thus allow us to estimate the triviality bound on $`m_H`$ from the following inequality:
$$Rf(\mathrm{},G)$$
(78)
However because $`g`$ becomes infinite at $`\mathrm{\Lambda }_{max}`$, the question of determining a bound on, say, the Higgs mass is highly a nonperturbative issue. Hasenfratz and Nager using the LPA of the Wegner-Houghton ERGE have shown how one can proceed to estimate that bound nonperturbatively (see also ).
##### Principle of naturalness.
Some authors have invoked a “concept of naturalness” to argue that fundamental scalar fields may not exist. Initiated by Wilson , this concept would require “the observable properties of a theory to be stable against minute variations of the fundamental parameters”.
A different concept of naturalness, brought up with a view to eliminate non asymptotically free field theories, would be that “the effective interactions ($`\mathrm{}`$) at a low energy scale $`\mu _1`$ should follow from the properties ($`\mathrm{}`$) at a much higher energy scale $`\mu _2`$ without the requirement that various different parameters at the energy scale $`\mu _2`$ match with an accuracy of the order of $`\frac{\mu _1}{\mu _2}`$. That would be unnatural. On the other hand, if at the energy scale $`\mu _2`$ some parameters would be very small, say $`\alpha (\mu _2)=O(\frac{\mu _1}{\mu _2})`$, then this may still be natural $`\mathrm{}`$.
Anyway, the two expressions have the same consequence for the scalar field theory which appears non natural. Let us illustrate this point with the help of the LPA.
With a view to make the scalar field theory in four dimensions (i.e. $`\varphi _4^4`$) non trivial, a non Gaussian fixed point is required. Assuming that a nontrivial fixed point exists in four dimensions, the procedure of construction of the resulting continuum limit “at” this fixed point would be similar to that described in fig. 3 for the purely massive field theory in three dimensions. We have seen that in order to approach the RT T<sub>0</sub>, at least one parameter ($`u_2(0)`$) of the (bare) initial action as to be finely tuned<sup>38</sup><sup>38</sup>38It is clear that the large number of digits in the determination of $`u_2^c`$ (e.g. $`0.299586913\mathrm{}`$ as indicated in the caption of fig. 2 or similarly in the determination of $`\sigma ^{}=0.228601293102\mathrm{}`$ in ) is not indicative of the accuracy in the determination of a physical parameter in the LPA but is required to get as close as possible to the critical surface (consistently within LPA). to a nonzero value ($`u_2^c`$). Moreover an infinitely small deviation from the actual $`u_2^c`$ results in a drastic change in the scale dependence of the effective renormalized parameter. This adjustment is unnatural: how can we justify the origin of numbers like $`u_2^c`$? This unnatural adjustment will be required each time one defines a continuum limit “at” a non trivial fixed point.
On the contrary, an asymptotically free field theory would appear natural because the adjustment of the initial action is made with respect to the Gaussian fixed point (the nonuniversal parameters like $`u_2(0)`$ are adjusted to zero!).
##### Dynamical generation of masses.
$`\mathrm{}`$ one essential element of this systematic theory (a satisfying synthesis of the theories of weak, electromagnetic and strong interactions) has remained obscure: we must take the mass of the leptons and quarks as input parameters, without any real idea of where they come from. $`\mathrm{}`$ the search for a truly natural theory of the quark masses must continue.
In order to give masses to the intermediate vector bosons while preserving symmetry, one has imagined the occurrence of a spontaneous symmetry breaking mechanism. In the standard model of electro-weak interaction the symmetry-breaking mechanism is associated with the introduction of a scalar field ($`\varphi _4^4`$) in the model and the vector bosons acquire masses via the Higgs mechanism . The conceptual difficulty with this model is that one has introduced a peculiar kind of interaction (the scalar field is self-interacting) which is not asymptotically free. The Higgs mechanism finally appears to be convenient in the range of energy scales over which the standard model seems to work but conceptually unnatural and not generalizable to higher energies ($`\varphi _4^4`$ does not make sense above some energy). The other possibility is that the symmetry-breaking mechanism occurs dynamically, that is to say without need for introducing scalars but simply because the gauge fields are interacting fields .
In the process of generating masses dynamically, the main interesting feature of a non-abelian-gauge-invariant field theory is not actually asymptotic freedom but its infrared diseases associated with the presence of IR “renormalons” in perturbative series (or equivalently the existence of a ghost in the infrared regime , for a review see ). It is very likely that those IR “renormalons” convey the lack of any infrared stable fixed point in the critical (i.e. massless) surface. This means that a scale dependent coupling constant $`G`$ of a purely massless (gauge invariant) theory is not defined below some momentum scale $`\mathrm{\Lambda }_{min}`$. It is thus expected that the appearance of massive particles below $`\mathrm{\Lambda }_{min}`$ can proceed from the existence of a symmetric (massless) theory at momentum-scales larger than $`\mathrm{\Lambda }_{min}`$ . Let us illustrate this point with the scalar theory and LPA.
The usual scalar field theory in three dimensions is well defined but does not present a great interest with respect to a mass generation because:
* either the theory is massive and the mass is a given parameter.
* or the theory is purely massless but defined as such at any (momentum) scale in the range $`]0,\mathrm{}[`$ (interpolation between two fixed points).
On the contrary, as shown in fig. 5, the massless theory becomes asymptotically free in the sector $`u_4<0`$ (the fact that the action has the wrong sign is not important for our illustrative purposes). But more importantly, there is no infrared stable fixed point to allow the scale dependence to be defined at any scale along this RT. Consequently, as close to the Gaussian fixed point as any trajectory would be initialized, the resulting trajectory will, after a finite “time”, end up going away from the critical surface, i.e. within the massive sector. Finally masses would have been generated from the momentum scale dependence of a purely massless theory.
This mechanism illustrated here for masses may also occur for any symmetry breaking parameter.
## 4 Further developments
### 4.1 Next-to-leading order in the derivative expansion
The LPA considered in the preceding sections is the zeroth order of a derivative expansion first proposed as a systematic expansion by Golner . To be fair, the first use of the derivative expansion (or rather the gradient expansion ) was by Myerson in conjunction with an expansion in powers of the field. A line of fixed points with $`\eta 0.045`$ was obtained.
The genuine derivative expansion is a functional power series expansion of the Wilson effective action in powers of momenta so that all powers of the field are included at each level of the approximation. The idea is to expand the action $`S[\varphi ;t]`$ in powers of momenta :
$$S[\varphi ;t]=S^{(0)}[\varphi ;t]+S^{(2)}[\varphi ;t]+\underset{i=1}{\overset{3}{}}S_i^{(4)}[\varphi ;t]+\mathrm{}$$
where
$`S_i^{(2k)}[\varphi ;t]`$ $`=`$ $`{\displaystyle \underset{n}{}}a_{in}^{(2k)}(t)H_{in}^{(2k)}[\varphi ]\text{,}`$
$`H_{in}^{(2k)}[\varphi ]`$ $`=`$ $`{\displaystyle _{q_1}}\mathrm{}{\displaystyle _{q_n}}h_i^{(2k)}(𝐪_1,\mathrm{},𝐪_n)\widehat{\delta }\left(𝐪_1+\mathrm{}+𝐪_n\right)\varphi _{q_1}\mathrm{}\varphi _{q_n}\text{,}`$
$`H_0^{(0)}[\varphi ]`$ $`=`$ $`\delta (0)`$
and the $`h_i^{(2k)}(𝐪_1,\mathrm{},𝐪_n)`$ are homogenous monomials in $`\left\{𝐪_j\right\}`$ of degree $`2k`$, with the index $`i`$ present when needed to keep track of degeneracies. Because of the momentum conserving $`\delta `$ function we have, for spatially isotropic systems, only one linearly independent functional of degree 2: $`h^{(2)}=𝐪_1𝐪_2`$, and three of degree 4: $`h_1^{(4)}=\left(𝐪_1𝐪_2\right)^2`$, $`h_2^{(4)}=\left(𝐪_1𝐪_2\right)\left(𝐪_1𝐪_3\right)`$, $`h_3^{(4)}=\left(𝐪_1𝐪_2\right)\left(𝐪_3𝐪_4\right)`$, since all powers of $`q_j^2`$ can be re-expressed in terms of powers of $`𝐪_i𝐪_j`$, $`ij`$. This is better seen in the position space where the expansion up to third order may be written as follows:
$`S[\varphi ]`$ $`=`$ $`{\displaystyle }d^dx\{V(\varphi ,t)+{\displaystyle \frac{1}{2}}Z(\varphi ,t)(_\mu \varphi )^2+H_1(\varphi ,t)(_\mu \varphi )^4`$
$`+H_2(\varphi ,t)(\mathrm{}\varphi )^2+H_3(\varphi ,t)(_\mu \varphi )^2(\mathrm{}\varphi )+\mathrm{}\}`$
on which expression the integrations by parts allow to easily identify the linearly dependent functionals (as previously the symbol $`\mathrm{}`$ stands for $`_\mu ^\mu `$).
It remains to substitute this expansion into the ERGE chosen among those described in section 2. To our knowledge the derivative expansion has only been explicitly written down up to the first order. This produces two coupled nonlinear partial differential equations for $`V`$ and $`Z`$.
For the sake of clarity we limit ourselves to a detailed discussion of the equations for the Polchinski version of the ERGE \[of section 2.5.2, eq. (34)\], the other forms of the ERGE<sup>39</sup><sup>39</sup>39The discussion of the Wilson formulation given by eq. (30) is very similar to that of Polchinski and will not be considered explicitly here (see for details). are considered in section 4.3.
The Polchinski version of the ERGE at first order of the derivative expansion yields the following coupled equations (see also ):
$`\dot{f}`$ $`=`$ $`2K^{}(0)ff^{}({\displaystyle K^{}})f^{\prime \prime }({\displaystyle p^2K^{}})Z^{}+{\displaystyle \frac{d+2\eta }{2}}f{\displaystyle \frac{d2+\eta }{2}}\phi f^{},`$
$`\dot{Z}`$ $`=`$ $`2K^{}(0)fZ^{}+4K^{}(0)f^{}Z+2K^{\prime \prime }(0)f^{}{}_{}{}^{2}({\displaystyle }K^{})Z^{\prime \prime }4K^1(0)K^{}(0)f^{}`$
$`\eta Z{\displaystyle \frac{d2+\eta }{2}}\phi Z^{},`$
with $`\phi \varphi _0`$ and $`f(\phi )V^{}(\phi )`$. As previously defined in section 2.5.2, $`K^{}`$ stands for d$`K(p^2)/`$d$`p^2`$ and $`K^{}_pK^{}(p^2)`$ etc….
It is convenient to perform the following rescalings
$$\phi \sqrt{K^{}}\phi ,f\frac{\sqrt{K^{}}}{K^{}(0)}f,ZK^1(0)Z;$$
so that,
$`\dot{f}`$ $`=`$ $`2ff^{}+f^{\prime \prime }+AZ^{}+{\displaystyle \frac{d+2\eta }{2}}f{\displaystyle \frac{d2+\eta }{2}}\phi f^{},`$ (79)
$`\dot{Z}`$ $`=`$ $`2fZ^{}4f^{}Z+2Bf^{}{}_{}{}^{2}+Z^{\prime \prime }+4f^{}\eta Z{\displaystyle \frac{d2+\eta }{2}}\phi Z^{},`$ (80)
where
$$A\frac{(K^{}(0))(p^2K^{})}{(K^{})},B\frac{K^{\prime \prime }(0)}{(K^{}(0))^2}.$$
Compared to , we have set $`K(0)=1`$. These conventions coincide also with those of .
Eqs. (79, 80) show that all cutoff (scheme) dependence at order $`p^2`$ is reduced to a two-parameter family $`(A,B)`$ while at zeroth order \[eq. (79) with $`Z^{}=0`$\] there is no explicit dependence. In general the scheme (cutoff) dependence can be absorbed into $`2k`$ parameters at $`k`$-th order in the derivative expansion . The set of eqs. (79, 80) has been considered first by Ball et al in with a view to study the scheme dependence of the estimates of critical exponents and reexamined by Comellas who emphasizes (following a remark by Morris ) the importance of the breaking of the reparametrization invariance in estimating the critical exponents (see also and section 2.2.4). Let us report on this important aspect of the eqs. (79, 80).
#### 4.1.1 Fixed points, $`\eta `$ and the breaking of the reparametrization invariance
The distribution of the fixed points of eqs. (79, 80) solution of $`\dot{f}^{}=\dot{Z}^{}=0`$ is identical to that of the leading order (LPA discussed in section 3.2) except for $`d=2`$ (see sections 3.2.4 and 4.1.3). Let us simply present the case of the Wilson-Fisher fixed point for $`d=3`$.
Following and in accordance with the discussion of section 3.2.2, to get the non trivial fixed point we impose the following boundary conditions:
$`f^{}(0)`$ $`=`$ $`0,`$ (81)
$`Z^{}(0)`$ $`=`$ $`0,`$ (82)
$`f^{}(\phi )`$ $``$ $`{\displaystyle \frac{2\eta }{2}}\phi +C\phi ^{\frac{d2+\eta }{d+2\eta }}+\mathrm{},\text{as }\phi \mathrm{}`$ (83)
$`Z^{}(\phi )`$ $``$ $`D+\mathrm{},\text{as }\phi \mathrm{}`$ (84)
where $`C`$ and $`D`$ are arbitrary constants. The first two conditions (81, 82) come from imposing $`Z_2`$-symmetry, while the last two come directly from the fixed point eqs. (79, 80), once we require the solutions to exist for the whole range $`0\phi <\mathrm{}`$.
Hence we have three free parameters ($`C`$, $`D`$, $`\eta `$) which are reduced to one after imposing eqs. (81, 82). The remaining arbitrary parameter, e.g. $`z=Z(0)`$, generates a line of (Wilson-Fisher) fixed points (one fixed point for each normalization $`z`$). In principle these fixed points are equivalent as a consequence of the reparameterization invariance (see section 2.2.4) and there is a corresponding unique value of $`\eta `$ (for any fixed point of the line). Consequently, if the reparametrization invariance was preserved one could get rid of the arbitrary parameter $`z`$ by setting it equal to 1. Unfortunately, due to the derivative expansion, this is not the case: eqs. (79, 80) violate the reparametrization invariance and the estimates of $`\eta `$ (and of $`\nu `$) depend on $`z`$.
In order to get the best estimates for $`\eta `$, one can adjust $`z`$ in such a way as to get an almost realized reparametrization invariance . The analysis is not simple due to the additional effects of the two cutoff parameters $`A`$ and $`B`$. Finally estimates of the critical and subcritical exponents are proposed ($`d=3`$ and $`N=1`$) :
$`\eta `$ $`=`$ $`0.042`$
$`\nu `$ $`=`$ $`0.622`$
$`\omega `$ $`=`$ $`0.754`$
It is interesting to compare these estimates with those obtained by Golner in from the Wilson version of the ERGE:
$`\eta `$ $`=`$ $`0.024\pm 0.007`$ (85)
$`\nu `$ $`=`$ $`0.617\pm 0.008`$ (86)
The equations are essentially similar in both cases and the difference in the estimations of $`\eta `$ surely originates from the way the cutoff is introduced and used in the Polchinski case.
Although those two sets of values are close to the best values (see footnote 31), the procedure which involves $`z`$ as adjustable parameter is less attractive than if $`\eta `$ was uniquely defined at each order of the derivative expansion. It is thus interesting to look for the conditions of preservation of the reparametrization invariance.
#### 4.1.2 Reparametrization invariance linearly realized and preserved
With a view to control the preservation of the reparametrization invariance, one may impose it evidently, i.e. linearly, via a particular choice of cutoff function and try to keep this realization through the derivative expansion . This is what has been done in for the Legendre version of the ERGE (see below). For the smooth cutoff version of the ERGE, the only acceptable cutoff function is power-law like (otherwise the cutoff should be sharp ). Unfortunately, for the Polchinski version, the symmetry is broken at finite order in the derivative expansion and the regulators do not regulate, at least not in a finite order in the derivative expansion . Now considering the Legendre version of the ERGE of section 2.6 is sufficient to overcome this difficulty .
##### The smooth cutoff Legendre version and the derivative expansion
By choosing a power-law cutoff function $`\stackrel{~}{C}(q^2)=q^{2k}`$ in eq. (39), one is sure that the derivative expansion will preserve the reparametrization invariance and that the exponent $`\eta `$ will be unambiguously defined.
Let us expand the Legendre (effective) action $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ as follows:
$$\mathrm{\Gamma }\left[\mathrm{\Phi }\right]=d^dx\left\{U(\phi ,t)+\frac{1}{2}Z(\phi ,t)(_\mu \mathrm{\Phi })^2\right\}$$
in which $`\phi `$ is independent on $`x`$.
For $`d=3`$ and $`k=1`$, the first order of the derivative expansion yields (after a long but straightforward computation) the following two coupled equations for $`U`$ and $`Z`$ :
$`\dot{U}`$ $`=`$ $`{\displaystyle \frac{1\eta /4}{\sqrt{Z}\sqrt{U^{\prime \prime }+2\sqrt{Z}}}}+3U{\displaystyle \frac{1}{2}}(1+\eta )\phi U^{}`$
$`\dot{Z}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(1+\eta )\phi Z^{}\eta Z+(1{\displaystyle \frac{\eta }{4}})\{{\displaystyle \frac{1}{48}}{\displaystyle \frac{24ZZ^{\prime \prime }19(Z^{})^2}{Z^{3/2}(U^{\prime \prime }+2\sqrt{Z})^{3/2}}}`$ (87)
$`{\displaystyle \frac{1}{48}}{\displaystyle \frac{58U^{\prime \prime \prime }Z^{}\sqrt{Z}+57(Z^{})^2+(Z^{\prime \prime \prime })^2Z}{Z(U^{\prime \prime }+2\sqrt{Z})^{5/2}}}+{\displaystyle \frac{5}{12}}{\displaystyle \frac{(U^{\prime \prime \prime })^2Z+2U^{\prime \prime \prime }Z^{}\sqrt{Z}+(Z^{})^2}{\sqrt{Z}(U^{\prime \prime }+2\sqrt{Z})^{7/2}}}\}`$
As expected, the search for a non trivial fixed point solution for these equations (a solution which is nonsingular up to $`\phi \mathrm{}`$) produces a unique solution with an unambiguously defined $`\eta `$ :
$$\eta =0.05393$$
(88)
The linearization about this fixed point yields the eigenvalues:
$`\nu `$ $`=`$ $`0.6181`$ (89)
$`\omega `$ $`=`$ $`0.8975`$ (90)
and also a zero eigenvalue $`\lambda =0`$ which corresponds to the redundant operator $`𝒪_1`$ \[eq. (24)\] responsible for the moving along the line of equivalent fixed points. This is, of course, an expected confirmation of the preservation of the reparametrization invariance.
A generalization of the above equations (87) to the $`O(N)`$ symmetric scalar field theory has been done by Morris and Turner in . There, estimates of $`\eta `$, $`\nu `$ and $`\omega `$ are provided for various values of $`N`$ and it is shown that the derivative expansion reproduces exactly known results at special values $`N=\mathrm{},2,4,\mathrm{}`$ and an interesting discussion on the numerical methods used is presented in their appendix.
##### The sharp cutoff Legendre version and the derivative expansion
The sharp cutoff is the other kind of regularization which allows a linear realization of the reparametrization invariance . As in the previous case of the power-law form of the cutoff function, the derivative expansion performed with the ERGE satisfied by the Wilson effective action $`S\left[\varphi \right]`$ with a sharp cutoff induces singularities which can be avoided by considering the Legendre transformed $`\mathrm{\Gamma }\left[\mathrm{\Phi }\right]`$ . But the Taylor expansion in the momenta must be replaced by an expansion in terms of homogeneous functions of momenta of integer degree (momentum-scale expansion). A systematic series of approximations — the $`O(p^M)`$ approximations — results .
Although not absolutely necessary, an additional expansion and truncation in powers of $`\phi `$ (avoiding the truncation of the potential) have been performed in due to the complexity of the equations<sup>40</sup><sup>40</sup>40Especially $`\phi ^8`$ terms and higher have been discarded in non-zero momentum pieces.. There, as in the previous case of the smooth cutoff, the zero eigenvalue corresponding to the redundant operator $`𝒪_1`$ is found. The estimates for the exponents, however, are worse than those obtained with smooth cutoff in \[see eqs.(88-90)\] presumably due to the truncation of the field dependence :
$`\eta `$ $`=`$ $`0.0660`$
$`\nu `$ $`=`$ $`0.612`$
$`\omega `$ $`=`$ $`0.91`$
The set of equations (87) together with the sharp cutoff version of the momentum expansion \[$`O(p^1`$)\] have also been studied in where, in particular, universal quantities other than the exponents (universal coupling ratios) have been estimated.
#### 4.1.3 Studies in two dimensions
In two dimensions it is expected that an infinite set of non-perturbative multicritical fixed points exists corresponding to the unitary minimal series of $`(p,p+1)`$ conformal field theories with $`p=3,4,\mathrm{},\mathrm{}`$ . As mentioned in section 3.2.4, this infinite set cannot be obtained at the level of LPA with which only periodic solutions could be obtained . Using the Legendre ERGE at first order of the derivative expansion with a power law cutoff \[the equations are obtained similarly to (87) but for $`d=2`$\], Morris in (see also ) has found the first ten fixed points (and only these) and computed the corresponding critical exponents (and other quantities). The comparison with the exact results of the conformal field theory is satisfactory (in consideration of the low — the lowest — order of approximation). A similar study has been done using the Polchinski ERGE (at first order of the derivative expansion) by Kubyshin et al using the same iteration technique as in .
### 4.2 A field theorist’s self-consistent approach
There is an efficient short cut for obtaining the ERGE satisfied by the (Legendre) effective action. It is based on the observation that this (exact) equation \[see (40 or 41)\] may be obtained from the one loop (unregularized, thus formal) expression of the effective action, which reads (up to a field independent term within the logarithm):
$$\mathrm{\Gamma }\left[\mathrm{\Phi }\right]=S\left[\mathrm{\Phi }\right]+\frac{1}{2}\text{tr}\mathrm{ln}\left(\frac{\delta ^2S}{\delta \varphi \delta \varphi }|_{\varphi =\mathrm{\Phi }}\right)+\text{higher loop-order,}$$
(91)
by using the following practical rules:
1. add the infrared cutoff function $`C(p,\mathrm{\Lambda })`$ of (38) to the action $`S`$, eq. (91) then becomes:
$$\mathrm{\Gamma }\left[\mathrm{\Phi }\right]=\frac{1}{2}_p\mathrm{\Phi }_p\mathrm{\Phi }_pC^1(p,\mathrm{\Lambda })+S\left[\mathrm{\Phi }\right]+\frac{1}{2}\text{tr}\mathrm{ln}\left(C^1+\frac{\delta ^2S}{\delta \varphi \delta \varphi }\right)|_{\varphi =\mathrm{\Phi }}+\mathrm{}$$
2. redefine $`\stackrel{~}{\mathrm{\Gamma }}\left[\mathrm{\Phi }\right]=\mathrm{\Gamma }\left[\mathrm{\Phi }\right]\frac{1}{2}_p\mathrm{\Phi }_p\mathrm{\Phi }_pC^1(p,\mathrm{\Lambda })`$, then:
$$\stackrel{~}{\mathrm{\Gamma }}\left[\mathrm{\Phi }\right]=S\left[\mathrm{\Phi }\right]+\frac{1}{2}\text{tr}\mathrm{ln}\left(C^1+\frac{\delta ^2S}{\delta \varphi \delta \varphi }\right)|_{\varphi =\mathrm{\Phi }}+\mathrm{}$$
3. perform the derivative with respect to $`\mathrm{\Lambda }`$, (only the cutoff function is concerned) and forget about the higher loop contributions:
$$_t\stackrel{~}{\mathrm{\Gamma }}=\frac{1}{2}\text{tr}\left[\frac{1}{C}\mathrm{\Lambda }\frac{C}{\mathrm{\Lambda }}\left(1+C\frac{\delta ^2S[\mathrm{\Phi }]}{\delta \mathrm{\Phi }\delta \mathrm{\Phi }}\right)^1\right]$$
4. replace<sup>41</sup><sup>41</sup>41This step is often referred to as the “renormalization group improvement” of the one loop effective action. the action $`S`$, in the right hand side of the latter equation, by the effective action $`\stackrel{~}{\mathrm{\Gamma }}`$ to get eqs. (40, 41), the dilatation part $`𝒢_{\text{dil}}\stackrel{~}{\mathrm{\Gamma }}`$ being obtained from usual (engineering) dimensional considerations<sup>42</sup><sup>42</sup>42But do not forget to introduce the anomalous dimension of the field in order to get an eventual nontrivial fixed point..
It is noteworthy that the above rules have been heuristically first used to obtain the local potential approximation of the ERGE for the (Legendre) effective action. However the main interest of the above considerations is that they allow introducing the (infra-red) cutoff function independently of $`S`$, via the so-called “proper time” (or “heat kernel” or “operator”) regularization . This kind of regularization is introduced at the level of eq. (91) via the general identity:
$$\text{tr}\mathrm{ln}\left(\frac{A}{B}\right)=_0^{\mathrm{}}\frac{\text{d}s}{s}\text{tr}\left(\text{e}^{sA}\text{e}^{sB}\right)$$
Forgetting again about the field-independent part (and, momentaneously, about the ultra-violet regularization needed for $`s0`$) one introduces an infrared cutoff function $`F_\mathrm{\Lambda }\left(s\right)`$ within the proper time integral representation of the logarithm of $`A=\frac{\delta ^2S}{\delta \varphi \delta \varphi }|_{\varphi =\mathrm{\Phi }}`$:
$$\frac{1}{2}\text{tr}\mathrm{ln}A\frac{1}{2}_0^{\mathrm{}}\frac{\text{d}s}{s}F_\mathrm{\Lambda }\left(s\right)\text{tr e}^{sA}$$
The function $`F_\mathrm{\Lambda }\left(s\right)`$ must tend to zero sufficiently rapidly for large values of $`s`$ in order to suppress the small momentum modes and should be equal to 1 for $`\mathrm{\Lambda }=0`$.
Then following the rules 3-4 above applied on $`\mathrm{\Gamma }`$ (i.e., not on $`\stackrel{~}{\mathrm{\Gamma }}`$), one obtains a new kind of ERGE<sup>43</sup><sup>43</sup>43Notice that, because one performs a derivative with respect to $`\mathrm{\Lambda }`$, the essential contribution to $`_t\mathrm{\Gamma }`$ comes from the integration over a small range of values of $`s`$ (corresponding to the rapid decreasing of $`F_\mathrm{\Lambda }\left(s\right)`$), hence an ultraviolet regularization is not needed provided that the resulting RG equation be finite. for the effective action :
$$_t\mathrm{\Gamma }=\frac{1}{2}_0^{\mathrm{}}\frac{\text{d}s}{s}\mathrm{\Lambda }\frac{F_\mathrm{\Lambda }\left(s\right)}{\mathrm{\Lambda }}\mathrm{exp}\left[s\frac{\delta ^2\mathrm{\Gamma }}{\delta \mathrm{\Phi }\delta \mathrm{\Phi }}\right]$$
(92)
There are apparently two advantages of using this kind of ERGE:
* the regularization preserves the symmetry of the action
* the derivative expansion is slightly easier to perform than in the conventional approach and one may preserve the reparametrization invariance .
In , Bonanno and Zappalà have considered the next-to-leading order of the derivative expansion of (92) (while in only a pseudo derivative expansion, in which the wave-function renormalization function $`Z(\varphi ,t)`$ is field-independent<sup>44</sup><sup>44</sup>44See below., was used). They have chosen $`F_\mathrm{\Lambda }\left(s\right)`$ in such a way that the integro-differential character of the ERGE disappears and they have tested the preservation of the reparametrization invariance. Moreover a scheme dependence parameter, related to the cutoff width, is at hand in this framework which, presumably, will allow someone to look at the best possible convergence of the derivative expansion when higher orders will be considered.
Another kind of regularization related to this “self-consistent” approach should be mentioned here. It consists in introducing the cutoff function in eq. (91) in-between the momentum integration \[expressing the trace\] and the logarithm. This procedure has been considered in at the level of the local potential approximation. However the ERGE keeps its integro-differential character and the study of has then been conducted within the constraining polynomial expansion method of section 3.3.
### 4.3 Other studies up to first order of the derivative expansion
In this section we mention studies of the derivative expansion which, although interesting, do not consider explicitly the reparametrization invariance.
Filippov and Radievskii have obtained a set of two coupled equations that look like eqs. (79, 80) but their numerical studies were based on an approximation which consists in neglecting the term corresponding to $`AZ^{}`$ in (79). They, nevertheless, present interesting estimates of critical exponents for several values of $`d`$ in the range $`]2,3.5]`$.
As already mentioned, Ball et al have studied the “scheme” dependence with the help of the Polchinski ERGE at first order in the derivative expansion (without considering explicitly the breaking of the reparametrization invariance). They have used an interesting simple iteration procedure to determined the fixed point.
Bonanno et al have presented a sharp version of the coupled differential equation for $`V`$ and $`Z`$ which however yields a negative value of $`\eta `$ in three dimensions. The authors claim that this failure is not to be searched in an intrinsic weakness of the sharp cutoff. A previous attempt had been done with the sharp cutoff version and a two loop perturbative anomalous dimension was obtained by means of a polynomial truncation in the field dependence . Re-obtention of two loop results from the derivative expansion of the ERGE satisfied by the effective (Legendre) action (with smooth cutoff) are also described in .
###### Pseudo derivative expansion
Tetradis and Wetterich have initiated an original strategy to obtain systematic accurate estimates on, say, critical quantities already from the lowest order of the derivative expansion. The idea is based on the smallness of $`\eta `$ and may be roughly described as follows. At lowest order of the derivative expansion (LPA), one assumes that $`Z(0,t)`$ already depends on $`t`$. One then determines an approximate $`t`$-dependence (assuming $`\eta `$ is small) from the momentum dependence of the exact propagator $`\mathrm{\Gamma }^{(2)}`$. Hence $`\eta `$ is not equal to zero even at the lowest order of the derivative expansion and this yields an “improved” LPA. The next order would amount to consider an explicit $`\phi `$-dependent $`Z(\phi ,t)`$ and the following order higher derivatives of the field in the action. This is not a genuine derivative expansion and it does not account for the reparametrization invariance. Nevertheless the approach seems efficient considering the estimates obtained at the leading order of that pseudo derivative expansion. Let us first quote, for $`d=3`$ and $`N=1`$, the results found with the supplementary help of a truncation in powers of the field associated to an expansion around the minimum of the potential : $`\nu =0.638`$, $`\eta =0.045`$, $`\gamma =1.247`$, $`\beta =0.333`$ and without truncation in the field dependence : $`\nu =0.643`$, $`\eta =0.044`$, $`\gamma =1.258`$, $`\beta =0.336`$, $`\delta =4.75`$. In this latter work, the scaled equation of state has been calculated using this pseudo derivative expansion. For more details on this approach see the review by Berges et al in this volume.
A study for $`d=2`$ has also been achieved following the spirit of (i.e. with a truncation) with a view to discuss the Kosterlitz-Thouless phase transition . The aim of the authors was to show the power of the ERGE compared to the perturbative approach (due to IR singularities). It is amazing to notice the excellent estimation of $`\eta `$ ($`0`$.$`24`$ instead of $`\frac{1}{4}`$) obtained in this work knowing that $`\eta `$ was assumed to be small and that the truncation was crude. As indicated by the authors, this result may be accidental. An extension of this work may be found in .
### 4.4 Convergence of the derivative expansion?
Comparing the estimates of the critical exponents obtained at first order of the derivative expansion to that of, e.g., the $`\epsilon `$-expansion, one can easily see that the derivative expansion is potentially much more effective than the perturbative (field theoretical) approach. But, to date, it is not known whether it converges or not. Morris and Tighe have considered this question at one and two loop orders for different cases of regularization (cutoff) functions and for either the Wilson (— Polchinski) or Legendre effective action. It is found that the Legendre flow equation converges at one and two loops: slowly with sharp cutoff (as a momentum-scale expansion), and rapidly in the case of a smooth exponential cutoff (but, in this latter case, the reparametrization invariance is not satisfied, see above). The Wilson (— Polchinski) version and the Legendre flow equation with power law cutoff function do not converge.
It is possible that the derivative expansion gives rise to asymptotic series which would be Borel summable. This is deduced from the knowledge of an exact solution for the effective potential for QED<sub>2+1</sub> in a particular inhomogeneous external magnetic field, from which it has been shown that the derivative expansion (known at any order) is a divergent but Borel summable asymptotic series .
The annoying perspective that the derivative expansion does not converge has prevailed on Golner to look for a method of successive approximations for the ERGE that is not based on power series expansion .
### 4.5 Other models, other ERGE’s, other studies…
Up to this point, we have presented in some details various aspects (derivation, invariances, approximations, truncations, calculations)<sup>45</sup><sup>45</sup>45Let us quote, in addition, the issue of scheme dependence, already mentioned in the text, several aspects of which are considered in . of the ERGE for scalar systems. Since these issues are also encountered for more complex systems (but with, potentially, a more interesting physical content), in this section we limit ourselves to mentioning the existence of studies based on the ERGE relative to models different from the pure scalar theory<sup>46</sup><sup>46</sup>46One may also refer to a recent review by Aoki .. Most often these models involve more structure due to supplementary internal degrees of freedom. Formally, the master equations keep essentially the same general forms as described in section 2 \[owing to the trace symbol as used in (40)\]. The equations involved in the studies actually show their differences when approximations are effectively considered. The studies are characterized by the action $`S`$ considered, the cutoff function chosen and the approximation applied on the ERGE. In consideration of the large number of publications and the variety of models studied, we choose to classify them according to the increasing degree of complexity of the model with respect to the field: scalar (or vector), spinor and gauge field.
###### Other ERGE’s involving pure scalars (or vectors)
We have already mentioned the Stiefel model studied in the LPA in let us quote also the Boson systems , the nucleation and spinodal decomposition , the interface unbinding transitions arising in wetting phenomena , the roughening transition , transitions in magnets (non-collinear spin ordering, frustrated) studied in , disordered systems , the quantum tunnelling effect , the well developed turbulence and even the one-quantum-particle system . More developed are the numerous studies of scalar theories at finite temperature . Reviews on the finite temperature framework may be found in though they cover more than scalars.
###### Other ERGE’s
It exists some studies involving pure spinors and also some mixing scalars and spinors , reviews may be found in .
In addition, a rich literature on ERGE deals with systems in presence of gauge fields: pure gauge fields , gauge fields with scalars and with spins , supersymmetric gauge fields and gravity . Reviews on this theme are listed in .
Despite the great number of studies done up to now on the ERGE, its systematic use in nonperturbative calculations and in describing nonuniversalities is still in its infancy. A better mastery of invariances within the truncation procedure, the extension of series (this has required some time in perturbation theory), the consideration of more complex and realistic models with a view to obtain estimates of useful physical quantities will necessitate much more investigations in the future.
Figure captions
1. Three solutions of the sharp cutoff fixed point equation for $`f(\varphi _0)=V^{}(\varphi _0)`$ and $`d=3`$ \[eq. (55) with $`\dot{f}=0`$\]. All (here two) but one ($`\sigma ^{}=0.4615337\mathrm{}`$) of the solutions are singular at some (not fixed) $`\phi _c`$. The parameter $`\sigma =V^{\prime \prime }(0)`$ is adjusted to $`\sigma ^{}`$ by requiring the physical fixed point to be defined for all $`\varphi _0`$ (in the text $`\phi `$ stands for $`\varphi _0`$).
2. Determination by the shooting method of the initial critical value $`u_2^c=0.299586913\mathrm{}`$ corresponding to the initial values $`u_4(0)=3`$ and $`u_n(0)=0`$ for $`n>4`$ \[from eq. (55) with $`d=3`$\]. Open circles indicate the initial points chosen in the canonical surface (representing simple actions) of $`𝒮`$. The illustration is made via projections onto the plane $`[u_2,u_4]`$. The determination of $`u_2^c`$ is made by iterations (shooting method) according to increasing labels. Arrows indicate the infrared direction (decreasing of the momentum-scale of reference). The RG trajectories follow two opposite directions according to whether $`u_2(0)>u_2^c`$ (labels 1, 3, 5) or $`u_2(0)<u_2^c`$ (labels 2, 4, 6). The Wilson-Fisher (once infrared unstable) fixed point (full circle) is only reached when $`u_2(0)=u_2^c`$ (dashed curve). The corresponding RG trajectory lies in the critical surface $`𝒮_c`$ of codimension 1.
3. Illustration of the simplest nonperturbative continuum limit in three dimensions \[from eq. (55) with $`d=3`$\]. Approach to the purely massive “renormalized trajectory” $`T_0`$ (dot-dashed curve) by RG trajectories initialized at $`u_4(0)=3`$ and $`u_n(0)=0`$ for $`n>4`$ and $`(u_2(0)u_2^c)0^+`$ (open circles). The trajectories drawn correspond to $`\mathrm{log}(u_2(0)u_2^c)=1,2,3,4,5,6`$. When $`u_2(0)=u_2^c`$ the trajectories do not leave the critical surface and approach the Wilson-Fisher fixed point (full circle), as in figure 2. But, “moving a little bit away from the critical manifold, the trajectory of the RG will to begin with, move towards the fixed point, but then shoot away along \[…\] the relevant direction towards the so-called high temperature fixed point (see text, section 2.10.1 and ).
4. Projection onto the plane $`(u_4,u_6)`$ of some remarkable RG trajectories for $`u_4(0)>0`$ \[from eq. (55) with $`d=3`$\]. Full lines represent trajectories on the critical surface $`𝒮_c`$. The arrows indicate the directions of the RG flows on the trajectories. The submanifold T<sub>1</sub> of one dimension to which are attracted the trajectories with small values of $`u_4(0)`$ and which links the Gaussian fixed point to the Wilson-Fisher fixed point corresponds to the renormalized trajectory on which is defined the continuum limit of the massless field theory in three dimensions. For larger values of $`u_4(0)`$ the RG trajectories approach the Wilson-Fisher fixed point from the opposite side, they correspond to the Ising model. The dotted line T<sub>2</sub> plunging into the Wilson-Fisher fixed point does not lie on $`𝒮_c`$ but represents a RG trajectory approaching the Wilson-Fisher fixed point along the second less irrelevant direction (lying in a space of codimension 2). The corresponding critical behavior is characterized by the absence of the first kind of correction to scaling (that corresponding to the exponent $`\omega `$), it would be representative of some Ising models with spin $`s=1/2`$. The two trajectories that leave the Wilson-Fisher fixed point (dashed lines) correspond to the unique relevant eigendirection (with two ways, due to the arbitrary normalization, associated with the two phases of the critical point, they also correspond to two massive RT’s). The open circles represent initial simple actions.
5. Projection onto the plane $`\{u_2,u_4\}`$ of some remarkable RG trajectories for $`u_4(0)<0`$ \[from eq. (55) with $`d=3`$\]. Black circles represent the Gaussian and Wilson-Fisher fixed points. The arrows indicate the directions of the RG flows on the trajectories. The ideal trajectory T<sub>1</sub> (dot line) which interpolates between the two fixed points represents the RT corresponding to the so-called $`\varphi _3^4`$ renormalized field theory in three dimensions (usual RT for $`u_4>0`$). White circles represent the projections onto the plane of initial critical actions. For $`u_4(0)>0`$, the effective actions (e.g. initialized at B’) run toward the Wilson-Fisher fixed point asymptotically along the usual RT. Instead, for $`u_4(0)<0`$ and according to the initial values of the parameters of higher order ($`u_6`$, $`u_8`$, etc.), the RG trajectories either (A) meet an endless RT emerging from the Gaussian fixed point T$`{}_{}{}^{\prime \prime }{}_{1}{}^{}`$ (dashed curve) and lying entirely in the sector $`u_4<0`$ or (B) meet the usual RT to reach the Wilson-Fisher fixed point. The frontier which separates these two very different cases (A and B) corresponds to initial actions lying on the tri-critical subspace (white square C) that are sources of RG trajectories flowing toward the Gaussian fixed point asymptotically along the tricritical (pseudo) RT. Notice that the coincidence of the initial point B with the RG trajectory starting at point A is not real (it is accidentally due to the projection onto a plane of the trajectories lying in a space of infinite dimension). See text for a discussion and .
6. RG trajectories on the critical surface $`𝒮_c`$ obtained from integration of eq. (55) with $`d=4`$ (projection onto the plane $`(u_4,u_6)`$). Open circles indicate the initial points chosen on the canonical surface of $`𝒮_c`$ (of codimension 1).The two lines which come from the upper side of the figure are RG trajectories initialized at $`u_4(0)=20`$ and $`u_4(0)=40`$ respectively. The arrows indicate the infrared direction. The trajectories are attracted to a submanifold of dimension one before plunging into the Gaussian fixed point. This pseudo renormalized trajectory (it has no well defined beginning) allows to make sense to the notion of effective (here massless) field theory. Strictly speaking, the continuum limit does not exist due to the lack of another (nontrivial fixed point) which would allow the scale dependence (of the renormalized parameter along the RT) to be defined in the whole range of scale $`]0,\mathrm{}[`$. See text for a discussion (from ). |
warning/0002/math0002230.html | ar5iv | text | # Gauge transformations on locally trivial quantum principal fibre bundles
## 1 Gauge transformations
In gauge transformation on a QPFB $`𝒫`$ are defined as automorphisms $`\alpha :𝒫𝒫`$ fulfilling the conditions
$`(\alpha id)\mathrm{\Delta }_𝒫`$ $`=`$ $`\mathrm{\Delta }_𝒫\alpha `$ (1)
$`\alpha \iota `$ $`=`$ $`\iota .`$ (2)
On trivial QPFB gauge transformations are in one to one correspondence with homomorphisms $`\tau _\alpha :HB`$ fulfilling
$`\tau _\alpha (1)`$ $`=`$ $`1`$
$`\tau _\alpha (h)a`$ $`=`$ $`a\tau _\alpha (h),hHaB`$
such that
$$\alpha (ah)=a\tau _\alpha (h_1)h_2.$$
To get nonclassical gauge transformations one needs a more general definition. First, let us define gauge transformation on trivial QPFB.
###### Definition 1
Let $`BH`$ be a trivial QPFB. A left (right) gauge transformation on $`BH`$ is a left (right) $`(B1)`$-module isomorphism $`\alpha :BHBH`$ satisfying
$`(\alpha id)(id\mathrm{\Delta })`$ $`=`$ $`(id\mathrm{\Delta })\alpha `$ (3)
$`\alpha (id1)`$ $`=`$ $`(id1).`$ (4)
Remark: Obviously, for every gauge transformation $`\alpha `$, the inverse $`\alpha ^1`$ is also a gauge transformation.
In the sequel, our standard notation will be $`\alpha _l`$ for left and $`\alpha _r`$ for right gauge transformations. The notation $`\alpha _{l,r}`$ will be used if something is true for both left and right gauge transformations.
###### Proposition 1
Left (right) gauge transformations $`\alpha _{l,r}`$ on a trivial bundle $`BH`$ are in one to one correspondence to linear maps $`\tau _{\alpha _{l,r}}:HB`$ satisfying
$`\tau _{\alpha _{l,r}}(1)`$ $`=`$ $`1`$
$`{\displaystyle \tau _{\alpha _l^1}(h_1)\tau _{\alpha _l}(h_2)}`$ $`=`$ $`{\displaystyle \tau _{\alpha _l}(h_1)\tau _{\alpha _l^1}(h_2)}=\epsilon (h)1`$
$`{\displaystyle \tau _{\alpha _r^1}(h_2)\tau _{\alpha _r}(h_1)}`$ $`=`$ $`{\displaystyle \tau _{\alpha _r}(h_2)\tau _{\alpha _r^1}(h_1)}=\epsilon (h)1`$
such that
$`\alpha _l(ah)`$ $`=`$ $`{\displaystyle a\tau _{\alpha _l}(h_1)h_2}`$ (5)
$`\alpha _r(ah)`$ $`=`$ $`{\displaystyle \tau _{\alpha _r}(h_1)ah_2}.`$ (6)
Proof: We define
$$\tau _{\alpha _{l,r}}(h)=(id\epsilon )\alpha _{l,r}(1h).$$
(7)
Because of formula (3) and $`(id\epsilon id)(id\mathrm{\Delta })=id`$ we obtain for left gauge transformations
$`\alpha _l(ah)`$ $`=`$ $`(id\epsilon id)(id\mathrm{\Delta })\alpha _l(ah)`$
$`=`$ $`(a1)(id\epsilon id)(\alpha _lid)(id\mathrm{\Delta })(1h)`$
$`=`$ $`(a1){\displaystyle (\tau _{\alpha _l}(h_1)h_2)}`$
$`=`$ $`{\displaystyle (a\tau _{\alpha _l}(h_1)h_2)}`$
and for right gauge transformations
$`\alpha _r(ah)`$ $`=`$ $`(id\epsilon id)(id\mathrm{\Delta })\alpha _r(ah)`$
$`=`$ $`(id\epsilon id)(\alpha _rid)(id\mathrm{\Delta })(1h)(a1)`$
$`=`$ $`{\displaystyle (\tau _{\alpha _r}(h_1)h_2)(a1)}`$
$`=`$ $`{\displaystyle (\tau _{\alpha _l}(h_1)ah_2)}.`$
Since the gauge transformations are left (right) $`(B1)`$-module isomorphisms, the properties claimed for $`\tau _{\alpha _{l,r}}`$ are easily verified. We leave the other direction of the proof to the reader. $`\mathrm{}`$
###### Definition 2
Let $`𝒫`$ be a locally trivial QPFB with local trivializations $`\chi _i:𝒫B_iH`$.
A left (right) gauge transformation on $`𝒫`$ is a left $`\iota (B)`$-module isomorphism $`\alpha _l,r:𝒫𝒫`$ such that there exists a family $`(\alpha _{l,r_i})_{iI}`$ of left (right) gauge transformation on the trivializations $`B_iH`$ satisfying
$$\chi _i\alpha _{l,r}=\alpha _{l,r_i}\chi _i.$$
(8)
###### Proposition 2
A left (right) gauge transformation fulfills
$$\mathrm{\Delta }_𝒫\alpha _{l,r}=(\alpha _{l,r}id)\mathrm{\Delta }_𝒫.$$
(9)
Proof: Using the definitions one calculates
$`(\chi _i\alpha _{l,r}id)\mathrm{\Delta }_𝒫`$ $`=`$ $`(\alpha _{l,r_i}\chi _iid)\mathrm{\Delta }_𝒫`$
$`=`$ $`((\alpha _{l,r_i}id)(id\mathrm{\Delta })\chi _i`$
$`=`$ $`(id\mathrm{\Delta })\alpha _{l,r_i}\chi _i`$
$`=`$ $`(id\mathrm{\Delta })\chi _i\alpha _{l,r}`$
$`=`$ $`(\chi _iid)\mathrm{\Delta }_𝒫\alpha _{l,r}.`$
(9) follows from $`_iker\chi _i=0`$. $`\mathrm{}`$
Remark: Left (right) gauge transformations $`\alpha _{l,r}`$ can be equivalently defined as left (right) $`\iota (B)`$ module isomorphisms $`\alpha _{l,r}:𝒫𝒫`$ fulfilling
$`\mathrm{\Delta }_𝒫\alpha _{l,r}`$ $`=`$ $`(\alpha _{l,r}id)\mathrm{\Delta }_𝒫`$
$`\alpha _{l,r}\iota `$ $`=`$ $`\iota `$
$`\alpha (ker\chi _i)`$ $`=`$ $`ker\chi _i;iI.`$
###### Proposition 3
The set $`G_{l,r}`$ of all left (right) transformations is a group with the composition of maps as group multiplication. $`\mathrm{}`$
###### Proposition 4
Left (right) gauge transformations $`\alpha _{l,r}`$ on a locally trivial QPFB $`𝒫`$ are in one to one correspondence to linear maps $`g_{\alpha _{l,r}}:H𝒫`$ satisfying
$`g_{\alpha _{l,r}}(1)`$ $`=`$ $`1`$ (10)
$`\mathrm{\Delta }_𝒫(g_{\alpha _l}(h))`$ $`=`$ $`{\displaystyle g_{\alpha _l}(h_2)S(h_1)h_3}`$ (11)
$`\mathrm{\Delta }_𝒫(g_{\alpha _r}(h)`$ $`=`$ $`{\displaystyle g_{\alpha _r}(h_2)h_3S^1(h_1)}`$ (12)
$`{\displaystyle g_{\alpha _l}(h_1)g_{\alpha _l^1}(h_2)}`$ $`=`$ $`{\displaystyle g_{\alpha _l^1}(h_1)g_{\alpha _l}(h_2)}=\epsilon (h)1`$ (13)
$`{\displaystyle g_{\alpha _r}(h_2)g_{\alpha _r^1}(h_1)}`$ $`=`$ $`{\displaystyle g_{\alpha _r^1}(h_2)g_{\alpha _r}(h_1)}=\epsilon (h)1.`$ (14)
The correspondence is given by
$`\alpha _l(f)`$ $`=`$ $`{\displaystyle f_0g_{\alpha _l}(f_1)}`$ (15)
$`\alpha _r(f)`$ $`=`$ $`{\displaystyle g_{\alpha _r}(f_1)f_0}.`$ (16)
Proof: We will give the proof only for left gauge transformations because it works for right gauge transformations with the same arguments.
Assume that there is given a linear map $`g:H𝒫`$ with the properties $`\mathrm{\Delta }_𝒫(g(h))=g(h_2)S(h_1)h_2`$ and $`g(1)=1`$ such that there exists a linear map $`g^1:H𝒫`$ fullfilling $`g^1(h_1)g(h_2)=g(h_1)g^1(h_2)=\epsilon (h)1`$. It is easy to verify that the linear map $`\alpha :𝒫𝒫`$ defined by
$$\alpha (f):=f_0g(f_1)$$
is a left gauge transformation.
The proof of the other direction is more complicated. By definition for a given left gauge transformation $`\alpha _l`$ there exists a family $`(\alpha _{l_i})_{iI}`$ of left gauge transformations on the trivializations $`B_iH`$. Since the linear maps $`\alpha _{l_i}`$ are left $`(B_i1)`$-module isomorphisms, there exist left $`(B_{ij}1)`$ module isomorphisms $`\alpha _{l_i}^j:B_{ij}HB_{ij}H`$ satisfying
$$\alpha _{l_i}^j(\pi _j^iid)=(\pi _j^iid)\alpha _{l_i}.$$
(17)
These $`\alpha _{l_i}^j`$ satisfy the identity
$$\alpha _{l_i}^j=\varphi _{ij}\alpha _{j_l}^i\varphi _{ji},$$
(18)
where $`\varphi _{ij}`$ are the isomorphisms $`\varphi _{ij}:B_{ij}HB_{ij}H`$ induced from the transition functions $`\tau _{ij}`$ of the bundle $`𝒫`$ (see ). (18) is proved as follows. Let $`f𝒫`$. We know,
$$(\pi _j^iid)\chi _i(f)=\varphi _{ij}(\pi _i^jid)\chi _j(f),$$
(19)
therefore
$$(\pi _j^iid)\chi _i(\alpha _l(f))=\varphi _{ij}(\pi _i^jid)\chi _j(\alpha _l(f)).$$
With (8) and (17) follow the equations
$`(\pi _j^iid)\alpha _{l_i}\chi _i(f)`$ $`=`$ $`\varphi _{ij}(\pi _i^jid)\alpha _{j_l}\chi _j(f)`$ (20)
$`\alpha _{l_i}^j(\pi _j^iid)\chi _i(f)`$ $`=`$ $`\varphi _{ij}\alpha _{j_l}^i(\pi _i^jid)\chi _j(f).`$ (21)
Inserting (19) in (21) one obtains
$$\alpha _{l_i}^j\varphi _{ij}=\varphi _{ij}\alpha _{j_l}^i,$$
which proves (18). Because of this formula, the linear maps $`\tau _{\alpha _{l_i}}:HB_i`$ corresponding to the $`\alpha _{l_i}`$ satisfy
$$\pi _j^i(\tau _{\alpha _{l_i}}(h))=\tau _{ij}(h_1)\pi _i^j(\tau _{\alpha _{j_l}}(h_2))\tau _{ji}(h_3).$$
(22)
Now we define a family of linear maps $`g_{\alpha _{l_i}}:HB_iH`$ by
$$g_{\alpha _{l_i}}(h):=\tau _{\alpha _{l_i}}(h_2)S(h_1)h_3.$$
It is easy to see that
$$\alpha _{l_i}(ah)=(ah_1)g_{\alpha _{l_i}}(h_2).$$
Because of formula (22) the family of linear maps $`(g_{\alpha _{l_i}})_{iI}`$ fulfills
$$(\pi _j^iid)(g_{\alpha _{l_i}}(h))=\varphi _{ij}(\pi _i^jid)(g_{\alpha _{j_l}}(h)),$$
(23)
i.e. there exists a unique linear map $`g_{\alpha _l}:H𝒫`$ satisfying $`\chi _i(g_{\alpha _l}(h))=g_{\alpha _{l_i}}(h)`$ and $`\alpha _l(f)=f_0g_{\alpha _l}(f_1)`$. The properties of $`g_{\alpha _l}`$ are now easily verified by using the properties of the family $`(g_{\alpha _{l_i}})_{iI}`$. $`\mathrm{}`$
###### Proposition 5
There is a bijection between left and right gauge transformations.
Proof: Left and right gauge transformations are related by $`g_{\alpha _r}:=g_{\alpha _r}S^1`$. $`\mathrm{}`$
###### Proposition 6
Let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure on $`𝒫`$ and let $`\alpha _{l,r}`$ be a left (right) gauge transformation on $`𝒫`$.
Then the formulas
$`\alpha _l(\gamma )`$ $`:=`$ $`{\displaystyle \gamma _0g_{\alpha _l}(\gamma _1)}`$ (24)
$`\alpha _r(\gamma )`$ $`:=`$ $`{\displaystyle g_{\alpha _r}(\gamma _1)\gamma _0}`$ (25)
define left (right) $`\mathrm{\Gamma }_m(B)`$-module isomorphisms $`\alpha _{l,r}:hor\mathrm{\Gamma }_c(𝒫)hor\mathrm{\Gamma }_c(𝒫)`$ .
($`\mathrm{\Gamma }_m(B)`$ is the maximal embeddable LC-differential algebra induced from $`\mathrm{\Gamma }_c(𝒫)`$, see .) We leave the proof to the reader.
The concept of gauge transformations can be carried over to associated vector bundles. In general one can define gauge transformations on a locally trivial QVB $`E`$ as follows:
###### Definition 3
Let $`((E,B,\kappa ),V,(\zeta _i,J_i)_{iI})`$ be a locally trivial QVB. A gauge transformation on $`E`$ is a automorphism $`\eta :EE`$ with the properties
$`\kappa (a)\eta `$ $`=`$ $`\eta \kappa (a),aB`$ (26)
$`\eta (ker\zeta _i)`$ $`=`$ $`ker\zeta _i.`$ (27)
The last condition of this definition has the consequence that there are gauge transformations $`\eta _i:B_iVB_iV`$ and $`\eta _{ij}^i:B_{ij}VB_{ij}V`$ such that
$`\eta _i\zeta _i`$ $`=`$ $`\zeta _i\eta `$
$`\eta _{ij}^i(\pi _j^iid)`$ $`=`$ $`(\pi _j^iid)\eta _i.`$
###### Proposition 7
Gauge transformations on a locally trivial QVB $`E`$ are in one-to-one correspondence with families of gauge transformations $`\eta _i:B_iVB_iV`$ satisfying
$$\eta _{ij}^i=\varphi _{ij_E}\eta _{ij}^j\varphi _{ji_E}.$$
(28)
We omit the proof because it is quite analogous to the proof of Proposition 4 of .
###### Proposition 8
Let $`E(𝒫,F)`$ be an associated vector bundle. Every left gauge transformation on $`𝒫`$ determines a gauge transformation on $`E(𝒫,F)`$.
Proof: Let $`\alpha _l`$ be a left gauge transformation on $`𝒫`$. The linear map $`\eta _{\alpha _l}:E(𝒫,F)E(𝒫,F)`$ defined by
$$\eta _{\alpha _l}:=\alpha _lid$$
is seen to be a gauge transformation on $`E(𝒫,F)`$. $`\mathrm{}`$
Let $`\mathrm{\Gamma }_m(B)`$ be the maximal embeddable LC-differential algebra over $`B`$ induced from the differential structure on $`𝒫`$, and let $`E_\mathrm{\Gamma }(𝒫,F)`$ be the locally trivial QVB constructed in terms of $`E(𝒫,F)`$ and $`\mathrm{\Gamma }_m(B)`$. Because of Proposition 6, one can extend e very gauge transformation on $`E(𝒫,F)`$ determined by a gauge transformation on $`𝒫`$ to a module automorphism $`\eta _\alpha `$ of $`E_\mathrm{\Gamma }(𝒫,F)`$ by
$$\eta _{\alpha _l}:=ϵ_\mathrm{\Gamma }(\alpha _lid)ϵ_\mathrm{\Gamma }^1.$$
We end up this section with some remarks about the general structure of gauge transformations in the case when the structure group of the locally trivial QPFB is a compact quantum group.
For the algebra $`P(G)`$ of polynomial functions over such a compact quantum group $`G`$ one can construct the following linear basis ( see , and ). Let $`M`$ be the set of all irreducible unitary matrix co-representations of $`G`$.
(A unitary matrix co-representation is defined by an $`P(G)`$-valued $`N\times N`$-matrix $`(u_{ij})_{i,j=1,2,\mathrm{},N}`$ with $`\mathrm{\Delta }(u_{ij})=_ku_{ik}u_{kl}`$ and $`u_{ij}^{}=S(u_{ji})`$.)
Two co-representations $`\rho `$ and $`\sigma `$ are equivalent if there exists an intertwining operator. This defines an equivalence relation $``$ in $`M`$. Now one can select in every class $`\alpha =M/`$ a matrix co-representation $`(u_{ij}^\alpha )_{i,j=1,2,\mathrm{}N^\alpha }`$. It is proved in , and that the set of elements $`(u_{ij}^\alpha )_{\alpha M/;i,j=1,2,\mathrm{},N^\alpha }`$ is a linear basis of $`P(G)`$. This leads us to the following conlusion.
###### Proposition 9
Let $`𝒫`$ be a locally trivial QPFB where the structure group is a compact quantum group. The set of left (right) gauge transformations is in one to one correspondence with sets of invertible $`B_i`$-valued matrices $`(b_{i_{kl}}^\alpha )_{k,l=1,2,\mathrm{}N^\alpha \alpha M/}`$ satisfying
$$\pi _j^i(b_{i_{kl}}^\alpha )=\underset{m,n}{\overset{N^\alpha }{}}\tau _{ij}(u_{km}^\alpha )\pi _i^j(b_{j_{mn}}^\alpha )\tau _{ji}(u_{ml}^\alpha ).$$
(29)
More precisely (see also formula (22)), one can construct a gauge transformation by mapping every matrix $`u_{ij}^\alpha `$ to a set of invertible $`B_i`$-valued $`N^\alpha \times N^\alpha `$-matrices satisfying (29). Doing this for each $`\alpha `$ one obtains linear maps $`\tau _i:P(G)B_i`$ defined by
$$\tau _i(u_{kl}^\alpha )=b_{i_{kl}}^\alpha .$$
Since the matrices $`(b_{i_{kl}}^\alpha )`$ are invertible the $`\tau _i`$ are convolution invertible. Since formula (29) is fulfilled for each $`\alpha `$ the $`\tau _i`$ satisfy also (22) and it follows that they determine a gauge transformation.
## 2 Gauge transformations and connections
Since we have extended gauge transformations only to the subalgebra of horizontal forms $`hor\mathrm{\Gamma }_c(𝒫)`$ it is not possible to transform a connection by a gauge transformation analogous to the classical case by transforming the connection form $`\omega _{l,r}`$ . What we can transform is the covariant derivation $`D_{l,r}:=hord`$ corresponding to the connection.
Let $`D_{l,r}`$ be the left (right) covariant derivative corresponding to a left (right) connection and let $`\alpha _{l,r}`$ be a left (right) gauge transformation. One defines the linear map $`D_{l,r}^{}:hor\mathrm{\Gamma }_c(𝒫)hor\mathrm{\Gamma }_c(𝒫)`$ by
$$D_{l,r}^{}:=\alpha _{l,r}D_{l,r}\alpha _{l,r}^1.$$
In the sequel we want to discuss this formula for connections on trivial QPFB. We know, that a left (right) connection corresponds to a linear map $`A_{l,r}:H\mathrm{\Gamma }^1(B)`$ satisfying (54), (55),(79) and (80) of respectively. We are interested in the transformed maps $`A_{l,r}^{}`$ belonging to the $`D_{l,r}^{}`$. One calculates
$`D_l^{}(ah)`$ $`=`$ $`\alpha _lD_l\alpha _l^1(ah)`$
$`=`$ $`\alpha _lhor_ld{\displaystyle (a\tau _{\alpha _l^1}(h_1)h_2)}`$
$`=`$ $`\alpha _lhor_l({\displaystyle ((da)\tau _{\alpha _l^1}(h_1))\widehat{}h_2}+{\displaystyle a\tau _{\alpha _l^1}(h_1)\widehat{}dh_2})`$
$`=`$ $`\alpha _l({\displaystyle (da\tau _{\alpha _l^1}(h_1))\widehat{}h_2}{\displaystyle a\tau _{\alpha _l^1}(h_1)A_l(h_2)\widehat{}h_3})`$
$`=`$ $`{\displaystyle ((da)\tau _{\alpha _l^1}(h_1))\tau _{\alpha _l}(h_2)\widehat{}h_3}{\displaystyle a\tau _{\alpha _l^1}(h_1)A_l(h_2)\tau _{\alpha _l}(h_3)\widehat{}h_4}`$
$`=`$ $`(da)\widehat{}h{\displaystyle a\tau _{\alpha _l^1}(h_1)d\tau _{\alpha _l}(h_2)\widehat{}h_3}{\displaystyle a\tau _{\alpha _l^1}(h_1)A_l(h_2)\tau _{\alpha _l}(h_3)\widehat{}h_4}.`$
This shows that the linear map $`A_{}^{}{}_{l}{}^{}:H\mathrm{\Gamma }^1(B)`$ defined by
$$A_{}^{}{}_{l}{}^{}(h):=(id\epsilon )D_l^{}$$
has the form
$$A_l^{}(h)=\tau _{\alpha _l^1}(h_1)A_l(h_2)\tau _{\alpha _l}(h_3)+\tau _{\alpha _l^1}(h_1)d\tau _{\alpha _l}(h_2).$$
(30)
The same calculation for right connections leads to
$$A_r^{}(h)=\tau _{\alpha _r}(h_3)A_r(h_2)\tau _{\alpha _r^1}(h_1)\tau _{\alpha _r}(h_2)d\tau _{\alpha _r^1}(h_1).$$
(31)
In general the linear maps $`A_{l,r}^{}`$ do not satisfy the conditions (79) and (80) of respectively, i.e. $`D_{l,r}^{}`$ do in general not define a connection. Only in the special case when $`\mathrm{\Gamma }(H)`$ is the universal differential algebra, which means every linear map $`A:H\mathrm{\Omega }^1(B)`$ satisfying $`A(1)=0`$ defines a left and a right connection, every gauge transformation transforms connections in connections.
Because of this problem it seems to be necessary to introduce the following definition.
###### Definition 4
Let $`𝒫`$ be locally trivial QPFB and let $`\alpha _{l,r}`$ be a left(right) gauge transformation respectively. A left connection defined by $`hor_l`$ is called $`\alpha _l`$-covariant if $`D_l^{}`$ defined by
$$D_l^{}:=\alpha _lhor_ld\alpha _l^1$$
defines the left covariant derivation of a left connection.
A right connection defined by $`hor_r`$ is called $`\alpha _r`$-covariant if $`D_r^{}`$ defined by
$$D_r^{}:=\alpha _rhor_rd\alpha _r^1$$
defines the right covariant derivation of a right connection.
A left (right) connection is called covariant, if it is $`\alpha _{l,r}`$-covariant for all gauge transformations $`\alpha _{l,r}`$.
###### Definition 5
Let $`𝒢_{l,r}G_{l,r}`$ be a subgroup. A left (right) connection is $`𝒢_{l,r}`$-covariant if it is $`\alpha _{l,r}`$-covariant for all $`\alpha _{l,r}𝒢_{l,r}`$.
###### Proposition 10
Let $`𝒫`$ be a locally trivial QPFB and let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure of $`𝒫`$ where the differential algebra $`\mathrm{\Gamma }(H)`$ is the universal one. All left (right) connections are covariant.
For the proof see the remarks above.
###### Proposition 11
Let $`𝒫`$ be a locally trivial QPFB and let $`\mathrm{\Gamma }(𝒫)`$ be a differential structure on $`𝒫`$ where $`\mathrm{\Gamma }(H)`$ is bicovariant, i.e. the corresponding right ideal $`R`$ is Ad-invariant. Let $`𝒬G_{l,r}`$ be the subgroup of all gauge transformations which are differentiable algebra isomorphisms. All left and right connections are $`𝒬`$-covariant.
Proof: It is sufficient to prove this assertion on a trivial bundel $`BH`$. As noted at the beginning of Section 1, an algebra automorphism $`\alpha :BHBH`$ which is a gauge transformation corresponds to a homomorphism $`\tau _\alpha :HB`$ with the property that $`\tau _\alpha (H)`$ lies in the center of $`B`$. We have assumed that $`\alpha `$ is differentiable with respect to $`\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H)`$. Let $`J(BH)\mathrm{\Omega }(BH)`$ be the differential ideal corresponding to $`\mathrm{\Gamma }(B)\widehat{}\mathrm{\Gamma }(H)`$. The assumption that $`\alpha `$ is differentiable means $`\alpha _\mathrm{\Omega }(J(BH))=J(BH)`$. As shown in , Proposition 4, $`J(BH)`$ is generated by the sets
$`(id1)_\mathrm{\Omega }(J(B));\{{\displaystyle (1S^1(r_2))d(1r_1)}|rR\};`$
$`\{(a1)d(1h)(d(1h))(a1)|aB,hH\}.`$
Applying $`\alpha _\mathrm{\Omega }`$ to these sets and using the Ad-invariance of $`R`$ gives the following identities in $`\mathrm{\Gamma }(B)`$.
$`{\displaystyle \tau _{\alpha ^1}(r_1)d\tau _\alpha (r_2)}`$ $`=`$ $`0,rR`$
$`(da)\tau _\alpha (h)`$ $`=`$ $`\tau _\alpha (h)da,aB,hH.`$
Let $`A_{l,r}`$ be the linear map corresponding to a connection on $`BH`$. Using formula (30) one obtains for the transformed linear maps $`A_{l,r}^{}`$
$$A_{l,r}^{}(h)=A_{l,r}(h_2)\tau _\alpha (S(h_1)h_3)+\tau _{\alpha ^1}(h_1)d\tau _\alpha (h_2).$$
Inserting an element $`rR`$ in this equation and using the Ad-invariance of R, i.e. $`r_2S(r_1)r_3RH`$, and the properties of $`\tau _\alpha `$ one obtains
$$A_{l,r}^{}(r)=0,rR.$$
Thus, $`A_{l,r}^{}`$ defines a left (right) connection again. $`\mathrm{}`$
###### Proposition 12
Let $`𝒫`$ be a locally trivial QPFB. Assume there is given a differential structure on $`𝒫`$. Let $`D_{l,r}`$ be a left (right) covariant derivative and $`\alpha _{l,r}`$ a left (right) gauge transformation. The Map $`D_{l,r}^{}`$ defined by
$$D_{l,r}^{}:=\alpha _{l,r}D_{l,r}\alpha _{l,r}^1$$
is a left (right) covariant derivative.
We leave the proof for the reader.
###### Proposition 13
Let $`𝒫`$ be a locally trivial QPFB, let $`D_{l,r}`$ be a covariant derivative and let $`\alpha _{l,r}`$ be a gauge transformation.
Then there are the following transformation formulas for the left (right) curvature form $`\mathrm{\Omega }_{l,r}`$:
$`\mathrm{\Omega }_{}^{}{}_{l}{}^{}(h)`$ $`=`$ $`{\displaystyle g_{\alpha _l^1}(h_1)\mathrm{\Omega }(h_2)g_{\alpha _l}(h_3)}`$ (32)
$`\mathrm{\Omega }_{}^{}{}_{r}{}^{}(h)`$ $`=`$ $`{\displaystyle g_{\alpha _r}(h_3)\mathrm{\Omega }(h_2)g_{\alpha _r^1}(h_1)}`$ (33)
Proof: Using formulas (105) and (106) of and (24) and (25) one obtains
$$(D_{l,r}^{})^2=\alpha _{l,r}(D_{l,r})^2\alpha _{l,r}^1$$
and
$`\alpha _l(D_l)^2\alpha _l^1(f)`$ $`=`$ $`{\displaystyle f_0g_{\alpha _l^1}(f_1)\mathrm{\Omega }_l(f_2)g_{\alpha _l}(f_3)}`$
$`\alpha _r(D_r)^2\alpha _r^1(f)`$ $`=`$ $`{\displaystyle g_{\alpha _r}(f_3)\mathrm{\Omega }_r(f_2)g_{\alpha _r^1}(f_1)f_0}.`$
for $`f𝒫`$. $`\mathrm{}`$
At the end of this section let us remark that on locally trivial QVB every gauge transformation $`\eta `$ transforms connections $``$ in connections $`^{}`$ by
$$^{}=\eta \eta ^1.$$
This is analogous to Proposition 12.
## 3 Example
This example is constructed to show that there exist nonclassical gauge transformations.
Here we glue together a noncommutative “tube” with the quantum group $`SU_\nu (2)`$ along the classical subspace $`S^1`$ and construct a $`SU_\nu (2)`$ bundle over this “base”. First, we define the algebra over the noncommutative “tube” as the algebra $`B_1`$ generated by the elements
$$x,x^{},y=y^{}$$
satisfying the relations
$`xx^{}`$ $`=`$ $`x^{}x=1`$
$`xy`$ $`=`$ $`qyx`$
$`x^{}y`$ $`=`$ $`q^1yx^{},`$
where $`q(0,1]`$. It is easy to see that there exists a surjective homomorphism $`\pi _2^1:BP(S^1)`$ defined by
$`\pi _2^1(x)`$ $`=`$ $`a`$
$`\pi _2^1(x^{})`$ $`=`$ $`a^{}`$
$`\pi _2^1(y)`$ $`=`$ $`0,`$
where $`a`$ is the generator of $`P(U(1))`$.
There also exists a surjective homomorphism $`\pi _1^2:P(SU_\nu (2))P(S^1)`$ defined by
$`\pi _1^2(\alpha )`$ $`=`$ $`a`$
$`\pi _1^2(\alpha ^{})`$ $`=`$ $`a^{}`$
$`\pi _1^2(\gamma )`$ $`=`$ $`\pi _1^2(\gamma ^{})=0,`$
where $`\alpha ,\gamma `$ are the usual generators of $`P(SU_\nu (2))`$. Our basis algebra $`B`$ is defined as the gluing of $`B_1`$ and $`P(SU_\nu (2))`$ by means of $`\pi _2^1`$ and $`\pi _1^2`$,
$$B:=\{(f_1,f_2)B_1P(SU_\nu (2))|\pi _2^1(f_1)=\pi _1^2(f_2)\}.$$
Now one chooses the transition functions $`\tau _{ij}:P(SU_\nu (2)P(S^1)`$ as follows:
$`\tau _{12}(\alpha )`$ $`=`$ $`a`$
$`\tau _{12}(\alpha ^{})`$ $`=`$ $`a^{}`$
$`\tau _{12}(\gamma )`$ $`=`$ $`\tau _{12}(\gamma ^{})=0`$
and obtains a locally trivial QPFB $`𝒫`$ with “structure group” $`SU_\nu (2)`$.
According to Proposition 4, to construct a left gauge transformation $`\alpha _l`$ we have to find a linear map $`g_{\alpha _l}:P(SU_\nu (2))𝒫`$ fulfilling (10) - (14.
First we define the linear maps $`\tau _1^{(n)}:P(SU_\nu (2))B_1`$ and $`\tau _2^{(n)}:P(SU_\nu (2))P(SU_\nu (2))`$, $`n1`$. $`\tau _1^{(n)}`$ is assumed to be a homomorphism defined by
$`\tau _1^{(n)}(\alpha )`$ $`=`$ $`x^n`$
$`\tau _1^{(n)}(\alpha ^{})`$ $`=`$ $`x_{}^{}{}_{}{}^{n}`$
$`\tau _1^{(n)}(\gamma )`$ $`=`$ $`\tau _1^n(\gamma ^{})=0.`$
$`\tau _2^{(n)}`$ is defined as follows:
$$\tau _2^{(n)}(h)=h_1h_2\mathrm{}h_n.$$
The linear maps $`\tau _i^{(n)}`$ are convolution invertible with convolution inverse
$$\tau _{1}^{(n)}{}_{}{}^{1}(h)=\tau _1^{(n)}(S(h))$$
and
$$\tau _{2}^{(n)}{}_{}{}^{1}(h)=S(h_1)S(h_2)\mathrm{}S(h_n).$$
By an easy calculation one obtains the identity
$$\pi _2^1(\tau _1^{(n)}(h))=\tau _{12}(h_1)\pi _1^2(\tau _2^{(n)}(h_2))\tau _{21}(h_3)=\pi _1^2(\tau _2^{(n)}(h)),$$
hence (see formula (22) ) there exists a convolution invertible map $`g_{\alpha _l}^{(n)}:P(SU_\nu (2))𝒫`$ which determines a left gauge transformation. |
warning/0002/cond-mat0002167.html | ar5iv | text | # Analysis of soft optical modes in hexagonal BaTiO3: transference of perovskite local distortions
\[
## Abstract
We have performed detailed first-principles calculations to determine the eigenvectors of the zone-center modes of hexagonal BaTiO<sub>3</sub> and shown that the experimentally relevant low-energy modes (including the non-polar instability) can be represented as suitable combinations of basic local polar distortions associated with the instability of the cubic perovskite phase. The hexagonal structure provides a testing ground for the analysis of the influence of the stacking of TiO<sub>6</sub> octahedra: the occurrence of relatively high-energy chains of dipoles highlights the importance of local effects related to the coherent hybridization enhancement between Ti and O ions. Our results provide simple heuristic rules which could be useful for the analysis of related compounds.
\]
Barium titanate has two structural polymorphs: the cubic perovskite type (c-BT) and its hexagonal modification (h-BT) with six formula units per unit cell. While c-BT has been one of the best studied ferroelectric materials for decades , most of the work on the structural and dielectric properties of h-BT is quite recent . As shown in Fig. 1, h-BT is also composed of TiO<sub>6</sub> groups, albeit with a different stacking than the perovskite form. It seems well established that the hexagonal polymorph undergoes two zone-center structural phase transitions: at $`222`$ K from the high temperature $`P6_3/mmc`$ hexagonal phase to a second non-polar $`C222_1`$ phase, and at $`74`$ K to a ferroelectric $`P2_1`$ phase. The first transition is associated with the softening of an optical mode and the second attributed to a shear strain instability; but a detailed analysis is lacking due to the absence of structural information on the two low-symmetry phases, and little is known about the microscopic origin of the instabilities.
On the other hand, the discovery of a giant LO-TO splitting in h-BT by Inoue et al. suggested that its ferroelectric modes have the same origin as those of c-BT. In the cubic phase the ferroelectric instabilities can be essentially described as chains of dipoles that originate in the movement of Ti ions relative to their surrounding O<sub>6</sub> octahedra with a minor distortion of the later. Additional evidence in support of this view is provided by the successful use of polar local modes by Zhong et al. in the construction of an effective Hamiltonian for c-BT. In view of the basic structural similarities between the cubic and hexagonal forms of BaTiO<sub>3</sub>, it is meaningful to ask whether the local modes that describe the unstable branches in c-BT can somehow be transferred to h-BT and serve as a basis to discuss the low-energy distortions of the structure.
Here we show the results of first-principles calculations that provide for the first time structural information on the low-symmetry phases of the hexagonal polymorph of BaTiO<sub>3</sub> and reveal the microscopic nature of the modes. Our analysis proves that the structure of the experimentally found zone-center optical soft modes in h-BT is indeed characterized by the same distortions of the TiO<sub>6</sub> octahedra that are relevant in c-BT, leading to similar chains of dipoles in the hexagonal structure.
In h-BT, the experimentally found optical soft modes are: the zone-center instability that drives the phase transition at $`222`$ K and transforms according to the $`E_{2u}`$ irreducible representation (irrep) of $`6/mmm`$, and the $`A_{2u}`$ ferroelectric mode that softens (though remaining stable) in the temperature range of the $`C222_1`$ phase and is responsible for the giant LO-TO splitting. Our calculations agree with this experimental evidence. We performed a full ab-initio relaxation of the thirty-atom h-BT structure, resulting in lattice parameters $`a`$=$`10.68`$ and $`c`$=$`26.053`$ a.u. (to be compared with experimental values of 10.77 and 26.451 , respectively). The five free internal coordinates are also in excellent agreement with the experimentally determined ones (within $`1\%`$). After computing the force-constant matrix at $`\mathrm{\Gamma }`$ and diagonalizing it within the subspaces of the appropriate symmetries, we found an unstable E<sub>2u</sub> mode and a soft but not unstable A<sub>2u</sub> mode. \[It should be noted that if the calculations are performed using the experimental lattice parameters (i.e., at a larger cell volume), the A<sub>2u</sub> mode is found to be unstable (and the E<sub>2u</sub> instability is more pronounced). The fact that the A<sub>2u</sub> mode is very close to being unstable could be particularly relevant for the phase transition at $`74`$ K.\]
A first analysis of the eigenvectors for both soft modes reveals that the Ba contribution is small (around $`10\%`$ of the total mode norm, compared to $`4\%`$ in the perovskite soft mode). Moreover, we checked that if the Ba ions are frozen at their high-symmetry positions the modes are still soft, so we do not consider them the following discussion, and focus on the distortions of the TiO<sub>6</sub> groups.
To make the comparison to the perovskite quantitative, let us consider the polar deformations of the TiO<sub>6</sub> groups of c-BT. These are shown in Fig. 2, where we assume that the Ti ion is located at the origin of coordinates, so only the displacement patterns of the O ions need to be considered. For each spatial direction $`\alpha `$=$`x,y,z`$ we have two symmetry-adapted distortions denoted by $`\widehat{s}_{1,\alpha }`$ and $`\widehat{s}_{2,\alpha }`$ and transforming according to the $`T_{1u}`$ (vector like) irrep of $`m3m`$, the point group of the regular octahedra of c-BT. In terms of this basis, the tetragonal ferroelectric distortion in c-BT (along $`x`$ for concreteness) can be written as $`0.69\widehat{s}_{1x}+0.73\widehat{s}_{2x}`$, with the distorted octahedra exhibiting point group symmetry $`4mm`$ . In h-BT there are two kinds of octahedra: those centered around Ti ions at $`2a`$ Wyckoff positions with $`\overline{3}m`$ point symmetry (denoted by TiO<sub>6</sub>(1) in Fig. 1) and those arranged around Ti ions at $`4f`$ Wyckoff positions with $`3m`$ point symmetry (TiO<sub>6</sub>(2) in the figure) . The octahedra in the first set are coordinated in the same way as in c-BT, i.e., by sharing O ions with six other octahedra. Those in the second set are linked to other $`1+3`$ octahedra by sharing one O<sub>3</sub> face and single O ions respectively.
Due to the low (as compared to the case of c-BT) symmetry of the octahedra in h-BT, their distortions associated to general $`E_{2u}`$ and $`A_{2u}`$ modes can be decomposed in a relatively large number of symmetry-adapted displacement patterns. Among all the possible ones, we restrict ourselves to those of c-BT type and check if they can actually account for the structure of the soft modes . For instance, a general $`A_{2u}`$ distortion leads the crystal to a phase with space group $`P6_3mc`$, in which the TiO<sub>6</sub>(1) groups reduce their point symmetry to $`3m`$ (see Table I). As shown in Fig. 3.b, c-BT distortions in the $`s_{ix}=s_{iy}=s_{iz}`$ component combination (Ti ions move towards O<sub>3</sub> faces as in the rhombohedric phase of c-BT) produce this symmetry breaking. In the case of an $`E_{2u}`$ distortion, TiO<sub>6</sub>(1) reduces its point symmetry to $`2`$, and the appropriate c-BT mode has the form $`s_{ix}=s_{iy}`$ (orthorhombic), as shown in Fig. 3.c.
For the two soft modes of h-BT, we considered separately the various classes of octahedra, computed from our ab-initio eigenvectors the displacement of the O ions relative to the Ti ion, and performed a projection of the resulting distortion field into the c-BT type symmetry-adapted modes. The results (last column of Table I) present two main features: First, almost $`100\%`$ of the total structural change associated with both soft modes can be described in terms of the c-BT type polar distortions (normalization is chosen in such a way that, for instance, for the first row in Table I we have $`(0.62^2+0.78^2)\times 100=99.3\%`$). Second, the components $`s_1;s_2`$ are always similar in magnitude to those of c-BT ($`0.69;0.73`$) and present a positive $`s_1/s_2`$ ratio, which implies that the O<sub>6</sub> octahedral cage moves almost rigidly relative to the Ti ion also in h-BT. We have also computed the Born effective charge associated with the ferroelectric $`A_{2u}`$ soft mode and found it unusually large ($`Z^{}=11.29`$), further confirming the relation with the (rhombohedric) ferroelectric instability of c-BT (for which $`Z^{}=9.956`$.
We have proved then that at a local level the soft modes in h-BT can be described by the same distortion vectors that determine the c-BT polar instability. In the crystal as a whole, these local polar distortions lead to chains of dipoles, which points at the $`E_{2u}`$ instability and the softness of the $`A_{2u}`$ mode of h-BT as being caused by Coulomb destabilizing forces, as it happens in the cubic perovskite. The ferroelectric $`A_{2u}`$ soft mode, polarized along $`z^{}`$, is roughly depicted in Fig. 1. In the $`E_{2u}`$ distortion the chains of dipoles lay on the $`x^{}y^{}`$ plane and alternate in orientation with a zero net polarization (the resulting $`C222_1`$ phase could be informally considered anti-ferroelectric rather than paraelectric).
From first-principles studies of the c-BT phase it is known that parallel dipole chains are very weakly coupled, so that a transverse modulation of a chain-like instability is not energetically relevant, and, therefore, unstable TO normal modes exist almost in the whole Brillouin Zone (BZ) . \[The only exception are $`k`$ points near $`𝐤_R=\frac{2\pi }{a}(1,1,1)`$, for which we have an anti-phase modulation of the Ti displacements (Ti$``$O$``$Ti–O) in the three spatial directions, so the long-range destabilizing forces are always canceled.\] If this view is taken to its logical conclusion, we could expect to find more zone-center soft modes in h-BT, corresponding to the other possible distributions of chains of dipoles. Table II enumerates all the possibilities. Apart from the already discussed $`A_{2u}`$ and $`E_{2u}`$ modes, our ab-initio calculations show that there is one $`E_{1g}`$ mode that is indeed rather low in energy, while the ferroelectric $`E_{1u}`$ and the $`E_{2g}`$ modes that are dominated by the movement of Ti ions are quite hard. In order to explain this result, let us remark that for the $`E_{2u}`$ and $`A_{2u}`$ soft modes the distortion is such that if an O ion is approached by one of its two Ti neighbors the second Ti ion moves away from it. This reflects the hybridization of the Ti $`3d`$ and O $`2p`$ electronic states, which has been shown to be essential for the occurrence of the c-BT ferroelectric instability . It can be checked that any other zone-center arrangement of the chains of dipoles results in either two Ti ions approaching one O ion (for example, if the two Ti ions in one of the O$`{}_{3}{}^{}`$Ti$``$O$`{}_{3}{}^{}`$Ti$``$O<sub>3</sub> groups depicted in Fig. 1 move in the same way in the $`x^{}y^{}`$ plane, there is at least one oxygen of the shared face that is approached by both) or in the second titanium not moving away from an oxygen. In the former case ($`E_{1u}`$ and $`E_{2g}`$) the effect of the hybridization is lost and the corresponding modes are hard. In the latter ($`E_{1g}`$ and $`B_{1g}`$), the hardening is not as strong. Thus, we conclude that the particular stacking of the TiO<sub>6</sub> groups in h-BT causes (through this local effect) the relatively high energy of some chain-like distortions. We can then formulate two “rules of thumb” for the characterization of a given locally polar distortion as low-energy: First ($``$1): “There need to be chains of dipoles (without regard for their transverse modulation)”. Second ($``$2): “The distribution of such chains must lead to a coherent hybridization enhancement between Ti and O ions, where the word coherent means that the destabilizing effect is lost when two Ti ions approach the same O ion”. These heuristic rules could be used to predict the occurrence of locally polar soft modes at other $`k`$ points of the BZ of h-BT, as well as in other structures with TiO<sub>6</sub> octahedra as basic building blocks.
In summary, first-principles calculations of the character of the zone-center modes of hexagonal BaTiO<sub>3</sub> support the physically appealing idea that the experimentally relevant soft modes (including the non-polar instability) can be represented as combinations of local polar distortions transferred directly from the cubic perovskite form of the compound. Our results lead also to heuristic rules that provide insight into the influence of the arrangement of TiO<sub>6</sub> octahedra on the low-energy dynamics of a structure .
This work was supported in part by the UPV research grant 060.310-EA149/95 and by the Spanish Ministry of Education grant PB97-0598. J.I. acknowledges financial support from the Basque regional government. |
warning/0002/astro-ph0002079.html | ar5iv | text | # Pinpointing the Position of the Post-AGB Star at the Core of RAFGL 2688 using Polarimetric Imaging with NICMOS
## 1 Introduction
The bipolar structures exhibited by a substantial fraction of the known planetary nebulae likely arise during the last, rapid, pre-planetary nebula (PPN) stage of evolution of intermediate-mass (1–8 M) stars off the asymptotic giant branch (AGB). A popular, albeit largely untested, model for such bipolarity is that the central AGB star possesses a companion that aids in the buildup of a dense, dusty equatorial torus surrounding the central star (e.g., Soker (1998)). Alternatively, the fossil remnant of a $`\beta `$ Pic-like main-sequence disk may bear responsibility for triggering bipolarity during post-main sequence evolution (Kastner & Weintraub (1995)). Whatever the mechanism that abets their formation, bipolar PPN typically show two bright reflection lobes separated by a dark dust lane. The star that illuminates the polar lobes presumably is located at or near the center of the equatorial, dust torus. While this geometry obscures the central star along our direct line of sight, photons readily escape the nebular core in the polar directions and subsequently are scattered by dust grains located primarily in the walls of the rarefied, expanding lobes. As even the lobe walls tend to be optically thin in the near-infrared, such photons can be singly scattered out of the nebula into our line of sight. Single scattering produces polarized light that contains a record of the original direction of the unpolarized light source; therefore, polarimetric maps of such polarized nebulae contain clues as to the locations of their illuminating sources, even if those stars lie hidden inside dust lanes.
Recent direct imaging of RAFGL 2688 (the Egg Nebula) with the Near Infrared Camera and Multi-Object Spectrometer (NICMOS) aboard the Hubble Space Telescope (HST) (Sahai et al. (1998)) revealed a compact red source just south of the bottom of the northern reflection lobe. However, initial analysis of the polarimetric maps from NICMOS indicated that this red source was not the primary illuminator of the reflection nebulosity; this object is most likely a companion to the post-AGB star that lurks in the core of the Egg Nebula. From a preliminary examination of the 2.0 $`\mu `$m polarimetric map, Sahai et al. suggested that the obscured, post-AGB star was located $``$ 750 AU (0$`\stackrel{}{\mathrm{.}}`$75) south of the red companion.
In this paper, we present a rigorous analysis of the 2.0 $`\mu `$m polarization map of RAFGL 2688 obtained by NICMOS. We determine the precise position of the post-AGB star in the core, assess the relationship of the red source to the illuminator star, and discuss the implications of this work for understanding the formation of the Egg Nebula and of other bipolar PPN.
## 2 Polarization Data Analysis
The data and data reduction methods used in this study were first presented by Sahai et al. (1998). In brief summary, RAFGL 2688 was imaged through the POL0L, POL120L and POL240L filters with camera 2 (NIC2) of NICMOS, using integration times of 1215 s for each filter. These filters are centered at 1.994 $`\mu `$m and have a full-width-half-maximum of 0.2025 $`\mu `$m. The field of view for these images is 19$`\stackrel{}{\mathrm{.}}`$5$`\times `$19$`\stackrel{}{\mathrm{.}}`$3 and the plate scale is 0$`\stackrel{}{\mathrm{.}}`$076/pixel (Thompson et al. (1998)). The calculations of fractional polarization (p) and polarization position angle ($`\theta `$) are carried out as described by Hines (1998)<sup>1</sup><sup>1</sup>1Note that the coefficients for polarimetric imaging calculations have been updated; see http://www.stsci.edu/instruments/nicmos/nicmos\_polar.html and Hines, Schmidt & Schneider 1999; however, we find that the best position angle calculations include the addition of a small, constant angle $`\varphi `$ to $`\theta `$, i.e.
$$\theta =\frac{1}{2}\mathrm{tan}^1\left(\frac{U}{Q}\right)+\varphi ,$$
where $`U`$ and $`Q`$ are the Stokes vectors obtained from the polarimetric images. The offset angle $`\varphi `$ could represent a systematic rotation of the filters in the polarization filter set from their nominal position angles. For example, if the three polarizing filters were designed to lie at position angles 0, 120, and 240, they actually are found at position angles 0 \+ $`\varphi `$, 120 \+ $`\varphi `$, and 240 \+ $`\varphi `$. Alternatively, $`\varphi `$ could represent uncertainties in our knowledge of the absolute position angles assumed for the polarization calibrators. We suggest that RAFGL 2688 represents the best absolute position angle calibrator for NICMOS polarimetric data. As explained in §4.1, we have determined empirically that $`\varphi `$ = 4.0 $`\pm `$ 0.2. A polarization map of RAFGL 2688, made with $`\varphi `$ = 4.0, is presented in Figure 1.
## 3 The Polarization Structure of the Nebula
A centrosymmetric pattern is the dominant single feature of the polarization map (Fig. 1); however, it is apparent by careful inspection of Fig. 1 that the polarimetric centroid is not spatially coincident with the source (labeled A) at the southern tip of the northern lobe (see §4). Overall, the nebula is very highly polarized, with virtually the entire southern lobe polarized with $`p`$ $`>`$ 0.50 (see Fig. 6 in Sahai et al. 1998 for a grayscale map of the polarized intensity). A second strong feature of the polarization structure is the apparent point symmetry of the polarization pattern around position B, which we describe below. The implication of such a symmetry for the origin of the bipolar lobes is discussed later (see §4.2).
The southern lobe is more highly polarized overall than the northern lobe (Figure 2). In the north, only 11 pixels show vectors with polarization amplitudes above 0.7; all of these vectors are on or west of the polar axis, with all but one at least 5<sup>′′</sup> from the center of the nebula (Figure 2a). In contrast, $``$300 pixels in the southern lobe have $`p`$ $`>`$ 0.7; these pixels are dominantly on the eastern side of the polar axis and all of them lie more than 4<sup>′′</sup> from the center of the nebula, demonstrating a strong point symmetry to the polarization pattern around the nebular core. An additional $``$1000 pixels are polarized with 0.6 $`<`$ $`p`$ $`<`$ 0.7 (Figure 2b). In the north, virtually all these vectors lie west of the polar axis, stretching inwards along the west limb of the reflection lobe from a distance of $``$7<sup>′′</sup> to just more than 3<sup>′′</sup> from the center. In the south, these vectors are uniformly spread across the lobe in the outer regions and more concentrated to the east of the polar axis closer to the core. Most of the rest of the southern lobe is polarized at a level $`p`$ $`>`$ 0.50 (Figure 2c). In the north, the polarization vectors in the range 0.4 $`<`$ $`p`$ $`<`$ 0.6 cover most of the center of the lobe (Figure 2c); the region covered by these vectors stretches radially away from the core along the eastern side; the polarization vectors in the range 0.4 $`<`$ $`p`$ $`<`$ 0.6 also cover the center of the southern lobe at small radial distances and then this region stretches outwards from the core along the western side. Finally, the outer edges of the northern lobe nearest to the nebular core are dominated by polarization amplitudes in the 0.15–0.40 range (Figure 2d).
## 4 The Polarimetric Centroid
### 4.1 Method of Determination
To determine the position of the source that illuminates the nebula, we have used the method presented by Weintraub & Kastner (1993), coded into a program in the software package IDL. This method takes advantage of the fact that a dust grain that singly scatters photons out of the nebula imparts a polarization position angle to the scattered light that is perpendicular to the scattering plane, i.e., perpendicular to the projected direction from that dust grain to the source of illumination. Thus, for every pair of polarization vectors in a map, we can draw perpendiculars to each vector and determine a point of intersection. Ideally, for noiseless data and purely singly scattered photons, all the pairs of vectors would have a unique intersection, the polarimetric centroid, which should mark the intersection between the polar axis and the disk midplane (assuming the illuminating source is identical with the central star of the nebula and that the central star lies at the geometric center of the nebula).
Even for noisy data and a mixture of singly and multiply scattered photons, one can use the method of intersections of polarization perpendiculars to determine the polarization centroid, albeit with finite positional error bars (Weintraub & Kastner (1993)). For a given data set, the accuracy with which we can determine the centroid depends on the absolute calibration of the position angles and thus depends on our knowledge of $`\varphi `$. If $`\varphi `$ is marginally inaccurate, the polarimetric centroid will be poorly determined while if $`\varphi `$ is quite inaccurate, there will be no polarimetric centroid in the map at all. Thus, we have determined $`\varphi `$ by examining a range of $`\varphi `$ values between $`10^{}`$ and $`+10^{}`$ and adopting the value that minimizes the uncertainty in determining the polarimetric centroid. In calculating the polarimetric centroid, we limit the calculation to the $`>`$8000 pixels containing flux levels with signal-to-noise ratios greater than six in all three of the POL0L, POL120L, and POL240L images.
Many of these pairs of vectors have nearly parallel position angles. For vector pairs with similar position angles, especially given even a small error in determining the true position angles, the intersection position is poorly determined. We therefore impose an additional constraint: we reject all vector pairs for which the angle between the vectors (modulo 180) is less than 20. This ensures that the small uncertainties in the position angle calculations do not produce large uncertainties in the actual position of the centroid. In practice, in addition to noise, many of the pixels, usually those with polarization vectors with lower polarization amplitudes, represent parts of the reflection nebula in which multiple scattering is probably dominant. Thus, for our final calculations, we placed a limit on the minimum allowable fractional polarization to be $`p_{min}`$ $``$ 0.15 in order to exclude lines of sight dominated by multiple scattering.
After calculating the intersection points for the complete set of allowable vectors and vector pairs, we calculate the statistical mean and the standard deviation of the mean ($`\sigma `$) for the polarization centroid. We then repeat this calculation, keeping only intersection points within a 3-$`\sigma `$ rejection threshold of the initially determined mean. We continue with this process, iteratively, until the solution converges on the polarimetric centroid (denoted B). We find that the initial calculation typically lies within 0.1 pixels ($`<`$ 0$`\stackrel{}{\mathrm{.}}`$01) of the final position and the calculation converges after only $``$five iterations and after rejecting only $``$2%–4% of the total possible intersections. Changing the rejection threshold appears to affect only the size of the uncertainty and the rate of convergence, not the position of the polarization centroid itself.
### 4.2 Results
In Figure 3, we present the same map as shown in Figure 1 but drawn with all the vectors perpendicular to the polarization position angles. These vectors clearly point to a single intersection point, the polarimetric centroid (labeled B in Fig. 1,3-6). In addition, this map illustrates, very clearly, the symmetry axis of the nebula, as seen in scattered light.
By examining solutions where $`p_{min}`$ ranges from 0.15 to 0.35, we find that the centroid lies 0$`\stackrel{}{\mathrm{.}}`$52 $`\pm `$ 0$`\stackrel{}{\mathrm{.}}`$02 west and 0$`\stackrel{}{\mathrm{.}}`$16 $`\pm `$ 0$`\stackrel{}{\mathrm{.}}`$03 south (Figure 4) of the isolated intensity peak at the southern tip of the north lobe (position A), well within the dark dust lane that cuts across the middle of the bipolar nebula. The positional uncertainty is dominated by the systematic differences between solutions found when selecting different values of $`p_{min}`$, rather than by the statistical errors in a single calculation (which are more than an order of magnitude smaller).
We have used the position of the polarimetric centroid combined with the vector pattern to determine the direction of the projection of the polar (major) axis of the Egg Nebula. One can see (Fig. 3) that the projected polar axis, drawn at a position angle of 12 (east of north), runs exactly parallel to the straight lines formed by the alignment of the perpendiculars (of the polarization vectors) along the central axis of both the north and south scattering lobes. A change in more than 1 in the position angle of the polar axis produces a clear error in the left-right symmetry of the lobes, as defined by the polarization vectors. Thus, we believe this determination of the projected position angle of the polar axis represents an improvement over the previously inferred angle of 15 (Ney et al. (1975)).
In projection, the centroid is located much closer to the northern than the southern lobe. The fact that B lies closer to the southern tip of the northern lobe than to the northern tip of the southern lobe is consistent with previous determinations that the polar axis of the system is inclined such that the northern lobe is tilted toward the observer. This geometry causes the optically thick equatorial torus to obscure the innermost part of the southern lobe but permits us to view most of the inner regions of the northern lobe.
It is interesting to note the point symmetry between the two scattering lobes. In the north, the majority of the total intensity of the nebula is east of the polar axis, including the brightest reflection peaks (see Fig. 1). In contrast, in the south, most of the reflection nebula is found to the west of the polar axis. In both lobes, the morphologically larger side of the nebula is the side showing lower overall polarization levels. We also see that the polar axis runs through the eastern side of the inward extension of the southern lobe and through the western side of the inward extension of the northern lobe. The simplest mechanism for producing point symmetric structure in the nebula is the operation of collimated bipolar outflows. Sahai & Trauger (1998) have argued, based on finding a high degree of point symmetry in the morphologies of their sample of young planetary nebulae, that such outflows are the primary agent for producing aspherical structure in planetary nebulae.
## 5 The Illuminator Star and its Surroundings
The polarimetric centroid presumably marks the position of the post-AGB star that illuminates most or all of both the northern and southern reflection lobes of the Egg Nebula. We now consider whether this illuminator and the intensity peak A constitute a widely spaced ($`>`$550 AU) binary system.
If a field star were at position A, such a star would reveal itself in an Airy pattern in the total intensity profile, as NICMOS generates such patterns even for very faint point sources. The absence of such a pattern indicates that the intensity peak A is an extended object. Such an object could be either a region of enhanced dust density that reflects light from B or a star embedded in the nebula that illuminates and heats the local pocket of dust around it.
In Figure 5, we present a polarization map of the same region as seen in Figure 4; however, in order to focus on the polarization behavior near A, we present in Fig. 5 only the polarization vectors with amplitudes $`p`$ $`<`$ 0.15. If a point source at position A suffers little local extinction, then it becomes a source of 2 $`\mu `$m photons which should generate some sort of centrosymmetric polarization pattern centered on A while the direct line of sight to A should show a low polarization level. Given the local presence of the illuminator star at B, we might expect this pattern to be distorted by the influence of a second photon source.
In examining Fig. 5, we find neither an indication of any kind of centrosymmetric pattern, even a strongly distorted one, centered on the position of the intensity peak at A, nor a simple, centrosymmetric pattern focused on the position of the illuminator at B, similar to that which characterizes the vectors in the rest of the nebula. Instead, close to A, we find a region marked by extremely low polarization levels and a disorganized polarization pattern, despite the fact that the signal-to-noise ratio is high. A somewhat more organized vector pattern is seen in the vectors that lie northeast, north and northwest of A, and which appear to define a centrosymmetric pattern centered on B.
If intensity peak A were simply a region of enhanced density of cold dust, we should see a pattern of highly polarized vectors at A suggesting direct illumination from position B, as is seen at other intermediate intensity peaks further out in the northern lobe. The absence of such a pattern suggests that peak A is self-luminous; however, the lack of any Airy profile as would be expected from a point source indicates that the source at A, at 2 $`\mu `$m, is seen as a small, extended nebula. At this position in the nebula, the local NICMOS point spread function generated by emission from the extended source at A, combined with the illumination of dust in this vicinity by B, generates a disorganized polarization pattern marked by relatively low polarization levels. This analysis therefore supports the suggestion that intensity peak A is a self-luminous near-infrared source.
What is the nature of the self-luminous source at peak A? Is it a deeply embedded star or a blob of warm dust? If it is a blob of warm dust, the only likely heat sources would be illumination from the former AGB star located at least 550 AU distant or shock heating. To produce significant thermal emission at 1.65 $`\mu `$m, the wavelength at which the blob begins to appear (Sahai et al. (1998)), would require dust with temperatures of at least 1000 K. The heating of a large amount of dust when the heat source is at least 550 AU away is highly unlikely, even for an AGB star with a luminosity of 10<sup>4</sup> L. In addition, some of the luminosity of the AGB star would be expected to show up in a reflection pattern at peak A, which we do not see. As for shock heating, the maps of H<sub>2</sub> emission (Fig. 6; also, see Sahai et al. (1998)) reveal no evidence of shocked gas within a few tenths of an arcsec (several hundred AU) of peak A. If the dust had been heated by a passing shock that is now 200 AU away, having moved past at 30 km s<sup>-1</sup>, it would have had at least 30 years to cool down. Thus, it appears more likely that intensity peak A is a star and that A and B most likely constitute a widely spaced, binary star system.
Assuming A and B are a binary, their minimum separation is 550 AU (taking d = 1 kpc). If A and B are both in the equatorial plane and the polar axis is tilted 15 out of the plane of the sky (Sahai et al. (1998) estimated a tilt of 10–20 from the axial ratio of the dust torus), then the star at A would lie $``$900 AU more distant than the star at B, making the true binary separation about 1000 AU. This separation is several orders of magnitude larger than that hypothesized (Morris (1987); Soker (1998)) for a central binary system that could trigger the formation of an equatorial disk and the consequent bipolar outflow.
It is remarkable that position B appears to be equidistant and point-symmetrically placed between the apex of the western loop (E1), the apex of the middle of the eastern loops (E3), and the most distant points in the polar lobes of molecular hydrogen emission (Fig. 6). Thus, the polarimetric and molecular hydrogen emission centroids are positionally coincident. This result strongly indicates that the nebular illuminator at B also generated the H<sub>2</sub> emission, where the H<sub>2</sub> emission regions are delineated by sharp outer boundaries suggestive of shocks. As shocks require fairly sudden changes — in this case, perhaps the rapid turning-on of a fast wind from the former AGB star, perhaps triggered by the quite quick stripping and ejection of the stellar envelope and the subsequent capture of a close companion — the relationship between position B and the H<sub>2</sub> emission lobes suggests that the shocks seen in the H<sub>2</sub> were caused by a very sudden event or series of events in the evolution of the central star.
Thus, while the presence of the A+B binary at the core of RAFGL 2688 does not lend support to the binary trigger hypothesis for the formation of bipolar planetary nebulae, the relationship between the central star at B and the H<sub>2</sub> lobes may support such a hypothesis. Specifically, absorption of a close binary companion by the atmosphere of the central AGB star may cause the ejection of high-velocity material; the ejected material produces the shocked H<sub>2</sub> emission and generates the bipolar structure of the Egg Nebula.
## 6 Summary
From a detailed analysis of the polarimetric images obtained using NICMOS and the HST, we have precisely determined the position of the post-AGB star in the waist of the Egg Nebula and the projected orientation of the polar axis (PA 12) of this bipolar system. This post-AGB star, which illuminates the Egg Nebula, falls point-symmetrically at the center of the molecular hydrogen emission regions that mark the waist and the polar lobes of the nebula. We find that this star lies 550 AU in projected distance, and perhaps 1000 AU in physical distance, from the star previously identified (Sahai et al. (1998)) at the southern tip of the northern polar lobe. Thus, these data provide clear evidence for the presence of an optically obscured, widely spaced binary system near the core of the bipolar, pre-planetary nebula RAFGL 2688. However, the separation between these components is orders of magnitude larger than required by models postulating that companions to AGB stars trigger the production of bipolar planetary nebulae.
DCH acknowledges support by NASA grant NAG5-3042 to the NICMOS instrument definition team. RS thanks NASA for support through grant GO-07423.01-96A from the Space Telescope Science Institute (which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555). |
warning/0002/hep-ph0002088.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Relatively little work has been done on heavy baryons using lattice QCD. The most complete study in a relativistic framework is the one done in the UKQCD collaboration for both charm and bottom flavoured baryons using an $`O(a)`$ improved fermion action in quenched approximation. In that work hyperfine splittings in both charm and bottom baryons were found to be considerably smaller than those predicted by phenomenological models and observed experimentally. A recent NRQCD calculation was able to resolve the hyperfine splittings for baryons with b quarks and found values in accord with phenomenological expectations. Since the hyperfine splitting is an important feature of the baryon spectrum further investigation of charmed baryons seems justified.
In this work we present the results of another quenched lattice simulation for charmed baryons. Due to limitations in computing resources it was not possible to use the same lattice spacing and volume as used in the UKQCD calculation. Rather we work on a more coarse lattice ($``$0.2fm) with a highly improved action. Past experience has shown that results of reasonable accuracy may be obtained with such lattices. In order to check the calculation, the spectrum of baryons in the light quark (u,d,s) sector was calculated at the same time. As well, meson masses for both heavy and light quarks were calculated. The results of all these calculations are in reasonable agreement with experimental values and with the results obtained at small lattice spacing.
The simulation reported here differs from in two other respects. In the interpolating operators used for the baryons were taken to have a form which emphasizes the heavy quark symmetry. Secondly, the spin 1/2 and spin 3/2 $`\mathrm{\Sigma }`$-like baryons are interpolated by the same operator, a Rarita-Schwinger spin-3/2 field. As is well known, the correlator of such a field propagates both $`J=1/2`$ and $`J=3/2`$ states and it is the $`J=1/2`$ projection of this correlator that is identified with the $`\mathrm{\Sigma }`$-like baryons in .
The procedure used here is different. As has been done in the context of a QCD sum rule calculation, the interpolating fields used for heavy baryons are taken to have the same form as those used in the light quark sector. The $`J=1/2`$ $`\mathrm{\Sigma }`$ baryon has an interpolating operator which is distinct from that used for the $`J=3/2`$ $`\mathrm{\Sigma }^{}`$, just as the nucleon is usually interpolated by a field that is distinct from the $`\mathrm{\Delta }`$. No assumption is made about heavy quark symmetry. An advantage of this is that we can use the same procedure (and computer code) at all masses which provides some check on the results.
The correlation function used to extract the masses of $`\mathrm{\Sigma }`$’s is not obtained from the projection of the $`J=3/2`$ field’s correlator. In fact, it seems that the spin 1/2 projection of this correlator has very poor overlap with the $`J=1/2`$ ground state. Our finding is that the spin 1/2 projected correlator is very small and noisy compared to the correlator calculated directly using a spin 1/2 field.
The final conclusion of this study is that in quenched lattice QCD the hyperfine splittings for both singly and doubly charmed baryons are in reasonable agreement with phenomenological expectations. Indications are that the splittings may be overestimated compared to experiment which seems to be a common tendency for quenched QCD simulations of baryons at all quark masses.
## 2 Method
The calculation is done on an anisotropic lattice using the gauge field action
$`S_G(U)`$ $`=`$ $`\beta [c_{ps}{\displaystyle \underset{ps}{}}(1{\displaystyle \frac{1}{3}}ReTrU_{ps})+c_{rs}{\displaystyle \underset{rs}{}}(1{\displaystyle \frac{1}{3}}ReTrU_{rs})`$ (1)
$`+c_{pt}{\displaystyle \underset{pt}{}}(1{\displaystyle \frac{1}{3}}ReTrU_{pt})+c_{rst}{\displaystyle \underset{rst}{}}(1{\displaystyle \frac{1}{3}}ReTrU_{rst})`$
$`+c_{rts}{\displaystyle \underset{rts}{}}(1{\displaystyle \frac{1}{3}}ReTrU_{rts})]`$
where ps and rs denote spatial plaquettes and spatial planar $`2\times 1`$ rectangles respectively. The plaquettes lying in the temporal-spatial planes are denoted by pt while rectangles with the long side in a spatial(temporal) direction are labeled by rst(rts). The c coefficients incorporate the aspect ratio $`\xi =a_s/a_t`$ and gauge link renormalization factors $`u_s`$ and $`u_t`$. These renormalization factors are estimated using the link expectation value in Landau gauge.
The fermion action is of the anisotropic D234 type and has the form
$`S_F`$ $`=`$ $`{\displaystyle \underset{x,i}{}}(c_{1i}D_{1i}(x)+c_{2i}D_{2i}(x))+{\displaystyle \underset{x}{}}(c_{1t}D_{1t}(x)+c_{2t}D_{2t}(x))`$ (2)
$`+{\displaystyle \underset{x,i<j}{}}c_{0s}\overline{\psi }(x)\sigma _{ij}F_{ij}(x)\psi (x)+{\displaystyle \underset{x,i}{}}c_{0t}\overline{\psi }(x)\sigma _{0i}F_{0i}(x)\psi (x)`$
$`{\displaystyle \underset{x}{}}\overline{\psi }(x)\psi (x),`$
where
$$D_{1i}(x)=\overline{\psi }(x)(1\xi \gamma _i)U_i(x)\psi (x+\widehat{i})+\overline{\psi }(x+\widehat{i})(1+\xi \gamma _i)U_i^{}(x)\psi (x),$$
(3)
$$D_{1t}(x)=\overline{\psi }(x)(1\gamma _4)U_4(x)\psi (x+\widehat{t})+\overline{\psi }(x+\widehat{t})(1+\gamma _4)U_4^{}(x)\psi (x),$$
(4)
$`D_{2i}(x)`$ $`=`$ $`\overline{\psi }(x)(1\xi \gamma _i)U_i(x)U_i(x+\widehat{i})\psi (x+2\widehat{i})`$ (5)
$`+\overline{\psi }(x+2\widehat{i})(1+\xi \gamma _i)U_i^{}(x+\widehat{i})U_i^{}(x)\psi (x),`$
$`D_{2t}(x)`$ $`=`$ $`\overline{\psi }(x)(1\gamma _4)U_4(x)U_4(x+\widehat{t})\psi (x+2\widehat{t})`$ (6)
$`+\overline{\psi }(x+2\widehat{t})(1+\gamma _4)U_4^{}(x+\widehat{t})U_4^{}(x)\psi (x).`$
The c coefficients in the fermion action include the aspect ratio, link renormalization and the hopping parameter factors and are shown in Table 1.
Hadron masses are calculated from zero-momentum correlation functions in the usual way. For mesons the interpolating fields were just the standard ones. For baryons some discussion is needed since the procedure used here differs from that used in . Start from the light quark (u,d,s) sector. A common choice for the proton operator in terms of u and d quark fields is
$$ϵ^{abc}[u_a^TC\gamma _5d_b]u_c$$
(7)
where a,b,c are colour indices and Dirac indices have been suppressed. For $`\mathrm{\Delta }`$ the operator
$$\frac{1}{\sqrt{3}}ϵ^{abc}\left\{2[u_a^TC\gamma _\mu d_b]u_c+[u_a^TC\gamma _\mu u_b]d_c\right\}$$
(8)
is used. This choice of operators is not unique but it allows an easy generalization to other baryons. The interpolating operators for the strange hyperons $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }^{}`$ are obtained by the replacement $`ds`$ in (7) and (8) respectively. Similarly the interpolators for $`\mathrm{\Xi }`$ and $`\mathrm{\Xi }^{}`$ are constructed by the replacement $`us`$. The $`\mathrm{\Lambda }`$ hyperon is more of a problem. In the SU(3) flavour limit it would be natural to use the octet lambda
$$\mathrm{\Lambda }_8=\frac{1}{\sqrt{6}}ϵ^{abc}\left\{2[u_a^TC\gamma _5d_b]s_c+[u_a^TC\gamma _5s_b]d_c[d_a^TC\gamma _5s_b]u_c\right\}$$
(9)
since it is degenerate with the nucleon and $`\mathrm{\Delta }`$. However, since SU(3) flavour is broken, this choice is not compelling. For example, in a combination $`\mathrm{\Lambda }`$ with both SU(3) octet and singlet components was defined. In this work a “heavy” $`\mathrm{\Lambda }`$ with a form
$$[u_a^TC\gamma _5d_b]s_c,$$
(10)
which is natural in the heavy quark limit, is also used.
The operators used to calculate the masses of the eight ground state singly-charmed baryons are taken to have the same structure as the operators given above. The $`\mathrm{\Sigma }_c`$ and $`\mathrm{\Sigma }_c^{}`$ baryons are obtained by the substitution $`dc`$ in (7) and (8) while $`\mathrm{\Omega }_c`$ and $`\mathrm{\Omega }_c^{}`$ are obtained by $`us,`$ d$`c`$. For the $`\mathrm{\Lambda }_c`$ and $`\mathrm{\Xi }_c`$ the operators (9) and (10) with the replacements $`sc`$ and $`ds,`$ $`sc`$ are used. Finally, for the remaining two states we take
$$\mathrm{\Xi }__c^{}=\frac{1}{\sqrt{2}}ϵ^{abc}\left\{[u_a^TC\gamma _5c_b]s_c+[s_a^TC\gamma _5c_b]u_c\right\}$$
(11)
and
$$\mathrm{\Xi }_c^{}=\frac{2}{\sqrt{3}}ϵ^{abc}\left\{[u_a^TC\gamma _\mu s_b]c_c+[s_a^TC\gamma _\mu c_b]u_c+[c_a^TC\gamma _\mu u_b]s_c\right\}.$$
(12)
Rather than using the above operators which have an explicit relativistic form one could consider the operators which survive in the limit of a static charm quark. The way to do this has been discussed in some detail in the context of QCD sum rule calculations. However there is no particular advantage to taking this limit here. The operators we use contain the leading heavy-quark components so the simulation will decide by itself whether they are dominant. The advantage of using the explicit relativistic forms is that it allows a unified analysis of hyperfine effects over the whole mass range from nucleon and $`\mathrm{\Delta }`$ to charmed baryons.
The operators such as (8) and(12) propagate both spin 1/2 and spin 3/2 states. At zero momentum the correlation function with spatial Lorentz indices has the general form
$$C_{ij}(t)=(\delta _{ij}\frac{1}{3}\gamma _i\gamma _j)C_{3/2}(t)+\frac{1}{3}\gamma _i\gamma _jC_{1/2}(t)$$
(13)
where the subscripts 3/2 and 1/2 denoted the spin projections. The quantity $`C_{3/2}(t)`$ was used to extract the mass of the spin 3/2 states. However, it was found that the spin 1/2 projection $`C_{1/2}(t)`$ was too noisy at large time separations to allow for the determination of a mass.
Hadron correlators were calculated using interpolating operators in local form at both source and sink and also applying a gauge invariant smearing to the quark propagators at the sink. The Gaussian smearing function, eqn(13) of , was used. Hadron masses were obtained by a simultaneous fit to local and sink-smeared correlators.
## 3 Results
The calculations were carried out at $`\beta =2.1`$ on lattices with a bare aspect ratio $`\xi `$ of 2. The Landau link tadpole factors were first determined iteratively to be $`u_s=0.7858`$ and $`u_t=0.9472`$ and these values were used in all subsequent calculations. The static potential was determined from both spatial and spatial-temporal Wilson loops. From this the lattice spacing and renormalized anisotropy were obtained. The results for the lattice spacings are $`a_s^1=(0.977\pm 0.003)GeV`$ and $`a_t^1=(1.914\pm 0.017)GeV`$ with a systematic uncertainty of $`0.01GeV`$ coming from uncertainty in the choice of parametrization of the short distance part of the potential. The renormalized anisotropy was found to be $`1.95\pm 0.02`$ which is compatible with other studies done with improved gluon actions at similar lattice spacings.
Fermion propagators were calculated on a $`10^3\times 30`$ lattice with Dirichlet boundary conditions on the fermion fields. A total of 420 configurations were analyzed. With some preliminary tuning it was found that $`\kappa =0.182`$ and $`\kappa =0.237`$ gave good values for the $`J/\psi `$ and $`\varphi `$ meson masses so these were the $`\kappa `$ values adopted for the charm and strange quarks for all calculations. Where necessary, masses were extrapolated to the physical up and down quark region using results from the set of hopping parameter values $`0.229,0.233,0.237`$ and $`0.241.`$ The value of the critical $`\kappa `$ is $`0.2429(2).`$
First consider the light quark (u,d,s) sector. The pion and $`\rho `$-meson masses were fixed at 0.140GeV and 0.770GeV which determines the hopping parameter for up and down quarks (taken to be degenerate) and the lattice scale $`a_\rho `$. It was found that $`a_\rho ^1=(1.99\pm 0.12)GeV`$ which is slightly larger than $`a_t^1`$ found from the static potential. This is an inevitable result of the quenched approximation. The $`\rho `$-meson mass scale was used in all subsequent calculations. The results in the light quark sector are given in Table 2 with statistical errors obtained by a bootstrap procedure. The dominant systematic error, a 6% uncertainty in the scale determination, is not shown explicitly in this and subsequent tables but should be kept in mind. For comparison, results from recent calculations (Table II in and Table XVII in ) done at small lattice spacing and extrapolated to the continuum are also shown. The results of our coarse lattice simulation are seen to be quite reasonable.
In Table 2 the $`\mathrm{\Lambda }`$ mass calculated with the operator (9) is given. The mass obtained using (10) was essentially identical. This was found to be true for all quark masses.
As check on how well charmed quarks are being simulated we first show the results for charmonium and D-mesons in Table 3. The hyperfine splitting between $`\eta _c`$ and $`J/\psi `$ is $`78(2)(5)MeV`$ which is in very good agreement with the results obtained in using a completely different fermion action. It is considerably smaller than the experimental value which is a well known feature of quenched QCD simulations for quarkonium. The $`D^{}D`$ splittings are compatible with results from NRQCD on similar lattices and are also smaller than experimental values although the suppression of hyperfine effects is less pronounced than in charmonium.
As mentioned in Section 2 the interpolating operators used for charmed baryons are taken to have the same form as those used for the light baryons. For example, the correlator for $`\mathrm{\Sigma }_c`$ is calculated directly using operators which are the same as used for the strange $`\mathrm{\Sigma }`$ hyperon except the mass is increased to charm. An alternative is to extract the masses of $`\mathrm{\Sigma }`$-like baryons (e.g., $`\mathrm{\Sigma }_c,\mathrm{\Omega }_c`$) from the spin-1/2 projection of a correlation function of spin 3/2 fields( see (13)). Figure 1 illustrates why this alternative is not used here. The correlation functions for $`\mathrm{\Omega }_c`$ and the spin-3/2 projected $`\mathrm{\Omega }_c^{}`$ are plotted as a function of lattice time. Also shown is the spin-1/2 projected correlator. This spin-1/2 projected piece has a very fast pre-asymptotic falloff and is therefore small and noisy in the time region in which one would want to extract the mass. This shows that the overlap of the spin-1/2 projection is small. It is also worth noting from Fig. 1 that even without any analysis one sees that the $`\mathrm{\Omega }_c`$ correlator has a less rapid falloff than that of $`\mathrm{\Omega }_c^{}`$ i.e., $`\mathrm{\Omega }_c`$ is less massive than $`\mathrm{\Omega }_c^{}`$.
The results for singly charmed baryons are given in Table 4. Overall the results are in reasonable close to the experimental values where they are known.
The masses of doubly charmed baryons were also calculated. There are no experimental data but a comparison with a selection of model calculations is given in Table 5
(A more complete tabulation from various models may be found in .) As might expected, without experimental constraints, the model calculations vary over a considerable range. The mass splittings from our quenched QCD simulation, 84(13)(5)MeV for $`\mathrm{\Xi }_{cc}^{}\mathrm{\Xi }_{cc}`$ and 78(7)(5)MeV for $`\mathrm{\Omega }_{cc}^{}\mathrm{\Omega }_{cc}`$, are substantial and lie in the middle of the range covered by the models listed in Table 5. They are only slightly smaller than the hyperfine splittings found for singly charmed baryons.
## 4 Summary
Hadron masses were calculated with an improved action on an anisotropic lattice with a spatial lattice spacing of about 0.2fm. Comparison with simulations done at small lattice spacings and extrapolated to the continuum indicate that lattice spacing errors are less than 10%.
The focus of this study is charmed baryons. For both singly and doubly charmed baryons spin splittings were found to be in agreement with expectations of quark models and other phenomenological approaches. The splittings for singly charmed baryons are somewhat larger than experimental values. This is in contrast to the small hyperfine effects for singly charmed baryons reported by the UKQCD collaboration.
The calculations reported here were done on a small lattice at a relatively large lattice spacing. By doing a unified study for light and heavy quark masses and comparing to continuum results where possible we have some confidence that the correct qualitative pattern of hyperfine effects in charmed baryons has been established for quenched lattice QCD. For a precise calculation, finite volume and lattice spacing issues have to be addressed. It is hoped that this can be done in the near future.
## Acknowledgements
It is a pleasure to thank H.R. Fiebig, D.B. Leinweber, R. Lewis, N.H. Shakespeare and H.D. Trottier for help and discussion and B.K. Jennings for the use of his computers. This work is supported in part by the Natural Sciences and Engineering Research Council of Canada.
## Figure caption
Fig. 1 Correlation function for the $`\mathrm{\Omega }_c`$ field (squares) and the $`\mathrm{\Omega }_c^{}`$ field ( spin 3/2 projection (triangles), spin 1/2 projection (circles)) as a function of lattice time. |
warning/0002/hep-th0002065.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Over the last few years, following the pioneering ideas of Ghoshal and Zamolodchikov, and others , much work has been done to investigate integrable quantum field theory with a boundary. In particular, the affine Toda field theories have offered a surprisingly rich structure which is just beginning to be understood. The classical affine Toda field theories are known to remain integrable in the presence of certain (generally quite restricted) boundary conditions confining them to a half-line, or to an interval . The corresponding quantum field theories are hardly explored although there has been progress in certain cases associated with the $`a_n^{(1)}`$ class of models . The first of these ($`n=1`$) is the sinh-Gordon model and, unlike all the others has a set of integrable boundary conditions depending on a pair of extra parameters, the so-called boundary parameters. However, even in this case, it remains to be seen precisely how the two boundary parameters influence quantities of interest such as the reflection factor.
The sinh-Gordon model is related to the sine-Gordon model and its reflection factors are related to the reflection factors of the lightest breather in the sine-Gordon model. Indeed, the results obtained by Ghoshal and Zamolodchikov allow a determination of the general form of the sinh-Gordon reflection factor and this is a very useful piece of information. However, apart from two special cases (Neumann and Dirichlet boundary conditions) Ghoshal and Zamolodchikov’s formulae fail to provide a relationship between the reflection factors and the boundary parameters themselves which appear as the data in a Lagrangian formulation of the model.
Recently, it was noticed that for certain ranges of the boundary parameters in the sinh-Gordon model there are real periodic classical finite-energy solutions called boundary breathers. The sinh-Gordon model has no real finite-energy solutions at all on the whole line (other than $`\varphi =0`$) but, once there is a boundary singularities may be hidden behind it allowing solutions restricted to a half-line to have finite energy. The existence of periodic solutions allows a semi-classical quantization determining the spectrum of boundary bound states. Once the spectrum of bound state energies is known it may be compared with a boundary bootstrap calculation of the energies of the same states and hence may be used, in principle, to find a relationship between the Lagrangian boundary parameters and the data in Ghoshal’s formula. In fact this is complicated and was carried out in for the simpler situation in which the two boundary parameters are equal and the $`\varphi \varphi `$ symmetry of the bulk theory is preserved. The corresponding calculation in the more general case is not yet completed. The semi-classical calculation may turn out to be exact (as it did for the energy spectrum of the breathers in the bulk sine-Gordon model itself; for a review see ) but, failing a proof of exactness the results may be checked against low-order perturbation theory .
If the two boundary parameters are different and the $`\varphi \varphi `$ symmetry is broken, the perturbation theory becomes substantially more complicated: the lowest energy static background is no longer the configuration $`\varphi =0`$ and therefore the perturbation theory must be developed within a non-trivial background, leading to additional cubic and higher odd order couplings as well as a substantially more intricate propagator. Nevertheless, as will be shown below, provided calculations are restricted to first order in the difference of the two boundary parameters, some of the complications disappear and we are able to calculate the correction to the reflection factor at one loop. The result we obtain is consistent with a conjecture for the relationship between the boundary parameters and Ghoshal’s formula in the more general setting. A full discussion, to all orders in the difference of the boundary parameters, even at one loop order in the bulk coupling, must be postponed.
## 2 sinh-Gordon model
The sinh-Gordon theory corresponds to the affine Toda field theory associated with $`a_1^{(1)}`$. The bulk Lagrangian density of the theory is:
$$=\frac{1}{2}_\mu \varphi ^\mu \varphi V(\varphi )$$
(2.1)
where
$$V(\varphi )=\frac{2m^2}{\beta ^2}\mathrm{cosh}(\beta \alpha \varphi )$$
(2.2)
The real constants $`m`$ and $`\beta `$ provide a mass scale and a coupling constant, respectively, and it is customary in affine Toda field theory to choose $`\alpha =\sqrt{2}`$.
Assuming two sinh-Gordon particles of relative rapidity $`\mathrm{\Theta }`$ scatter from each other elastically, the bulk S-matrix characterizing this process is given by
$$S(\mathrm{\Theta })=\frac{1}{(B)(2B)}$$
(2.3)
where
$$B=\frac{1}{2\pi }\frac{\beta ^2}{1+\beta ^2/4\pi }$$
(2.4)
and the symbol $`()`$ denotes the hyperbolic building block:
$$(x)=\frac{\mathrm{sinh}(\mathrm{\Theta }/2+\frac{i\pi x}{4})}{\mathrm{sinh}(\mathrm{\Theta }/2\frac{i\pi x}{4})}$$
(2.5)
The S-matrix is invariant under the following weak-strong coupling transformation
$$\beta 4\pi /\beta $$
(2.6)
a property known as weak-strong coupling duality.
A sinh-Gordon theory on the half-line is described by the following Lagrangian density:
$$\overline{}=\theta (x)\delta (x),$$
(2.7)
where $``$ is regarded as a functional of the field only, not its time derivative. Moreover, the generic form of $``$ is given by
$$=\frac{m}{\beta ^2}\left(\sigma _0e^{\frac{\beta }{\sqrt{2}}\varphi }+\sigma _1e^{\frac{\beta }{\sqrt{2}}\varphi }\right)$$
(2.8)
Note, the coefficients $`\sigma _0`$ and $`\sigma _1`$ are a pair of real numbers, essentially free (but see ), which represent the extra parameters permitted at the boundary $`x=0`$.
The equation of motion and the boundary condition for the sinh-Gordon model on the half-line become (after rescaling the mass scale to unity):
$$^2\varphi =\frac{\sqrt{2}}{\beta }\left(e^{\sqrt{2}\beta \varphi }e^{\sqrt{2}\beta \varphi }\right)x<0$$
(2.9)
$$\frac{\varphi }{x}=\frac{\sqrt{2}}{\beta }\left(\sigma _1e^{\beta \varphi /\sqrt{2}}\sigma _0e^{\beta \varphi /\sqrt{2}}\right)x=0$$
(2.10)
## 3 Reflection factor
When a sinh-Gordon particle is moving towards the boundary located at $`x=0`$ it will reflect elastically from it meaning that the in- and out- one particle states, conveniently labelled by rapidity, will be related by a reflection factor
$$|\theta >_{\mathrm{out}}=K(\theta )|\theta >_{\mathrm{in}}.$$
(3.1)
The general form of this reflection factor is known although its detailed dependence on the boundary parameters $`\sigma _0`$ and $`\sigma _1`$ appearing in (2.8) is known only in special cases. The sequence of arguments determining its form is somewhat indirect, stemming from the work of Ghoshal and Zamolodchikov as follows. Solving the boundary Yang-Baxter equation, and using general constraints implementing unitarity and a form of crossing symmetry, it proved possible to calculate the reflection factor for the sine-Gordon soliton. Since breathers are soliton-anti-soliton bound states, a subsequent set of calculations using the boundary bootstrap led Ghoshal to conjecture reflection factors for each member of the full tower of breathers. Finally, the reflection factor for the lightest breather is supposed to be identical with that of the sinh-Gordon particle provided the sine-Gordon coupling $`\beta `$ is replaced by $`i\beta `$.
Thus, suitably transformed in the manner described, Ghoshal’s formula for the sinh-Gordon reflection factor is given by:
$$K_q(\theta )=\frac{(1)(2B/2)(1+B/2)}{(1E(\sigma _0,\sigma _1,\beta ))(1+E(\sigma _0,\sigma _1,\beta ))(1F(\sigma _0,\sigma _1,\beta ))(1+F(\sigma _0,\sigma _1,\beta ))}$$
(3.2)
Actually, Ghoshal’s notation was a little different and made use of two other quantities $`\eta ,\vartheta `$ defined by $`E=B\eta /\pi ,F=iB\vartheta /\pi `$.
There are special cases in which $`E`$ and $`F`$ have been conjectured. For example, the Neumann boundary condition which is defined by
$$\frac{\varphi }{x}=0\text{when}x=0,$$
(3.3)
has been argued by Ghoshal and Zamolodchikov to need a reflection factor containing
$$F=0,E=1B/2.$$
(3.4)
More recently , the boundary breather states of the sinh-Gordon model restricted to a half-line were investigated and their energy spectrum calculated in two ways. First, by using the bootstrap equations, and then by finding a set of periodic finite-energy solutions which could be quantized using a WKB approximation. Insisting that the two methods agreed with each other provided strong evidence for a relationship between the boundary parameters, the bulk coupling constant, and the parameters appearing in the reflection factor. For technical reasons this work was restricted to the case $`\sigma _0=\sigma _1`$ but, nevertheless, yielded an expression for $`E`$, with $`F=0`$. Specifically, setting $`\sigma _0=\sigma _1=\mathrm{cos}a\pi `$ and restricting $`a`$ to the range $`1>a>1/2`$, $`E`$ is given by,
$$E=2a(1B/2).$$
(3.5)
In the limit $`a1/2`$ from above, this is in agreement with (3.4).
Both of these conjectures are underpinned by low order perturbation theory. Kim has calculated the one loop correction to the classical Neumann boundary condition reflection factor and found agreement with Ghoshal’s formula to $`O(\beta ^2)`$. On the other hand, a more general calculation was carried out, also to one loop, agreeing with the conjecture (3.5), and indeed, preceding it . In fact, the perturbative calculation in agrees with (3.5) for $`a`$ in the range $`1>a>0`$.
However, when $`\sigma _0\sigma _1`$ the situation becomes much more complicated. Indeed, all that is known up to now is the behaviour of $`E`$ and $`F`$ in the limit $`\beta 0`$, deduced from a direct calculation of the classical reflection factor in the general case . To describe this limit it is convenient to set $`\sigma _0=\mathrm{cos}a_0\pi `$ and $`\sigma _1=\mathrm{cos}a_1\pi `$, with $`|a_i|1,i=0,1`$. Then,
$$E(0,\sigma _0,\sigma _1)=a_0+a_1,F(0,\sigma _0,\sigma _1)=a_0a_1.$$
(3.6)
and, the classical reflection factor itself is given in term of the basic factors (2.5) by
$$K_c(\theta )=\frac{(1)^2}{(1a_0a_1)(1+a_0+a_1)(1a_0+a_1)(1+a_0a_1)}.$$
(3.7)
The expression (3.7) is an essential ingredient of the basic two-point function, or propagator, which takes the form ,
$`G(x,t;x^{},t^{})`$ $`=`$ $`{\displaystyle \frac{d\omega }{2\pi }\frac{dk}{2\pi }\frac{ie^{i\omega (tt^{})}}{(\omega ^2k^24+i\rho )}f(k,x^{})}`$ (3.8)
$`\times \left(f(k,x)e^{ik(xx^{})}+K_cf(k,x)e^{ik(x+x^{})}\right)`$
where
$$f(k,x)=\frac{ik2\mathrm{coth}2(xx_0)}{ik+2}.$$
(3.9)
The parameter $`x_0`$ enters the static background
$$e^{\beta \varphi _0/\sqrt{2}}=\frac{1+e^{2(xx_0)}}{1e^{2(xx_0)}}$$
(3.10)
and determines the point at which the background becomes singular. For this reason, it is crucial that $`x_00`$. Actually, $`x_0`$ is determined by the boundary condition (2.10) and satisfies
$$\mathrm{coth}x_0=\sqrt{\frac{1+\sigma _0}{1+\sigma _1}}.$$
(3.11)
On the understanding that $`\sigma _1\sigma _0`$, it is enough to consider this situation since (2.10) is symmetric under $`\varphi \varphi `$ and the interchange of $`\sigma _0`$ with $`\sigma _1`$.
In the expression (3.8) for the propagator, the classical reflection factor appears as the coefficient of the reflection part of the free field two-point function calculated within the classical static background. To check Ghoshal’s formula, the strategy introduced by Kim and developed in is to calculate perturbative corrections to the two-point function and then to identify corrections to the classical reflection factor by picking out the coefficient of $`e^{ik(x+x^{})}\text{as}x,x^{}\mathrm{}.`$
## 4 First order quantum corrections to the reflection factor
Previously, as noted above, the first order perturbation calculation was achieved for the specially symmetric case $`\sigma _0=\sigma _1`$. In this section, this restriction will be relaxed and enough of the ingredients of the perturbation expansion will be calculated to allow a calculation (in principle) up to $`\beta ^2`$. We will discover that generally the theory needs three- and four- point couplings which depend upon the $`x`$-dependent static background.
Expanding the bulk potential (2.2) around the background solution to the equation of motion, $`\varphi _0(x)`$, we derive the three- and four-point couplings:
$$C_{\mathrm{bulk}}^{(3)}=\frac{2\sqrt{2}}{3}\beta \mathrm{sinh}(\sqrt{2}\beta \varphi _0),$$
(4.1)
and
$$C_{\mathrm{bulk}}^{(4)}=\frac{1}{3}\beta ^2\mathrm{cosh}(\sqrt{2}\beta \varphi _0).$$
(4.2)
On the other hand, the static background is (3.10). So, after some manipulation, we have
$$C_{\mathrm{bulk}}^{(3)}=\frac{4\sqrt{2}}{3}\beta \mathrm{cosh}2(xx_0)\left(\mathrm{coth}^22(xx_0)1\right),$$
(4.3)
and
$$C_{\mathrm{bulk}}^{(4)}=\frac{1}{3}\beta ^2\left(2\mathrm{coth}^22(xx_0)1\right).$$
(4.4)
In the same manner, using (2.8) we can derive the three- and four-point couplings associated with the boundary term:
$$C_{\mathrm{boundary}}^{(3)}=\frac{\sqrt{2}\beta }{12}\left(\sigma _1e^{\beta \varphi _0/\sqrt{2}}\sigma _0e^{\beta \varphi _0/\sqrt{2}}\right),$$
(4.5)
and
$$C_{boundary}^{(4)}=\frac{\beta ^2}{48}\left(\sigma _1e^{\beta \varphi _0/\sqrt{2}}+\sigma _0e^{\beta \varphi _0/\sqrt{2}}\right).$$
(4.6)
For future reference, it will be useful to know these to first order in the difference of the two boundary parameters $`ϵ=\sigma _0\sigma _1`$:
$$C_{\mathrm{bulk}}^{(3)}=\frac{2\sqrt{2}}{3}\beta \frac{ϵ}{1+\sigma _1}e^{2x}+\mathrm{}$$
(4.7)
$$C_{\mathrm{bulk}}^{(4)}=\frac{1}{3}\beta ^2+\mathrm{}$$
(4.8)
and similarly,
$$C_{\mathrm{boundary}}^{(3)}=\frac{\sqrt{2}\beta }{12}\left(\frac{ϵ}{1+\sigma _1}\right)+\mathrm{}$$
(4.9)
$$C_{\mathrm{boundary}}^{(4)}=\frac{\beta ^2}{48}(2\sigma _1+ϵ)+\mathrm{}$$
(4.10)
We shall also need the expressions for $`f(k,x)`$ (3.9) and the classical reflection factor $`K_c`$ (3.7) to the same order in $`ϵ`$ since both contribute to the propagator (3.8). In fact,
$$f(k,x)=1+O(ϵ^2)$$
(4.11)
and the classical reflection factor (3.7) reduces to
$$K_c=\frac{ik+2\sigma }{ik2\sigma }+\frac{2ik}{(ik2\sigma )^2}ϵ+O(ϵ^2).$$
(4.12)
It is also convenient to write
$$K_c=K_0+ϵK_1,$$
(4.13)
where $`K_0`$ is the classical reflection factor when the two boundary parameters are equal.
To calculate quantum corrections to the classical reflection factor at one loop order we use perturbative methods generalised to the affine Toda field theory on a half-line . (For earlier references on boundary perturbation theory in general see; for affine Toda perturbation theory see , or the review .) The $`O(\beta ^2)`$ correction to $`K_0`$ has been calculated before and the purpose of this article is to calculate the corrections to $`K_1`$ to the same order.
The possible diagrams to $`O(\beta ^2)`$ are:
I II III
Figure 1: Three basic Feynman diagrams in one loop order.
These will be computed in configuration space noting that each vertex may either be situated at the boundary or within the bulk. In effect, there are ten contributions which need to be calculated.
However, there is a simplifying feature provided we are content to work to first order in $`ϵ`$. To recognise this it is enough to note that because the three-point coupling is already $`O(ϵ)`$, implying the type II and III diagrams involve $`ϵ^2`$, only diagrams of type I need concern us.
Thus, there are two contributions which need to be calculated in order to be able to deduce the quantum corrections to the classical reflection factor, both of type I. The first is directly related to the boundary, when the interaction vertex coincides with the boundary, and it takes the form
$$\frac{i\beta ^2}{4}(2\sigma _1+ϵ)_{\mathrm{}}^+\mathrm{}𝑑t^{\prime \prime }G(x,t;0,t^{\prime \prime })G(0,t^{\prime \prime };0,t^{\prime \prime }),G(0,t^{\prime \prime };x^{},t^{}),$$
(4.14)
including the correct coupling constant and combinatorial factors.
The second contribution refers to the bulk potential which means the interaction vertex is located inside the bulk region $`x<0`$. This contribution is given by
$$4i\beta ^2_{\mathrm{}}^+\mathrm{}𝑑t^{\prime \prime }_{\mathrm{}}^0𝑑x^{\prime \prime }G(x,t;x^{\prime \prime },t^{\prime \prime })G(x^{\prime \prime },t^{\prime \prime };x^{\prime \prime },t^{\prime \prime })G(x^{\prime \prime },t^{\prime \prime };x^{},t^{}).$$
(4.15)
Again, the combinatorial factor has been included.
Let us first calculate the boundary contribution (4.14). The loop propagator is given to $`O(ϵ)`$ by
$$G(0,t^{\prime \prime };0,t^{\prime \prime })=i\frac{d\omega ^{\prime \prime }}{2\pi }\frac{dk^{\prime \prime }}{2\pi }P_0(\omega ^{\prime \prime },k^{\prime \prime })\left(1+\frac{ik^{\prime \prime }+2\sigma _1}{ik^{\prime \prime }2\sigma _1}+\frac{2ik^{\prime \prime }ϵ}{(ik^{\prime \prime }2\sigma _1)^2}\right),$$
(4.16)
where we have defined
$$P_0(\omega ^{\prime \prime },k^{\prime \prime })=\frac{1}{\omega ^{\prime \prime 2}k^{\prime \prime 2}4+i\rho }.$$
(4.17)
Note, the integral is clearly divergent but the divergence is removed by a renormalization of the boundary parameter. In effect, rearranging the part of the integrand containing the offending terms as follows,
$$1+\frac{ik^{\prime \prime }+2\sigma _1}{ik^{\prime \prime }2\sigma _1}=2+\frac{4\sigma }{ik^{\prime \prime }2\sigma _1}$$
(4.18)
and making a minimal subtraction of the divergent portion, (i.e. deleting the ‘2’), renders the integral finite.
The part of the integral independent of $`ϵ`$ has been calculated before , therefore we may write,
$$G(0,t^{\prime \prime };0,t^{\prime \prime })=\frac{a_1\mathrm{cos}a_1\pi }{\mathrm{sin}a_1\pi }+i\frac{d\omega ^{\prime \prime }}{2\pi }\frac{dk^{\prime \prime }}{2\pi }P_0(\omega ^{\prime \prime },k^{\prime \prime })\frac{2ik^{\prime \prime }ϵ}{(ik^{\prime \prime }2\sigma _1)^2},$$
(4.19)
and concentrate on the $`O(ϵ)`$ part.
As far as the remaining integral is concerned, focusing on the energy variable and closing the integration contour in the upper half-plane, we encounter a simple pole at $`\sqrt{k^{\prime \prime 2}+4}`$. Therefore, we need to evaluate
$$\frac{1}{2}\frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}\frac{2ik^{\prime \prime }}{(ik^{\prime \prime }2\sigma _1)^2}ϵ.$$
(4.20)
The $`k^{\prime \prime }`$ integration may be performed by closing the contour in the upper (lower) half-plane depending on whether $`\sigma _1>0`$ ($`\sigma _1<0`$). However, because of the branch points at $`\pm 2i`$ it is convenient to locate the associated branch cuts along the imaginary axis from $`\pm 2i`$ to $`\mathrm{}`$. Evaluating the integral and assembling the pieces, we obtain
$$i\frac{d\omega ^{\prime \prime }}{2\pi }\frac{dk^{\prime \prime }}{2\pi }P_0(\omega ^{\prime \prime },k^{\prime \prime })\frac{2ik^{\prime \prime }}{(ik^{\prime \prime }2\sigma _1)^2}ϵ=\frac{ϵ}{2}\frac{a_1}{\mathrm{sin}^3a_1\pi }+\frac{ϵ\mathrm{cos}a_1\pi }{2\pi }\frac{1}{\mathrm{sin}^2a_1\pi }.$$
(4.21)
At this stage, the contribution to (4.14) is
$`i{\displaystyle \frac{\beta ^2}{4}}(2\sigma _1+ϵ)\left({\displaystyle \frac{a_1\mathrm{cos}a_1\pi }{\mathrm{sin}a_1\pi }}{\displaystyle \frac{ϵ}{2}}{\displaystyle \frac{a_1}{\mathrm{sin}^3a_1\pi }}+{\displaystyle \frac{ϵ\mathrm{cos}a_1\pi }{2\pi }}{\displaystyle \frac{1}{\mathrm{sin}^2a_1\pi }}\right)`$ (4.22)
$`\times `$ $`i{\displaystyle 𝑑t^{\prime \prime }\frac{d\omega }{2\pi }\frac{dk}{2\pi }e^{i\omega (tt^{\prime \prime })}P_0(\omega ,k)e^{ikx}\left(\frac{2ik}{ik2\sigma _1}+\frac{2ikϵ}{(ik2\sigma _1)^2}\right)}`$
$`\times `$ $`i{\displaystyle \frac{d\omega ^{}}{2\pi }\frac{dk^{}}{2\pi }e^{i\omega ^{}(t^{\prime \prime }t^{})}P_0(\omega ^{},k^{})e^{ik^{}x^{}}\left(\frac{2ik^{}}{ik^{}2\sigma _1}+\frac{2ik^{}ϵ}{(ik^{}2\sigma _1)^2}\right)}`$
The integration over $`t^{\prime \prime }`$ ensures energy conservation at the interaction vertex and creates a Dirac delta function which immediately removes one of the energy variables, for example $`\omega ^{}`$. The remaining integral over the momenta $`k`$ and $`k^{}`$ can be performed by completing the contours in the upper half-plane and taking into account the poles at $`k=k^{}=\sqrt{\omega ^24}\widehat{k}`$. However, if $`\sigma _1>0`$ it is evident that the expressions for $`K_0`$ and $`K_1`$ have no poles inside the contour. If $`\sigma _1<0`$, there is an additional pole but its contribution turns out to be exponentially decreasing in the asymptotic region $`x,x^{}\mathrm{}`$.
Finally, we obtain the boundary contribution (4.14) in the form
$`i{\displaystyle \frac{\beta ^2}{4}}{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega (tt^{})}e^{i\widehat{k}(x+x^{})}\{{\displaystyle \frac{2a_1\mathrm{cos}^2a_1\pi }{\mathrm{sin}a_1\pi }}{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^2}}`$
$`+\left({\displaystyle \frac{a_1\mathrm{cos}a_1\pi }{\mathrm{sin}a_1\pi }}+{\displaystyle \frac{a_1\mathrm{cos}a_1\pi }{\mathrm{sin}^3a_1\pi }}{\displaystyle \frac{\mathrm{cos}^2a_1\pi }{\pi \mathrm{sin}^2a_1\pi }}\right){\displaystyle \frac{ϵ}{(i\widehat{k}2\sigma _1)^2}}`$
$`+{\displaystyle \frac{4a_1\mathrm{cos}^2a_1\pi }{\mathrm{sin}a_1\pi }}{\displaystyle \frac{ϵ}{(i\widehat{k}2\sigma _1)^3}}\},`$ (4.23)
where $`\widehat{k}=2\mathrm{sinh}\theta `$.
Next, we need to calculate the contribution (4.15) which to $`O(ϵ)`$ is:
$`4i\beta ^2{\displaystyle 𝑑t^{\prime \prime }_{\mathrm{}}^0𝑑x^{\prime \prime }\frac{d\omega }{2\pi }\frac{dk}{2\pi }e^{i\omega (tt^{\prime \prime })}iP_0(\omega ,k)}`$
$`\left(e^{ik(xx^{\prime \prime })}+{\displaystyle \frac{ik+2\sigma _1}{ik2\sigma _1}}e^{ik(x+x^{\prime \prime })}+{\displaystyle \frac{2ikϵ}{(ik2\sigma _1)^2}}e^{ik(x+x^{\prime \prime })}\right)`$
$`{\displaystyle \frac{d\omega ^{\prime \prime }}{2\pi }\frac{dk^{\prime \prime }}{2\pi }iP_0(\omega ^{\prime \prime },k^{\prime \prime })\left(1+\frac{ik^{\prime \prime }+2\sigma _1}{ik^{\prime \prime }2\sigma _1}e^{2ik^{\prime \prime }x^{\prime \prime }}+\frac{2ik^{\prime \prime }ϵ}{(ik^{\prime \prime }2\sigma _1)^2}e^{2ik^{\prime \prime }x^{\prime \prime }}\right)}`$
$`{\displaystyle \frac{d\omega ^{}}{2\pi }\frac{dk^{}}{2\pi }iP_0(\omega ^{},k^{})e^{i\omega ^{}(t^{\prime \prime }t^{})}}`$
$`\left(e^{ik^{}(x^{\prime \prime }x^{})}+{\displaystyle \frac{ik^{}+2\sigma _1}{ik^{}2\sigma _1}}e^{ik^{}(x^{\prime \prime }+x^{})}+{\displaystyle \frac{2ik^{}ϵ}{(ik^{}2\sigma _1)^2}}e^{ik^{}(x^{\prime \prime }+x^{})}\right).`$ (4.24)
The integral over $`t^{\prime \prime }`$ yields a delta function which replaces $`\omega ^{}`$ by $`\omega `$. Furthermore, to calculate the integration over $`x^{\prime \prime }`$, it is convenient to use the following device
$$_{\mathrm{}}^0𝑑x^{\prime \prime }e^{ikx^{\prime \prime }+\tau x^{\prime \prime }}=\frac{i}{ki\tau }$$
(4.25)
where the small positive quantity $`\tau `$ will be taken to zero at the final stage of the calculations.
The loop integral which corresponds to the middle propagator of (4.15), is obviously logarithmically divergent and this divergence will be removed by the infinite renormalization of the mass parameter in the bulk potential. So, after making the subtraction and integrating $`x^{\prime \prime }`$ and $`\omega ^{\prime \prime }`$, and as before concentrating on the $`O(ϵ)`$ piece, we obtain the contribution
$`{\displaystyle \frac{i}{2}}{\displaystyle \frac{d\omega }{2\pi }\frac{dk}{2\pi }\frac{dk^{}}{2\pi }e^{ı\omega (tt^{})}e^{i(kx+k^{}x^{})}iP_0(\omega ,k)iP_0(\omega ,k^{})}`$
$`{\displaystyle }{\displaystyle \frac{dk^{\prime \prime }}{2\pi }}{\displaystyle \frac{1}{\sqrt{k^{\prime \prime 2}+4}}}\{{\displaystyle \frac{2ikϵ}{(ik2\sigma _1)^2}}{\displaystyle \frac{ik^{\prime \prime }+2\sigma _1}{ik^{\prime \prime }2\sigma _1}}({\displaystyle \frac{1}{k+k^{}2k^{\prime \prime }i\tau }}`$
$`{\displaystyle \frac{1}{k+k^{}+2k^{\prime \prime }+i\tau }}{\displaystyle \frac{ik^{}+2\sigma _1}{ik^{}2\sigma _1}})`$
$`+{\displaystyle \frac{2ik^{}ϵ}{(ik^{}2\sigma _1)^2}}{\displaystyle \frac{ik^{\prime \prime }+2\sigma _1}{ik^{\prime \prime }2\sigma _1}}\left({\displaystyle \frac{1}{k^{}+k2k^{\prime \prime }i\tau }}{\displaystyle \frac{1}{k^{}+k+2k^{\prime \prime }+i\tau }}{\displaystyle \frac{ik+2\sigma _1}{ik2\sigma _1}}\right)`$
$`+{\displaystyle \frac{2ik^{\prime \prime }ϵ}{(ik^{\prime \prime }2\sigma _1)^2}}({\displaystyle \frac{1}{k+k^{}2k^{\prime \prime }i\tau }}+{\displaystyle \frac{1}{kk^{}2k^{\prime \prime }i\tau }}{\displaystyle \frac{ik^{}+2\sigma _1}{ik^{}2\sigma _1}}`$
$`{\displaystyle \frac{1}{kk^{}+2k^{\prime \prime }+i\tau }}{\displaystyle \frac{ik+2\sigma _1}{ik2\sigma _1}}{\displaystyle \frac{1}{k+k^{}+2k^{\prime \prime }+i\tau }}{\displaystyle \frac{ik+2\sigma _1}{ik2\sigma _1}}{\displaystyle \frac{ik^{}+2\sigma _1}{ik^{}2\sigma _1}})\}.`$ (4.26)
In order to evaluate the integral over $`k^{\prime \prime }`$, we encounter the following two types:
$$\frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}\left(\frac{ik^{\prime \prime }+2\sigma _1}{ik^{\prime \prime }2\sigma _1}\right)\frac{1}{(k+k^{}2k^{\prime \prime }i\tau )}$$
(4.27)
and
$$\frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}\frac{2ik^{\prime \prime }ϵ}{(ik^{\prime \prime }2\sigma _1)^2}\frac{1}{(k+k^{}2k^{\prime \prime }i\tau )}.$$
(4.28)
Both of these may be performed by closing an appropriate contour in the upper half-plane ensuring that it runs around the branch cut located from $`k^{\prime \prime }=2i`$ to infinity along the imaginary axis. Note, if $`\sigma _1>0`$, then there is no pole inside the contour; however, if $`\sigma _1>0`$, there is an extra pole but its residue integrated over $`k`$ and $`k^{}`$ will give a contribution vanishing in the limit $`x,x^{}`$ $``$ $`\mathrm{}`$.
For example, the integral (4.28) evaluates to
$`{\displaystyle \frac{dk^{\prime \prime }}{2\pi }\frac{1}{\sqrt{k^{\prime \prime 2}+4}}\frac{2ik^{\prime \prime }ϵ}{(ik^{\prime \prime }2\sigma _1)^2}\frac{1}{(k+k^{}2k^{\prime \prime }i\tau )}}={\displaystyle \frac{ϵ}{\pi }}{\displaystyle \frac{k+k^{}}{(ik+ik^{}4\sigma _1)^2}}{\displaystyle \frac{a_1\pi }{\mathrm{sin}a_1\pi }}`$
$`{\displaystyle \frac{ϵ}{\pi }}{\displaystyle \frac{i\sigma _1}{(4\sigma _1ikik^{})}}{\displaystyle \frac{1}{\mathrm{sin}^2a_1\pi }}+{\displaystyle \frac{ϵ}{\pi }}{\displaystyle \frac{i\sigma _1^2}{(4\sigma _1ikik^{})}}{\displaystyle \frac{a_1\pi }{\mathrm{sin}^3a_1\pi }}`$
$`+{\displaystyle \frac{ϵ}{\pi }}{\displaystyle \frac{2i(k+k^{})}{(ik+ik^{}4\sigma _1)^2}}{\displaystyle \frac{1}{\sqrt{\frac{(k+k^{})^2}{4}+4}}}\mathrm{ln}\left\{{\displaystyle \frac{1+\frac{i(k+k^{})}{4}+\frac{i}{2}\sqrt{\frac{(k+k^{})^2}{4}+4}}{1+\frac{i(k+k^{})}{4}\frac{i}{2}\sqrt{\frac{(k+k^{})^2}{4}+4}}}\right\}`$ (4.29)
Let us divide the bulk contribution (4) in two parts: $`B_1`$ containing integrals over $`k^{\prime \prime }`$ of type (4.27) and $`B_2`$ whose $`k^{\prime \prime }`$ integrations are of type (4.28). For both it is necessary, after performing the integration over $`\omega ^{\prime \prime }`$, to do the $`k`$ and $`k^{}`$ integrals via closing contours in the upper half-plane to pick up the poles at $`k`$ or $`k^{}=\sqrt{\omega ^24}=\widehat{k}`$. All other pole contributions lead to exponentially damped terms in the limit $`x,x^{}\mathrm{}`$.
After some manipulation, $`B_1`$ is found to be equal to
$`B_1`$ $`=`$ $`2\beta ^2{\displaystyle \frac{d\omega }{2\pi }e^{i\omega (tt^{})}e^{i\widehat{k}(x+x^{})}\frac{1}{(2\widehat{k})^2}\frac{2i\widehat{k}ϵ}{(i\widehat{k}2\sigma _1)^2}}`$ (4.30)
$`\left\{{\displaystyle \frac{i}{4}}+{\displaystyle \frac{ia_1}{\mathrm{sin}a_1\pi }}{\displaystyle \frac{i\widehat{k}}{i\widehat{k}2\sigma _1}}+{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{1}{\sqrt{\widehat{k}^2+4}}}\left({\displaystyle \frac{i\pi }{2}}\theta \right)\right\}.`$
Notice that $`B_1`$, in the last term inside the braces, depends on $`\theta `$ in a manner which potentially is very inconvenient for a comparison with Ghoshal’s formula. Fortunately, this term will be canceled by a matching term in $`B_2`$.
After somewhat lengthier calculations, $`B_2`$ is given by
$`B_2`$ $`=`$ $`2\beta ^2{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega (tt^{})}e^{i\widehat{k}(x+x^{})}{\displaystyle \frac{1}{(2\widehat{k})^2}}\{{\displaystyle \frac{2iϵ}{(i\widehat{k}2\sigma _1)^2}}{\displaystyle \frac{\sigma _1^2}{\pi \mathrm{sin}^2a_1\pi }}`$ (4.31)
$`+{\displaystyle \frac{2iϵ}{(i\widehat{k}2\sigma _1)^2}}{\displaystyle \frac{a_1\sigma _1^3}{\mathrm{sin}^3a_1\pi }}+{\displaystyle \frac{2iϵ\widehat{k}}{(i\widehat{k}2\sigma _1)^2}}{\displaystyle \frac{1}{\pi \sqrt{\widehat{k}^2+4}}}\theta `$
$`+{\displaystyle \frac{i\widehat{k}+2\sigma _1}{i\widehat{k}2\sigma _1}}({\displaystyle \frac{iϵ}{2\pi \mathrm{sin}^2a_1\pi }}+{\displaystyle \frac{iϵa_1\sigma _1}{2\mathrm{sin}^3a_1\pi }})\}.`$
Assembling the $`O(ϵ)`$ part of (4) with (4.30) and (4.31) we obtain,
$`{\displaystyle \frac{i\beta ^2ϵ}{2}}{\displaystyle }{\displaystyle \frac{d\omega }{2\pi }}e^{i\omega (tt^{})}e^{i\widehat{k}(x+x^{})}\{{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^2}}({\displaystyle \frac{1}{2\pi }}+{\displaystyle \frac{a_1\sigma _1}{2\mathrm{sin}a_1\pi }})`$
$`+{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^3}}(2a_1\mathrm{sin}a_1\pi )+{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^2}}{\displaystyle \frac{2}{\widehat{k}}}({\displaystyle \frac{i}{4}}+{\displaystyle \frac{i}{2\sqrt{\widehat{k}^2+4}}})\},`$ (4.32)
whence we can deduce the correction to the quantity $`K_1`$ in (4.13). Explicitly, we have,
$`\delta K_1`$ $`=`$ $`i\beta ^2ϵ\widehat{k}\{{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^2}}({\displaystyle \frac{1}{2\pi }}+{\displaystyle \frac{a_1\sigma _1}{2\mathrm{sin}a_1\pi }})+{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^3}}(2a_1\mathrm{sin}a_1\pi )`$ (4.33)
$`+{\displaystyle \frac{1}{(i\widehat{k}2\sigma _1)^2}}{\displaystyle \frac{2}{\widehat{k}}}({\displaystyle \frac{i}{4}}+{\displaystyle \frac{i}{2\sqrt{\widehat{k}^2+4}}})\}.`$
The correction to $`K_0`$ which was calculated before in is,
$`\delta K_0`$ $`=`$ $`{\displaystyle \frac{i\beta ^2}{8}}K_0(\widehat{k})\mathrm{sinh}\theta \{({\displaystyle \frac{1}{\mathrm{cosh}\theta +1}}{\displaystyle \frac{1}{\mathrm{cosh}\theta }})`$ (4.34)
$`+2a_1({\displaystyle \frac{1}{\mathrm{cosh}\theta \mathrm{sin}a_1\pi }}{\displaystyle \frac{1}{\mathrm{cosh}\theta +\mathrm{sin}a_1\pi }})\}.`$
This completes the collection of ingredients we need.
## 5 Comparison with Ghoshal’s formula
In this section, the corrections to the classical reflection factor calculated above will be compared with the formula of Ghoshal quoted in (3.2).
Using (4.13), the relative correction to the classical reflection factor $`K_c`$ is given in terms of the corrections $`\delta K_0`$ and $`\delta K_1`$ by
$$\frac{\delta K_c}{K_c}=K_0^1\delta K_0+ϵ\left(K_0^1\delta K_1K_1K_0^2\delta K_0\right).$$
(5.1)
Hence, using (4.33) and (4.34) we have,
$`{\displaystyle \frac{\delta K_c}{K_c}}`$ $`=`$ $`{\displaystyle \frac{i\beta ^2}{8}}\mathrm{sinh}\theta \{({\displaystyle \frac{1}{\mathrm{cosh}\theta +1}}{\displaystyle \frac{1}{\mathrm{cosh}\theta }})`$ (5.2)
$`+2a_1({\displaystyle \frac{1}{\mathrm{cosh}\theta \mathrm{sin}a_1\pi }}{\displaystyle \frac{1}{\mathrm{cosh}\theta +\mathrm{sin}a_1\pi }})\}`$
$`+{\displaystyle \frac{i\beta ^2ϵ\mathrm{sinh}\theta }{8\mathrm{sin}a_1\pi }}\{{\displaystyle \frac{1}{\pi }}({\displaystyle \frac{1}{\mathrm{cosh}\theta \mathrm{sin}a_1\pi }}{\displaystyle \frac{1}{\mathrm{cosh}\theta +\mathrm{sin}a_1\pi }})`$
$`+a_1\mathrm{cos}a_1\pi ({\displaystyle \frac{1}{(\mathrm{cosh}\theta \mathrm{sin}a_1\pi )^2}}+{\displaystyle \frac{1}{(\mathrm{cosh}\theta +\mathrm{sin}a_1\pi )^2}})\}.`$
On the other hand, Ghoshal’s formula (3.2) for the reflection factor up to one loop order is given by:
$$K_q(\theta )K_c(\theta )\left(1\frac{i\beta ^2}{8}\mathrm{sinh}\theta (\theta )\right),$$
(5.3)
where
$`(\theta )`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{cosh}\theta +1}}{\displaystyle \frac{1}{\mathrm{cosh}\theta }}`$ (5.4)
$`+{\displaystyle \frac{e_1}{\mathrm{cosh}\theta +\mathrm{sin}(e_0\pi /2)}}{\displaystyle \frac{e_1}{\mathrm{cosh}\theta \mathrm{sin}(e_0\pi /2)}}`$
$`+{\displaystyle \frac{f_1}{\mathrm{cosh}\theta +\mathrm{sin}(f_0\pi /2)}}{\displaystyle \frac{f_1}{\mathrm{cosh}\theta \mathrm{sin}(f_0\pi /2)}}.`$
In calculating (5.4) we have made use of the expansions of $`E`$ and $`F`$ to $`O(\beta ^2)`$:
$$Ee_0+e_1\frac{\beta ^2}{4\pi }Ff_0+f_1\frac{\beta ^2}{4\pi },$$
(5.5)
with
$$e_0=a_0+a_1\text{and}f_0=a_0a_1.$$
(5.6)
Since $`K_q=K_c+\delta K_c`$, we deduce that
$$\frac{\delta K_c}{K_c}=\frac{i\beta ^2}{8}\mathrm{sinh}\theta (\theta ).$$
(5.7)
Thence, expanding to $`O(ϵ)`$, we find,
$`(\theta )`$ $`=`$ $`\{{\displaystyle \frac{1}{\mathrm{cosh}\theta +1}}{\displaystyle \frac{1}{\mathrm{cosh}\theta }}+{\displaystyle \frac{e_1}{\mathrm{cosh}\theta +\mathrm{sin}a_1\pi }}{\displaystyle \frac{e_1}{\mathrm{cosh}\theta \mathrm{sin}a_1\pi }}`$ (5.8)
$`+{\displaystyle \frac{e_1ϵ\mathrm{cos}a_1\pi }{2\mathrm{sin}a_1\pi }}\left({\displaystyle \frac{1}{(\mathrm{cosh}\theta +\mathrm{sin}a_1\pi )^2}}+{\displaystyle \frac{1}{(\mathrm{cosh}\theta \mathrm{sin}a_1\pi )^2}}\right)`$
$`+{\displaystyle \frac{ϵf_1}{\mathrm{sin}a_1\pi \mathrm{cosh}^2\theta }}\}.`$
Comparing (5.2) with (5.8) we see a pleasing similarity. In fact the two formulae are identical, to $`O(ϵ)`$, provided we choose $`e_1`$ and $`f_1`$ suitably. In other words, we may deduce that
$$e_1=2a_1+\frac{ϵ}{\pi \mathrm{sin}a_1\pi }(a_0+a_1)+O(ϵ^2),$$
(5.9)
and $`f_1`$ is proportional to $`ϵ`$. Unfortunately, the calculation does not allow anything more detailed to be learned about $`f_1`$. To do better would need a correction to the reflection factor to $`O(ϵ^2)`$.
## 6 Discussion
The purpose of this calculation was to test a little more deeply the expression for the sinh-Gordon particle reflection factor given in and to learn additional information concerning its dependence on the boundary parameters $`\sigma _0`$ and $`\sigma _1`$. The result of the investigation is gratifying because it agrees with alternative derivations of the reflection factor and it also agrees with the following conjecture. Everything we have learned so far is consistent with quite simple expressions for $`E`$ and $`F`$:
$$E=(a_0+a_1)(1B/2)F=(a_0a_1)(1B/2),$$
(6.1)
where the coupling constant dependence enters via the expression for $`B`$ given in (2.4). Similar expressions for these parameters have been arrived at via other arguments by Zamolodchikov .
If (6.1) is correct then the reflection factor is invariant under the interchange $`a_0a_1`$. In effect, this invariance restores the $`𝐙_\mathrm{𝟐}`$ bulk symmetry which apparently was broken by the boundary condition and replaced by a symmetry under the simultaneous interchange of $`\varphi `$ with $`\varphi `$ and $`a_0`$ with $`a_1`$. The reflection factor is also invariant if $`a_0`$ and/or $`a_1`$ is replaced by its negative, as it should be given the definitions of $`\sigma _0`$ and $`\sigma _1`$. It is consistent with what is known at the special value of the coupling constant, known as the ‘free-fermion’ point in the sine-Gordon model, where $`B=2`$ and the S-matrix is unity. There, the restrictions on the parameters in the reflection factor can be solved exactly and are in agreement with (6.1) .
Note also, that with the expressions (6.1) the reflection factor (3.2) has a weak-strong coupling symmetry matching the symmetry of the S-matrix under $`\beta 4\pi /\beta `$. To see this, note that setting
$$(a_0^{},a_1^{},\beta ^{})=\frac{4\pi }{\beta ^2}(a_0,a_1,\beta )$$
(6.2)
defines a new triple of coupling constants with the property that
$$K_q(\theta ,a_0,a_1,\beta )=K_q(\theta ,a_0^{},a_1^{},\beta ^{}).$$
(6.3)
If (6.1) is correct, implying the duality symmetry (6.2), then we are faced with other puzzles. For example, it is known that the supersymmetric version of the sinh-Gordon model is only integrable when restricted to a half-line with some very special boundary conditions (either $`a_0=a_1=0`$ or, $`a_0=a_1=1`$) (see ), and this restriction would appear to be incompatible with a weak-strong coupling symmetry without modifying (6.1).
It is also known that the other affine Toda field theories constructed from data in the $`ade`$ series, when restricted to a half-line, allow only a finite number of possible boundary conditions. In fact, the $`a_1^{(1)}`$ or sinh-Gordon model is apparently the only example within this series which allows continuous boundary parameters. Expressions for the associated reflection factors for the other models are largely unknown but, it will be interesting to discover if they too can permit a duality symmetry in the presence of a boundary which will match the symmetry of their bulk S-matrices.
## 7 Acknowledgement
One of us (AC) wishes to thank the Ministry of Culture and Higher Education of Iran for financial support. The other (EC) thanks G.W. Delius and Al.B. Zamolodchikov for discussions and the European Commission for partial support under a TMR grant, number FMRX-CT96-0012. |
warning/0002/hep-ph0002192.html | ar5iv | text | # NUC-MINN-00/05-TFebruary 2000 Modification of Z Boson Properties in Quark-Gluon Plasma
## Acknowledgments
We thank P. Stankus and R. Rusack for prodding us to investigate this problem. This work was supported by the US Department of Energy under grant DE-FG02-87ER40328. |
warning/0002/cond-mat0002182.html | ar5iv | text | # Anomalous metallicity and electronic phase separation in the CsC60 polymerized fulleride
## I INTRODUCTION
When the face centered cubic (f.c.c) phase of $`\mathrm{A}_1\mathrm{C}_{60}`$(A=K, Rb, Cs) compounds is slowly cooled from 400K, one-dimensional polymerization of C<sub>60</sub> molecules spontaneously occurs along the (110) cubic direction and leads to an orthorhombic phase. A drastic change in the electronic properties is observed at the structural transition. Indeed, the f.c.c. phase was shown to be a Mott insulator, whereas in the orthorhombic phase a plasma frequency was measured in optical experiments . However, the density of carriers is likely to be rather low or their effective mass very large since the plasma frequency is equal to 0.1eV in $`\mathrm{KC}_{60}`$ polymer and even lower for $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$. In addition, the low frequency conductivity of both $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$ decreases smoothly over a broad temperature range, being at variance with $`\mathrm{KC}_{60}`$ which remains conducting down to 4.2K. Furthermore, the temperature dependence of the $`{}_{}{}^{13}\mathrm{C}`$ spin-lattice relaxation rate shows that strong magnetic fluctuations are present up to room temperature in $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$ and the sharp decrease of the uniform static susceptibility (measured from EPR line intensity) below 50K for $`\mathrm{RbC}_{60}`$ and 40K for $`\mathrm{CsC}_{60}`$, suggests that both compounds undergo magnetic transitions at these respective temperatures. The occurrence of spin ordering is also evident from NMR experiments: the slowing down of magnetic fluctuations gives rise to a divergent relaxation rate below 40K. However, the nature of the spin order is less obvious. On one hand, EPR experiments suggest the onset of a spin density wave ground state as a result of a possible one-dimensional (1D) character of the band structure. On the other hand, $`\mu \mathrm{SR}`$ studies show a gradual transition towards a highly disordered magnetic phase and do not rule out the possibility of a random spin freezing below 40K.
In a recent NMR work, we have shown that some of the $`{}_{}{}^{133}\mathrm{Cs}`$ sites remain unaffected by the onset of the spin-ordering in the low temperature state, magnetic and nonmagnetic domains being spatially distributed. At the temperature of 13.8K the occurrence of a charge redistribution and a concomitant decrease of the local electronic susceptibility inside these nonmagnetic domains have been observed . In agreement with this latter result, detailed analysis of the EPR linewidth at ambient pressure also suggest that two distinct magnetic environments coexist in the low temperature state of $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$ polymers and insofar as a charge redistribution occurs in the nonmagnetic domains at 13.8K, the spontaneous thermal contraction recently observed at 14K by X-Ray diffraction in $`\mathrm{CsC}_{60}`$ strongly supports the fact that these inhomogeneities are intrinsic.
In this manuscript, we give experimental evidence showing that the “conducting” state of the $`\mathrm{CsC}_{60}`$ polymerized phase cannot be understood within the framework of an electronic band conductor as claimed earlier . We first report the temperature dependence of the spin lattice relaxation rate $`(T_1)^1`$ for both $`{}_{}{}^{13}\mathrm{C}`$ and $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei at different pressures up to 9 kbar, indicating that in the temperature domain above 80K two different mechanisms govern the relaxation of $`{}_{}{}^{13}\mathrm{C}`$ and $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei respectively. As far as $`{}_{}{}^{133}\mathrm{Cs}`$ is concerned, $`{}_{}{}^{133}(T_1)_{}^{1}`$ decreases linearly down to about 80K though remaining pressure independent up to 9kbar. This behavior is in sharp contrast with the $`{}_{}{}^{13}\mathrm{C}`$ nuclei for which $`{}_{}{}^{13}(T_1)_{}^{1}`$ strongly decreases under pressure up to 9kbar while remaining almost temperature independent. The difference between $`{}_{}{}^{133}\mathrm{Cs}`$ and $`{}_{}{}^{13}\mathrm{C}`$ nuclei exists independently of the nature of the ground state of the system. More insight into these peculiar properties is then obtained using quadrupolar echo experiments performed on the $`{}_{}{}^{133}\mathrm{Cs}`$ nucleus which enable us to analyze with great accuracy the temperature dependence of the NMR spectrum at 1bar. We show that the NMR spectrum of the two phases (magnetic and nonmagnetic) is motional narrowed above 100K because of the fast motion of the local environment around the $`{}_{}{}^{133}\mathrm{Cs}`$ sites. The evolution of the lineshape with temperature reveals that the static coexistence of two different $`{}_{}{}^{133}\mathrm{Cs}`$ sites below 15K arises from a gradual freezing of these fluctuations in the local environment.
## II EXPERIMENTAL DETAILS
The measurements have been conducted on two powdered samples with entirely consistent results, one of them (10%) $`{}_{}{}^{13}\mathrm{C}`$ enriched. The pressure set up is a homemade double-stage copper-beryllium cell using fluor-inert as the pressure medium. This enables us to correct for each temperature the loss of pressure within the sample chamber due to the gradual freezing of the fluor-inert.
The spin-lattice relaxation were measured by monitoring the recovery of the magnetization after saturation with a series of $`\pi /2`$ pulses. The recovery curve is exponential for $`{}_{}{}^{133}\mathrm{Cs}`$ and $`{}_{}{}^{13}\mathrm{C}`$ at room temperature. At ambient pressure, the recovery curve gradually becomes bi-exponential for both nuclei below 40K. A large distribution of short relaxation rates is observed giving raise to a recovery curve of the following shape $`1e^{(t/T_1)^\beta }`$ with a value of $`\beta `$ of the order of 0.5 at the lowest temperature investigated i.e 4K. At 5kbar the recovery curve is for $`{}_{}{}^{133}\mathrm{Cs}`$ exponential down to 4K. Not so for $`{}_{}{}^{13}\mathrm{C}`$ since a nonexponential recovery is observed below 20K. Different fit procedures did not help us to determine without ambiguity the shape of the recovery but no significant change were observed on the qualitative temperature dependence of $`{}_{}{}^{13}(T_1)_{}^{1}`$. The relaxation rates $`{}_{}{}^{13}(T_1)_{}^{1}`$ shown on Fig.1b at 5kbar are therefore deduced below 20K from a fit of the recovery curve assuming it to be exponential as above 20K. At 9kbar, the recovery curve is exponential for $`{}_{}{}^{133}\mathrm{Cs}`$ and $`{}_{}{}^{13}\mathrm{C}`$ in the all temperature range investigated.
Finally we should point out that despite the presence of a static quadrupole splitting for the NMR line of $`{}_{}{}^{133}\mathrm{Cs}`$ in the orthorhombic phase, the smallness of the quadrupole frequency which is of the order of 5kHz enables us to saturate all the transitions at once. Therefore the nuclear levels are initially equally populated establishing a well-defined spin temperature (equal to infinite). In that case no deviation from an exponential behavior is expected for the relaxation of the magnetization which perfect exponential recovery at room temperature is a proof of the homogeneity of the samples.
## III $`{}_{}{}^{13}\mathrm{C}`$ and $`{}_{}{}^{133}\mathrm{Cs}`$-NMR UNDER PRESSURE
We report on Fig.1a and Fig.1b the temperature dependence of the relaxation rate for $`{}_{}{}^{133}\mathrm{Cs}`$ and $`{}_{}{}^{13}\mathrm{C}`$ nuclei at ambient pressure, 5kbar and 9kbar. The large enhancement of $`{}_{}{}^{133}(T_1)_{}^{1}`$ and $`{}_{}{}^{13}(T_1)_{}^{1}`$ below 40K is due to a slowing down of magnetic fluctuations which is completely suppressed at 5kbar. At this pressure, both $`{}_{}{}^{133}(T_1)_{}^{1}`$ and $`{}_{}{}^{13}(T_1)_{}^{1}`$ decrease exponentially below 20K revealing the opening of a spin-gap at $`T_C20K`$, the ground state being homogeneous and nonmagnetic. The effect of an applied pressure on this long range order has been carefully investigated by $`{}_{}{}^{133}\mathrm{Cs}`$-NMR. The temperature dependence of $`{}_{}{}^{133}(TT_1)_{}^{1}`$ at 5, 5.5, 5.7 and 9kbar is shown in Fig.2.The well-defined instability at 5kbar gives rise to a sharp peak on $`{}_{}{}^{133}(TT_1)_{}^{1}`$ at $`T_C`$. Quite remarkably, a slight increase of the applied pressure strongly reduces the amplitude of the spin-gap, without any significant change in $`T_C`$ itself (as given by the position of the $`{}_{}{}^{133}(TT_1)_{}^{1}`$ peak, see on Fig.2). However, a smooth decrease of the temperature $`\mathrm{T}_{\mathrm{Mag}}`$ at which the slowing down of magnetic fluctuations occurs has been observed by EPR experiments under pressure up to 4kbar. This fact is also evident from the temperature dependence of the linewidth of the $`{}_{}{}^{133}\mathrm{Cs}`$ NMR line shown at different pressures on Fig.3. Thus, as the pressure increases $`\mathrm{T}_{\mathrm{Mag}}`$ drops continuously along a transition line which does not exist for the case of the spin-singlet ground state. Henceforth we can infer that the sharp suppression of the spin-gap below 20K which in turn gives rise to a metallic state, is not due to continuous changes in the magnitude of the electronic interactions but may reflect some structural changes above 5kbar as suggested by DC conductivity measurements performed under pressure.
Another striking feature in the response of the polymerized phase $`\mathrm{CsC}_{60}`$ to high pressure appears at glance in Fig.1a and Fig.1b. Indeed, a clear distinction has to be made between the two temperature domains 4.2-80K and 80K-300K. Below 80K, $`{}_{}{}^{133}(T_1)_{}^{1}`$ and $`{}_{}{}^{13}(T_1)_{}^{1}`$ exhibit a similar pressure and temperature dependence. This is, however, no longer true above 80K, where $`{}_{}{}^{133}(T_1)_{}^{1}`$ and $`{}_{}{}^{13}(T_1)_{}^{1}`$ behave in complete different ways. In particular, we can see in Fig.1a that above 80K, $`{}_{}{}^{133}(T_1)_{}^{1}`$ shows no pressure dependence up to 9kbar unlike $`{}_{}{}^{13}(T_1)_{}^{1}`$, which is shown on Fig.1b. Within the first five kilobars, the relaxation of $`{}_{}{}^{13}\mathrm{C}`$ nuclei is strongly affected by pressure in two manners:(i) an overall depression is observed under pressure following the depression of the uniform spin susceptibility ($`\chi `$) measured by EPR which drops at a rate of about 10% per kbar, (ii) a weakly temperature dependent contribution to $`{}_{}{}^{13}(T_1)_{}^{1}`$ (20% decrease from 300 to 40K) is suppressed at 5 kbar.
Broadly speaking, the spin-lattice relaxation rate for a given nuclei and the static electronic spin susceptibility are linked together by the following relation : $`(T_1T)^1_q|A(\stackrel{}{q})|^2\chi {}_{}{}^{}{}_{}{}^{^{\prime \prime }}(\stackrel{}{q})`$ where $`A(\stackrel{}{q})=_iA_ie^{i\stackrel{}{q}.\stackrel{}{r}_i}`$ is the form factor of the hyperfine interaction between a given nuclei and the electronic spins located at the neighboring sites. Unlike $`{}_{}{}^{133}\mathrm{Cs}`$ which environment is octahedral, there is no particular symmetry for $`{}_{}{}^{13}\mathrm{C}`$ sites. If both nuclei are coupled to the same electronic spins then, that $`{}_{}{}^{13}(T_1)_{}^{1}`$ and $`{}_{}{}^{133}(T_1)_{}^{1}`$ display a different pressure and temperature dependence above 80K, might be attributed to the presence of a spatially dependent electronic spin susceptibility which dominates the relaxation of $`{}_{}{}^{13}\mathrm{C}`$ nuclei. However, as previously shown for $`\mathrm{RbC}_{60}`$ , the decrease of $`{}_{}{}^{13}(T_1)_{}^{1}`$ follows the decrease of the uniform spin susceptibility deduced from EPR within at least the first five kilobars. This reveals that in the low pressure regime, the dominant contribution to the relaxation of $`{}_{}{}^{13}\mathrm{C}`$ above 80K is due to enhanced magnetic fluctuations at the wave vector $`\stackrel{}{q}=0`$ and therefore, the differences described above between $`{}_{}{}^{13}\mathrm{C}`$ and $`{}_{}{}^{133}\mathrm{Cs}`$ cannot be ascribed to the form factor of the alkali site in the polymerized phase.
As it is, one can draw the following conclusions. First, the absence of pressure dependence observed for $`{}_{}{}^{133}(T_1)_{}^{1}`$ above 80K shows that the dominant contribution to the fluctuating field at $`{}_{}{}^{133}\mathrm{Cs}`$ site in this temperature range is unrelated to the electronic spins involved in the relaxation of $`{}_{}{}^{13}\mathrm{C}`$ nuclei. Secondly, the fact that above 80K, $`{}_{}{}^{13}(T_1)_{}^{1}`$ is weakly temperature dependent at ambient pressure and constant at 5kbar suggests that the electronic spins are localized. This latter conclusion is in agreement with the calculated band structure of the polymer $`(\mathrm{C}_{60}^{})^n`$ which displays a dispersionless 1D half-filled band at the Fermi level but in apparent contradiction with transport measurements performed in the similar compound $`\mathrm{RbC}_{60}`$.
One can therefore conclude that a model based on a single electron specie is inadequate for describing the electronic properties of the polymerized phases $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$.
In a previous work, we have shown that the use of quadrupolar spin echoes of $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei enables to reveal the presence of nonmagnetic domains within a magnetic background. However, whether this inhomogeneous state results from the existence of static structural defects along the chains or is purely electronically driven e.g. as proposed for underdoped cuprates and spin-ladders compounds, remained an open question. In what follows, we address this problem again with the aid of quadrupolar spin echoes in order to determine how the inhomogeneous state at low temperature arises from the high temperature one.
## IV $`{}_{}{}^{133}\mathrm{Cs}`$-NMR AT AMBIENT PRESSURE
In a similar way than in the reference, the spin echoes of $`{}_{}{}^{133}\mathrm{Cs}`$ have been obtained after a ($`\pi /2\tau \pi /8`$) in-phase RF pulse sequence, maintaining fixed the echo delay $`\tau `$ at 40$`\mu s`$. Half of the spin-echo at 3$`\tau `$ is then Fourier transformed. This procedure gives rise to a spectrum containing two lines $`5/23/2`$ and $`3/25/2`$ split by an amount $`4\nu _Q`$, where $`\nu _Q`$ is the quadrupole frequency of $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei in the polymerized phase . The evolution of the $`{}_{}{}^{133}\mathrm{Cs}`$ spectrum is displayed on Fig.4 at different temperatures between 100 and 4.2K. The expected doublet spectrum corresponding to a single $`{}_{}{}^{133}\mathrm{Cs}`$ site is observed at 100K, but as $`T`$ approaches 40K, the shape becomes asymmetric and a fine structure gradually develops. At 25K, the coexistence of two different $`{}_{}{}^{133}\mathrm{Cs}`$ sites is evident in Fig.4, with a frequency difference in the local field of the order of $`4\nu _Q`$. This means therefore that two distinct magnetic environments are spatially distributed at this temperature. As the temperature is further lowered, the situation with a single quadrupolar split is recovered and thus only one $`{}_{}{}^{133}\mathrm{Cs}`$ site contributes to the spin echo signal below 15K. The amplitude of the spin echo refocused at 3$`\tau `$ is proportional to $`e^{3\tau \gamma \mathrm{\Delta }H(T)}`$ where $`\mathrm{\Delta }H(T)`$ is the width of the local field distribution due to the static electronic moments at a given temperature $`T`$. Considering two distinct populations of $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei below 30K, $`N_m`$ and $`N_{nm}`$ which are coupled to the local field inhomogeneity $`\mathrm{\Delta }H(T)`$ and located inside the nonmagnetic domains respectively, the total number of $`{}_{}{}^{133}\mathrm{Cs}`$ sites contributing to the spin echo signal at 3$`\tau `$ can be expressed as: $`N(T)=N_m/(1+(3\tau \gamma \mathrm{\Delta }H(T))^2)+N_{nm}`$. If the condition $`3\tau \gamma \mathrm{\Delta }H(T)1`$ is fulfilled, only a fraction $`N_{nm}`$ of the nuclei contribute to a spin echo at $`3\tau `$ since this experiment selects those $`\mathrm{Cs}`$ sites which are entirely decoupled from the onset of local magnetism. Let $`I(T)`$ be the integrated intensity of the Fourier transform performed on this spin echo. The temperature dependence of $`N(T)`$ (equal to $`I(T).T`$) is reported on Fig.5. We observe that a majority of the $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei is gradually wiped out of the signal below 40K. A minimal value for $`N_{nm}`$ is reached at 15K and amounts to about 10% of the total number of nuclei at $`40K`$. However, the estimated ratio between the two phases from the $`{}_{}{}^{13}\mathrm{C}`$ spectrum suggests that approximately half of the $`{}_{}{}^{13}\mathrm{C}`$ sites do not see the magnetic moment distribution in the low temperature state. We may solve this puzzle by considering that the $`{}_{}{}^{13}\mathrm{C}`$ spins probe the very local properties within each $`\mathrm{C}_{60}`$ chains carrying the electronic spins whereas only $`{}_{}{}^{133}\mathrm{Cs}`$ sites far from any magnetic domain will contribute to the echo signal refocused at $`3\tau `$. This would mean that the boundary surface is large compare to the domains size suggesting that the phase separation sets on a microscopic scale.
To gain insight into the driving force of this process more attention must be paid to what happens above the spin ordering temperature. In particular, we see on Fig.4 that the splitting of the $`{}_{}{}^{133}\mathrm{Cs}`$ spectrum displays a fine structure near 40K although the NMR spectrum corresponding to two $`{}_{}{}^{133}\mathrm{Cs}`$ sites is not yet resolved. This can be understood if we assume that the local field of a $`{}_{}{}^{133}\mathrm{Cs}`$ nucleus jumps randomly from one value to the other in the “conducting” state. Indeed, using only the difference between the resonance frequencies $`\delta \omega `$ and the hopping time $`\tau _h`$, we can propose the following scenario. At high temperature, $`\delta \omega \tau _h1`$ and the spectrum is motional narrowed, which means that only one doublet is visible. When the temperature is lowered, the jump frequency ($`1/\tau _h`$) decreases and the condition $`\delta \omega \tau _h1`$ becomes fulfilled with a fine structure developing in the NMR spectrum. Finally, when $`\delta \omega \tau _h1`$, the quadrupolar splitting of the two sites are well resolved, i.e. one for $`{}_{}{}^{133}\mathrm{Cs}`$ sites in the magnetic domains and the other for $`{}_{}{}^{133}\mathrm{Cs}`$ sites in the nonmagnetic ones. We simulate each of the three cases and our simulations at fixed $`\delta \omega `$ are shown in Fig.6 for different correlation times $`\tau _h`$ and superimposed (dotted line) on the experimental spectra on Fig.4. Clearly, the calculated spectra bear a strong resemblance with the experimental ones displayed on Fig.3 between 100 and 30K. We can thus infer the existence of a thermally activated change in the local environment of $`{}_{}{}^{133}\mathrm{Cs}`$ sites which may become the dominant contribution to $`{}_{}{}^{133}(T_1)_{}^{1}`$ when the frequency $`1/\tau _h`$ is of the order of the Larmor frequency (43 MHz) of the $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei. Therefore, from the results exposed in this section we can conclude to the existence of another degree of freedom aside from the fluctuations of the electronic spins located on the $`\mathrm{C}_{60}`$ molecules, and possibly related to spontaneous local structural changes in the polymerized phase.
## V DISCUSSION
As emphasized above, one of the difficulty aroused by our work is to bring together the conducting nature of the polymerized phase established by optical and transport measurements with the pressure and temperature dependence of $`{}_{}{}^{13}(T_1)_{}^{1}`$ which strongly suggest that electrons are localized. It therefore turns out natural to question ourselves about the possible relationship between the local structural change around $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei and the presence of charge degrees of freedom like polarons in the polymerized phase. On the basis of the above NMR results and anticipating results described further on, we suppose that the mobility of a charge carrier in the polymerized phase mainly depends upon the occurrence of a local structural distortion in its vicinity. From a point of view which is somewhat naive, one may consider that at thermal equilibrium, the charge carriers diffuse through the lattice under the action of a random force $`F(t)`$ which takes on only two discrete values $`\pm f_0`$. For our particular purpose, the relevant physical quantity to be consider is the spectral density $`F(\omega )`$ defined as the Fourier transform of the correlation function $`F(t)F(t+\tau )`$, the brackets indicating an ensemble average. In our case, $`F(t)F(t+\tau )`$ can be assumed to be of the form: $`f_0^2e^{|t|/\tau _h}`$, which leads to the following spectral density : $`F(\omega )=\tau _h/(1+(\omega \tau _h)^2)`$. Because any excited state of the charge carriers is to relax due to the random force $`F(t)`$, the spectral density $`F(\omega )`$ will lead to a strong frequency dependence in the response function of the carriers to external oscillating fields. It is therefore of a great interest to focus on AC resistivity measurements performed at ambient pressure in both $`\mathrm{KC}_{60}`$ and $`\mathrm{RbC}_{60}`$. For $`\mathrm{KC}_{60}`$ which does not exhibit a slowing down of spin fluctuations, AC and DC resistivities display a similar temperature dependence. This is however not true for $`\mathrm{RbC}_{60}`$ since a frequency dependent peak is clearly observed on AC resistivity. The peak shifts from 35K at 1.1kHz down to 25K at 43Hz, the order of magnitude of these frequencies being in good agreement with the value we deduced from our simulate spectra in the same temperature range for $`\mathrm{CsC}_{60}`$ (c.f.Fig.6). The fact that the electronic properties of $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$ display similar electronic and structural features as opposed to $`\mathrm{KC}_{60}`$ allows us to extrapolate the results obtained by Zhou et al for $`\mathrm{RbC}_{60}`$ to the case of $`\mathrm{CsC}_{60}`$. Thus experiments show that in the two polymerized phases $`\mathrm{RbC}_{60}`$ and $`\mathrm{CsC}_{60}`$, the dissipation reaches a maximum when the hopping frequency of the local environment of the alkali ion becomes equal to the AC frequency. Such a coincidence can be hardly fortuitous and suggests that the mobility of the charge carriers in the polymerized phase is strongly coupled to the environment of the alkali ion. In this context, it is worthwhile to mention that polaron-like distortions such as $`\mathrm{C}_{60}^{1x}\mathrm{C}_{60}^{1+x}`$ have been predicted to be energetically favorable in the charged polymer $`(\mathrm{C}_{60}^{})^n`$ which exhibits a tendency to undergo a charge density wave transition . In that particular case, the conduction mechanism would be due to an intramolecular property of the polymer itself and that would drastically change our expectations regarding the pressure effect on the electronic properties of the polymerized phase. However, on the sole basis of the NMR experiment above described we cannot address the microscopic mechanism at the origin of the spontaneous formation of polarons in the polymerized phase.
In the light of the above considerations, it is interesting to shortly reconsider the pressure effect on the spin-lattice relaxation rate $`{}_{}{}^{13}(T_1)_{}^{1}`$ of $`{}_{}{}^{13}\mathrm{C}`$ nuclei in the low pressure regime. As mentioned above, $`{}_{}{}^{13}(T_1)_{}^{1}`$ shows at room temperature a similar pressure decrease than the electronic spin susceptibility deduced from EPR which suggests that magnetic fluctuations at the wave vector $`\stackrel{}{q}=0`$ dominate $`{}_{}{}^{13}(T_1)_{}^{1}`$ at ambient pressure. One possible explanation for the origin of these enhanced uniform fluctuations might be that polarons acting as local defects, induce disorder in the AF exchange coupling $`J`$ along the chain leading to the formation of spin clusters. It was indeed shown theoretically that the low energy magnetic fluctuations (i.e when $`TJ`$) of a disordered AF spins chain are merely governed by clusters with an odd number of spins, each one acting as a nearly free localized (1/2) spin. In such a case the reversible suppression at 5kbar of a weak temperature dependent term in $`{}_{}{}^{13}(T_1)_{}^{1}`$ could be ascribed to the suppression with applied pressure of disorder in the magnetic coupling along the chain which presence would be henceforth closely related to the slowing down of spin fluctuations in the low temperature state. Much more experimental inputs are however required to go beyond this statement.
As it is, the phase separation occurring in the low temperature state at ambient pressure appears to be the logical outcome of the twofold nature of the polymerized phase $`\mathrm{CsC}_{60}`$ that is : mobile polarons spontaneously form aside from localized electrons and compete with a 3D magnetic order imposed by the transverse dipolar coupling between the chains. Note that the presence of nonmagnetic domains is in itself a strong hint that polarons are not randomly spatially distributed within the magnetic background but may form collective structures developing a long range order below 14K as suggested by NMR and X-ray experiments.
## VI Conclusion
The work described in this manuscript deals with the electronic properties of the polymerized phase $`\mathrm{CsC}_{60}`$ extensively studied by NMR of $`{}_{}{}^{13}\mathrm{C}`$ and $`{}_{}{}^{133}\mathrm{Cs}`$ nuclei. The salient result is that the electronic properties of the polymerized phase $`\mathrm{CsC}_{60}`$ involve two degrees of freedom : one related to localized spins, the other related to mobile charges which mobility is strongly entangled to the local environment of the Cs ion. The polymerized phase $`\mathrm{CsC}_{60}`$ is therefore dynamically inhomogeneous and as shown by NMR under pressure, this feature persists up to 9kbar.
At ambient pressure static inhomogeneities gradually develop below 40K concomitantly with a slowing down of spin fluctuations. At 5kbar, the polymerized phase $`\mathrm{CsC}_{60}`$ undergoes a nonmagnetic transition at $`T_c`$ equal to 20K. The ground state is homogeneous and a spin gap opened below 20K. Finally, a dramatic decrease of the amplitude of the spin gap is observed above 5kbar without any significant decrease of $`T_c`$. The presence of magnetism therefore appears to be closely related to the occurrence of static inhomogeneities. How does the applied pressure suppress these inhomogeneities and stabilize a homogeneous nonmagnetic ground state? That cannot be addressed by the present work but remains an important issue to be solved.
## VII ACKNOWLEDGMENT
It is a pleasure to thank F. Rachdi for the $`{}_{}{}^{13}\mathrm{C}`$ enriched $`\mathrm{C}_{60}`$. We are also very grateful to C. Berthier, S. Brasovski, P. Carretta, P. Sotta and P. Wzietek for illuminating discussions and to J. P. Cromières and M. Nardone for technical assistance.
One of the authors (L.F) is grateful for the support of the Swiss National Science Foundation. |
warning/0002/hep-th0002148.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Random matrix models enjoy a wide range of applications in physics due to their property of being universal (for reviews see ). This property manifests itself in the independence of correlation functions from the choice of the distribution function $`𝒫(M)\mathrm{exp}[n\text{Tr}V(M)]`$, where $`M`$ is an $`n\times n`$ matrix and $`V`$ is a polynomial. Different classes of universality are found depending on the way the large-$`n`$ limit is taken and which part of the spectrum is investigated . However, not in all applications a distribution function of the above form is realistic. Here, all the eigenvalues (=energy levels) of the matrix $`M`$ are coupled through the Jacobian after diagonalization. There are situations as for example in applications in Nuclear Physics, where in contrast to that the Hamiltonian of the model couples only few energy levels and is still very well described by the above random matrix model . It is therefore very interesting to study deformations and generalizations of the above distribution $`𝒫(M)`$ and to investigate first, if the correlations remain unchanged and second, if the property of universality is maintained. In this work we study a deformation which preserves the symmetry of the matrix model, which will be the unitary transformations of the Hermitian matrix $`M`$ in our case. The symmetry of the model is directly related to the properties of the Hamiltonian under rotations and time-reversal . There exists an interesting relation between the restricted trace ensembles (RTEs) which we will consider and the so-called Wigner ensembles , where different matrix elements are weighted with different distribution functions, without being invariant under unitary transformations.
It has been known only quite recently that there exist symmetry preserving deformations of the distribution function that destroy the property of universality . Two examples are the trace squared ensembles which were originally introduced in the context of Quantum Gravity and the generalized RTEs which were introduced, in their simplest pure-quadratic form, by Bronk and Rosenzweig in the context of Nuclear Physics. The non-universality does not necessarily spoil the applicability to physical systems. In fact the deformation may be introduced for physical reasons: indeed, trace-squared terms have been recognized as corresponding to higher order intrinsic curvature terms in the string action. Such terms have been added in order to cure an intrinsic instability of the theory related to a crumpled surface phase of the string world sheet . Moreover, in the framework of random surfaces, these additional terms are interpreted as touching interaction terms which make the random surfaces touch each other. The generalized RTEs permit the same graphical interpretation as they have been shown to be a limiting case of a special trace squared ensemble . The RTEs are defined such that the exponential weight function gets replaced by a constraint $`𝒫_\delta (M)\delta (A^21/n\text{Tr}V(M))`$ (or the $`\theta `$ Heaviside step function instead). By using the following representation of the $`\delta `$-function, $`\delta (x)=lim_l\mathrm{}\sqrt{\pi /l}\mathrm{exp}[lx^2]`$, the distribution $`𝒫_\delta (M)`$ can be easily seen to contain trace squared terms. The relation between RTEs and the trace squared ensembles was discussed in great detail for the spectral density in using saddle point techniques, and the canonical ensemble, the trace squared ensemble and the RTE were shown to agree to leading order. In also the two-point function was derived for the trace squared ensembles and shown to be non-universal using saddle point equations. Here, we will present a scheme to calculate general $`k`$-point functions of the RTEs in different large-$`n`$ regimes.
Namely, one has to distinguish two different types of large-$`n`$ limits - the macroscopic and microscopic limit - which may not all be affected by the deformation of the distribution function $`𝒫(M)`$. In the example of generalized RTEs which we will study here we find that this is precisely the case. While the macroscopic universality is destroyed by the global constraint, the microscopic correlations at short distances remain unchanged. In a sense, the canonical ensemble $`𝒫(M)`$ is replaced by its micro-canonical counterpart $`𝒫_\delta (M)`$. Therefore we have an explicit model of statistical mechanics at hand, where the correlation functions of both ensembles can be calculated analytically and then be compared for discrepancies. Another remarkable property of RTEs is that they possess a finite support already at finite $`n`$. Similar to the canonical Gaussian matrix model the RTEs allow for an explicit calculation at finite $`n`$. This has been shown already in for the one-and two-point function and will be given here for general $`k`$-point functions. Comparing the finite-$`n`$ and $`n\mathrm{}`$ results the differences can be thought of as finite-size corrections, when interpreting $`n\mathrm{}`$ as the continuum limit. In the RTE these corrections appear in a different way than in the canonical model, due to the finite support at finite-$`n`$. Let us explain now in more detail how the large-$`n`$ limit can be taken.
In the macroscopic large-$`n`$ limit no restrictions are made on the distance between different eigenvalues. This leads to smooth, universal two- and higher $`k`$-point correlation functions for the canonical, unconstrained models . Applications can be found in two-dimensional Quantum Gravity as well as the theory of transport properties of mesoscopic wires . For the generalized RTEs we have shown in a previous publication together with our collaborators , that for a certain class of potentials $`V(M)`$ the 2-point correlator is no longer universal. In this work we will extend these results to arbitrary polynomial potentials (see also ) and to all higher $`k`$-point resolvents. Since in ref. it was also shown that to leading order all $`k`$-point resolvents of the $`\delta `$\- and $`\theta `$-measure are equivalent we will restrict ourselves here to the former one.
In contrast to that, in the microscopic large-$`n`$ limit correlations of eigenvalues $`\lambda ,\mu `$ at a distance of the mean level spacing $`|\lambda \mu |1/n`$ are calculated. It is this kind of limit that finds a wide range of applications in nuclear physics, condensed matter physics (see e.g. ) and has initiated exact analytical solutions in the study of Dirac spectra in QCD . We will show that in generalized RTE with purely monomial potential $`V(M)=M^{2p}`$ (which includes of course the quadratic case for $`p=1`$) the constraint does not change the local properties at short distances, in the sense that the connected two-point correlator behaves according to the well-known “sine-law” of the Gaussian canonical ensemble . Indeed, this result holds also for higher $`k`$-point correlation functions. This does not come as a surprise since a global constraint should not change the local properties. We believe that the same is true for more general potentials and that universality holds here as well.
Consequently the present paper is split into two different parts. Section 2 is devoted to the macroscopic large-$`n`$ limit where we use loop equation techniques, closely following . Although the loop equations are originally designed to calculate higher orders in the $`1/n`$-expansion we will restrict ourselves to the planar limit. Since we find non-universality to leading order in all connected correlation functions we do not calculate the likewise non-universal higher orders in $`1/n`$. The difficulty to deal with the constraint will be treated in a similar way as the situation where the spectral density has a support consisting of several intervals . In these multi-band phases additional constraints have to be imposed to make the solution unique . We will restrict ourselves to hermitian matrices $`M`$ only, for non-hermitian matrices see . Since only the planar solution is needed for our results they can be easily extended to orthogonal and symplectic matrices using . Extensions to the complex matrix model are straightforward as well since the same loop equation techniques exist in the literature . Very recently finite-$`n`$ results have been obtained for the eigenvalue density of Gaussian ensembles with real symmetric and complex matrices . In the second part section 3 the microscopic large-$`n`$ limit is investigated for the RTE with a monomial potential $`V(M)=M^{2p}`$. Here, we generalize existing results for finite-$`n`$ of previous publication , by improving a technique already used in ref. . Rescaling variables in the microscopic limit and using the inverse Laplace transform we are able to match the connected $`k`$-point correlator to the well known “sine-law” behavior of the canonical ensemble, thereby proving the microscopic RTE universality. Let us stress that for the RTEs no orthogonal polynomial techniques are applicable.
## 2 The macroscopic limit: non-universality
In order to calculate all correlation functions for the constrained matrix model with an arbitrary potential $`V(M)=_{j=1}^{\mathrm{}}\frac{g_j}{j}M^j`$ we introduce an auxiliary potential $`W(M)`$ inside the partition function<sup>1</sup><sup>1</sup>1We added a trivial factor of $`n^2`$ inside the delta function.
$`𝒵_\delta `$ $``$ $`{\displaystyle 𝒟M\mathrm{exp}[n\text{Tr}W(M)]\delta \left(n^2A^2n\text{Tr}V(M)\right)},`$ (2.2)
$`W(M){\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{t_j}{j}}M^j,`$
where the two sets of variables $`\{t_j\}`$ and $`(\{g_j\},A)`$ are taken to be independent. All $`k`$-point resolvent operators can then be obtained by taking functional derivatives of $`𝒵_\delta `$ with respect to $`W(p)`$ as given below, where we then eventually set the auxiliary potential $`W`$ to zero at the end. A similar trick has been used in ref. in order to investigate a multi-trace random matrix ensemble , showing their non-universality as well. Furthermore we use the complex representation of the $`\delta `$-function to obtain
$$𝒵_\delta =\frac{d\alpha }{2\pi }𝒟M\mathrm{exp}\left[n\text{Tr}\left(W(M)+i\alpha (V(M)A^2)\right)\right].$$
(2.3)
If we had used instead the following representation of the $`\delta `$-function, $`\delta (x)=lim_l\mathrm{}\sqrt{\pi /l}\mathrm{exp}[lx^2]`$, we would have obtained the trace squared ensemble: $`𝒫_l(M)\mathrm{exp}[2lnA^2\text{Tr}V(M)l\text{Tr}V(M)^2]`$, with the strength of the touching interaction being proportional to $`l`$. However, it is not straightforward to derive and solve loop equations for such an ensemble. In particular, we cannot directly employ the non-universality results of , where it was crucial that the single- and multi-trace potentials were different. It is in the form eq. (2.3) that we can actually derive and solve the loop equations for the constrained model. Throughout the paper the same notation as in is used, which is redisplayed here for completeness. The resolvent or 1-loop correlator is defined as
$$G(p)\frac{1}{n}\text{Tr}\frac{1}{pM}_\delta =\frac{1}{n}\underset{k=0}{\overset{\mathrm{}}{}}\frac{\text{Tr}M^k_\delta }{p^{k+1}}$$
(2.4)
and higher $`k`$-point resolvents are given by
$$G(p_1,\mathrm{},p_k)n^{k2}\text{Tr}\frac{1}{p_1M}\mathrm{}\text{Tr}\frac{1}{p_kM}_{\delta ,conn}$$
(2.5)
where $`conn`$ stands for the connected part of the expectation value with respect to eq. (2.3). They are defined such that the leading part is of the order O(1). If we define the free energy $`_\delta `$ as follows
$$𝒵_\delta \mathrm{exp}[n^2_\delta ]$$
(2.6)
all resolvents can be obtained from it by successive applications of the loop insertion operator
$`{\displaystyle \frac{d}{dW}}(p)`$ $``$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{j}{p^{j+1}}}{\displaystyle \frac{d}{dt_j}},`$ (2.7)
$`G(p_1,\mathrm{},p_k)`$ $`=`$ $`{\displaystyle \frac{d}{dW}}(p_k){\displaystyle \frac{d}{dW}}(p_{k1})\mathrm{}{\displaystyle \frac{d}{dW}}(p_1)_\delta +\delta _{k,1}{\displaystyle \frac{1}{p_1}}.`$ (2.8)
In particular all higher resolvents can be derived from the 1-point resolvent alone.
### 2.1 The loop equation in the planar limit
The loop equation is derived in the usual way from the partition function eq. (2.3) by shifting variables $`MM+ϵ/(pM)`$ and requiring it to be invariant under this shift, i.e. $`\frac{d𝒵_\delta }{dϵ}|_{ϵ=0}=0`$,
$$_𝒞\frac{d\omega }{2\pi i}\frac{V_{eff}^{}(\omega )}{p\omega }G(\omega )=G(p)^2+\frac{1}{n^2}\frac{d}{dW}(p)G(p),$$
(2.9)
where we have defined the effective potential
$$V_{eff}(M)=W(M)+i\overline{\alpha }V(M).$$
(2.10)
In eq.(2.9) the integration contour $`𝒞`$ encircles the support $`[y,x]`$ of the spectral density $`\rho (\lambda )`$ counterclockwise in the complex plane, not including the argument $`p[y,x]`$. The parameter $`i\overline{\alpha }`$ inside the effective potential $`V_{eff}`$ is determined by the constraint $`\text{Tr}V(M)=nA^2`$ as a function of all coupling constants, as we will see in more detail below. This constraint can also obtained by requiring the invariance of $`𝒵_\delta `$ under the shift $`\alpha \alpha +ϵ`$.
Let us stress again that the resolvents are given by differentiating with respect to $`W(p)`$ and not the effective potential $`V_{eff}(p)`$. Due to the $`\alpha `$-integral in the partition function $`𝒵_\delta `$ eq. (2.3) we have $`\text{Tr}M^k_\delta i\alpha \text{Tr}M^k_\delta `$ which is needed to determine $`G(p)`$ eq. (2.4). Furthermore let us note that eq. (2.9) looks almost identical to the loop equation of the unconstrained hermitian matrix model defined in the next section eq. (3.1). However, the role of the above mentioned auxiliary potential $`W`$ as well as the constraint will modify the results of . The constraint leads to similar complications as in the situation where the support of the spectral density consists of several intervals .
In order to solve the loop equation we introduce a $`1/n^2`$ expansion for all $`k`$-point resolvent operators (2.5)
$$G(p_1,\mathrm{},p_k)=\underset{g=0}{\overset{\mathrm{}}{}}\frac{1}{n^{2g}}G_g(p_1,\mathrm{},p_k),$$
(2.11)
where the leading part with genus $`g=0`$ (planar) is of the order O(1). In ref. it was shown that for monomial potentials the expectation values of the RTEs possess such an expansion in $`1/n^2`$ as well. Here we assume that the same holds true for all polynomials potentials. Inserting this expansion into the loop equation (2.9) and taking the large-$`n`$ limit we obtain to leading order
$$_𝒞\frac{d\omega }{2\pi i}\frac{V_{eff}^{}(\omega )}{p\omega }G_0(\omega )=G_0(p)^2.$$
(2.12)
If we make the Ansatz that $`G_0(p)`$ has just one cut in the complex plane or equivalently the support of the eigenvalues consists of the single interval $`[y,x]`$ we obtain
$$G_0(p)=\frac{1}{2}\left(V_{eff}^{}(p)(p)\sqrt{(px)(py)}\right),$$
(2.13)
where the analytic function $`(p)`$ is given by
$$(p)=_𝒞_{\mathrm{}}\frac{d\omega }{2\pi i}\frac{V_{eff}^{}(\omega )}{(\omega p)\sqrt{(\omega x)(\omega y)}}.$$
(2.14)
For details of the derivation see for example ref. . The final result can be written as follows, after deforming back the integration contour,
$$G_0(p)=\frac{1}{2}_𝒞\frac{d\omega }{2\pi i}\frac{V_{eff}^{}(\omega )}{p\omega }\sqrt{\frac{(px)(py)}{(\omega x)(\omega y)}}.$$
(2.15)
To make the solution complete we still have to determine the endpoints $`x`$ and $`y`$ as well as the parameter $`i\overline{\alpha }`$ in terms of the coupling constants of $`W`$ and $`V`$ and the parameter $`A`$. The first two equations can be obtained from the asymptotic behavior of $`G_0(p)`$. According to the definition (2.4) we have
$$\underset{p\mathrm{}}{lim}G(p)=\frac{1}{p}.$$
(2.16)
Since the leading term does not depend on $`n`$ it comes from the planar part $`G_0(p)`$ and we obtain the conditions
$`0`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}{\displaystyle \frac{V_{eff}^{}(\omega )}{\sqrt{(\omega x)(\omega y)}}},`$
$`1`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}{\displaystyle \frac{\omega V_{eff}^{}(\omega )}{\sqrt{(\omega x)(\omega y)}}}.`$ (2.17)
The third equation needed we obtain from the constraint on $`\text{Tr}V(M)`$ which we rewrite in terms of the spectral density
$`\rho (\lambda )`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}\underset{ϵ0}{lim}\left(G_0(\lambda iϵ)G_0(\lambda +iϵ)\right)`$
$`=`$ $`{\displaystyle \frac{1}{2\pi }}(\lambda )\sqrt{(x\lambda )(\lambda y)}.`$ (2.18)
The constraint then reads
$$A^2=_y^x𝑑\lambda \rho (\lambda )V(\lambda )=_y^x\frac{d\lambda }{2\pi }(\lambda )\sqrt{(x\lambda )(\lambda y)}V(\lambda ),$$
(2.19)
which determines the parameter $`i\overline{\alpha }`$ contained in $`(\lambda )`$. This equation together with the boundary conditions eq. (2.17) determines the planar resolvent $`G_0(p)`$ eq. (2.4) completely as a function of the coupling constants of $`W`$ and $`V`$ and of the parameter $`A`$.
### 2.2 Higher planar $`k`$-point resolvents
Starting from the the planar resolvent $`G_0(p)`$ we can obtain all higher planar $`k`$-point resolvents by successively applying the loop insertion operator $`\frac{d}{dW}`$ to it, as it is given in eq. (2.8). Here we use the fact that all resolvents have the same expansion in $`1/n^2`$ eq. (2.11). For this purpose we introduce a set of new parameters $`M_k`$ and $`J_k`$, $`k`$ N<sub>+</sub>. These moments usually play the role of universal parameters of the higher $`k`$-point resolvents encoding all the information of the potential in addition to the endpoints $`x`$ and $`y`$ . We will then rewrite the loop insertion operator in terms of these new variables. This is done in order to make the successive application of $`\frac{d}{dW}`$ to an algebraic procedure. Finally we calculate explicitly the non-universal planar 2-point resolvent $`G_0(p,q)`$ and comment on the general situation.
Let us begin by defining
$`M_k`$ $``$ $`{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}V_{eff}^{}(\omega ){\displaystyle \frac{\varphi (\omega )}{(\omega x)^k}},\varphi (\omega ){\displaystyle \frac{1}{\sqrt{(\omega x)(\omega y)}}},`$
$`J_k`$ $``$ $`{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}V_{eff}^{}(\omega ){\displaystyle \frac{\varphi (\omega )}{(\omega y)^k}},k=1,2,\mathrm{}.`$ (2.20)
Expanding the poles at $`x`$ and $`y`$ the moments can be explicitly written as functions of the coupling constants. Because of $`M_k=\frac{1}{(k1)!}\frac{d^{k1}}{d\lambda ^{k1}}M(\lambda )|_{\lambda =x}`$ and similarly for $`y`$ the moments also characterize the multi-critical points of the model. We now rewrite the loop insertion operator eq. (2.7) in the following way
$`{\displaystyle \frac{d}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{}{W}}(p)+{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{}{x}}+{\displaystyle \frac{dy}{dW}}(p){\displaystyle \frac{}{y}}+{\displaystyle \frac{di\overline{\alpha }}{dW}}(p){\displaystyle \frac{}{i\overline{\alpha }}}`$ (2.21)
$`+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}\left({\displaystyle \frac{dM_k}{dW}}(p){\displaystyle \frac{}{M_k}}+{\displaystyle \frac{dJ_k}{dW}}(p){\displaystyle \frac{}{J_k}}\right),`$
$`{\displaystyle \frac{}{W}}(p)`$ $``$ $`{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{j}{p^{j+1}}}{\displaystyle \frac{}{t_j}}.`$ (2.22)
While the $`\frac{dM_k}{dV}(p)`$ and $`\frac{dJ_k}{dV}(p)`$ can be obtained in a straightforward way from the definition eq. (2.20) (see e.g. ) the remaining unknown quantities are derived by applying $`\frac{d}{dW}(p)`$ to the boundary conditions eqs. (2.17) and (2.19) and solving a linear set of equations. This is done in the Appendix A with the result reading
$`M_1{\displaystyle \frac{dx}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{\varphi (p)}{px}}+{\displaystyle \frac{4}{B(xy)}}(G_0(p)\varphi (p)),`$
$`J_1{\displaystyle \frac{dy}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{\varphi (p)}{py}}{\displaystyle \frac{4}{B(xy)}}(G_0(p)\varphi (p)),`$
$`{\displaystyle \frac{1}{i\overline{\alpha }}}{\displaystyle \frac{di\overline{\alpha }}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{1}{B}}(G_0(p)\varphi (p)),`$ (2.23)
where
$`B`$ $``$ $`i\overline{\alpha }{\displaystyle _y^x}d\lambda V(\lambda ){\displaystyle \frac{1}{2\pi i}}\underset{ϵ0}{lim}\times `$
$`\times \left[G_0(\lambda iϵ)\varphi (\lambda iϵ)(G_0(\lambda +iϵ)\varphi (\lambda +iϵ))\right].`$
All quantities are expressed by elementary functions and the planar resolvent $`G_0(p)`$ eq. (2.15) and we have set already the auxiliary potential $`W0`$. The result for general $`W`$ can be derived from eqs. (A.2) and (A.3).
We are now ready to apply the loop insertion operator $`\frac{d}{dW}`$ in the form (2.21) to $`G_0(p)`$ eq. (2.15), which does not explicitly depend on the moments:
$`G_0(p,q)`$ $`=`$ $`{\displaystyle \frac{d}{dW}}(q)G_0(p)`$ (2.25)
$`=`$ $`{\displaystyle \frac{1}{2(pq)^2}}\left({\displaystyle \frac{\varphi (q)}{\varphi (p)}}1\right)+{\displaystyle \frac{1}{4(qp)}}{\displaystyle \frac{\varphi (q)}{\varphi (p)}}\left({\displaystyle \frac{1}{qx}}+{\displaystyle \frac{1}{qy}}\right)`$
$`+{\displaystyle \frac{di\overline{\alpha }}{dW}}(q){\displaystyle \frac{1}{2}}{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}{\displaystyle \frac{V^{}(\omega )}{p\omega }}\sqrt{{\displaystyle \frac{(px)(py)}{(\omega x)(\omega y)}}}`$
$`+{\displaystyle \frac{1}{4\varphi (p)}}\left({\displaystyle \frac{1}{px}}M_1{\displaystyle \frac{dx}{dW}}(q)+{\displaystyle \frac{1}{py}}J_1{\displaystyle \frac{dy}{dW}}(q)\right),`$
after performing some contour integrals. If we set $`W0`$ we can use the results eq. (2.23) and finally obtain
$$G_0(p,q)=G_0^{can}(p,q)\frac{1}{B}(G_0(p)\varphi (p))(G_0(q)\varphi (q)),$$
(2.26)
the planar connected 2-point resolvent of the constrained matrix model. The first part is the well known universal 2-point resolvent of the corresponding unconstrained or canonical matrix model (can)
$$G_0^{can}(p,q)=\frac{1}{4(pq)^2}\left(\frac{(px)(qy)+(py)(qx)}{\sqrt{(px)(py)(qx)(qy)}}2\right)$$
(2.27)
whereas the second part contains the non-universal terms $`G_0(p)`$ and $`G_0(q)`$. Still, the result is given in closed form for an arbitrary polynomial potential $`V(M)`$. Eq. (2.26) can be compared with the earlier result
$$G_0(p,q)=G_0^{can}(p,q)\frac{1}{2p}_p\left(pG_0(p)\right)_q\left(qG_0(q)\right)$$
(2.28)
for the special case of monomial potentials $`V(M)=M^{2p}`$, $`p`$ N<sub>+</sub>. For $`p=1,2`$ we have checked explicitly that the two results eq. (2.26) and eq. (2.28) agree. The corresponding resolvents can be found in <sup>2</sup><sup>2</sup>2In eq. (2.18) of ref. the factor of $`1/2p`$ is missing. and the parameter $`i\overline{\alpha }`$ was already determined in . From the procedure described above it is clear that also the higher $`k`$-point resolvents will remain non-universal since the derivative $`\frac{d}{dW}(p_i)G_0(p_j)`$ always contains terms proportional to $`G_0(p_i)`$ and $`G_0(p_j)`$.
Let us also briefly comment on higher genus contributions. Expanding the loop equation in $`1/n^2`$ together with eq. (2.11) one obtains for genus one
$$\left(\widehat{𝒦}2G_0(p)\right)G_1(p)=G_0(p,p),$$
(2.29)
where
$$\widehat{𝒦}f(p)_𝒞\frac{d\omega }{2\pi i}\frac{V_{eff}^{}(\omega )}{p\omega }f(\omega ).$$
(2.30)
The right hand side of eq. (2.29) is easily obtained from eq. (2.26) by setting $`p=q`$. However, it is no longer a rational function in contrast to $`G_0^{can}(p,p)`$. Its non-universality will then translate to $`G_1(p)`$ after inverting the integral operator $`(\widehat{𝒦}2G_0(p))`$ and thus to higher genera through
$$\left(\widehat{𝒦}2G_0(p)\right)G_g(p)=\underset{g=1}{\overset{g1}{}}G_g(p)G_{gg}(p)+\frac{d}{dW}(p)G_{g1}(p),g1.$$
(2.31)
For this reason we do not go through the tedious procedure of finding a basis for $`(\widehat{𝒦}2G_0(p))`$ now including also square roots and inverting it.
## 3 The microscopic limit: universality
In this section we investigate correlations of eigenvalues at the distance of the mean level spacing $`D1/n`$, the so-called microscopic large-$`n`$ limit. We will heavily exploit the knowledge about correlations in the unconstrained canonical model for finite as well as infinite $`n`$. Our main result is an explicit relation between the canonical and RTE correlations for finite-$`n`$ which then serves to determine the microscopic RTE correlations and prove their universality.
Let us recall the known results about the canonical ensemble which also fixes our notation. The partition function reads
$$𝒵𝒟M\mathrm{exp}[n\text{Tr}\stackrel{~}{V}(M)],$$
(3.1)
where $`\stackrel{~}{V}(M)`$ is a polynomial. For the generalized RTE or micro-canonical ensemble
$$𝒵_\delta 𝒟M\delta \left(A^2\frac{1}{n}\text{Tr}V(M)\right),$$
(3.2)
where we do not need to introduce an auxiliary potential in contrast to the previous section. In order to obtain the same microscopic correlations for the two models we will eventually have to relate the coupling constants of the respective polynomial potentials $`V`$ and $`\stackrel{~}{V}`$ (as in the macroscopic limit, see ). The $`k`$-point density correlation function is defined as
$$\rho (\lambda _1,\mathrm{},\lambda _k)\frac{1}{n}\text{Tr}\delta (\lambda _1M)\mathrm{}\frac{1}{n}\text{Tr}\delta (\lambda _kM),$$
(3.3)
and similarly for the delta-measure. Its connected part $`conn`$ is related in the following way to the $`k`$-point resolvent defined in the previous section eq. (2.5)
$`\rho ^c(\lambda _1,\mathrm{},\lambda _k)`$ $``$ $`{\displaystyle \frac{1}{n}}\text{Tr}\delta (\lambda _1M)\mathrm{}{\displaystyle \frac{1}{n}}\text{Tr}\delta (\lambda _kM)_{conn},`$
$`=`$ $`{\displaystyle \frac{1}{n^{2k2}}}\left({\displaystyle \frac{1}{2\pi i}}\right)^k\times `$
$`\times \underset{ϵ0}{lim}{\displaystyle \underset{\sigma _i=\pm }{}}\left({\displaystyle \underset{i}{}}\sigma _i\right)G(\lambda _1+\sigma _1iϵ,\mathrm{},\lambda _k+\sigma _kiϵ),`$
where we have not yet taken the large-$`n`$ limit. In this section we deal with density correlations instead of resolvents because they can be given more explicitly for finite-$`n`$. Namely in the canonical model all density correlators can be expressed in terms of the Kernel $`K_n(\lambda ,\mu )`$ of a set of orthonormal polynomials $`P_l(\lambda )`$ at finite-$`n`$
$`\rho (\lambda _1,\mathrm{},\lambda _k)`$ $`=`$ $`\underset{1i,jk}{det}\left[K_n(\lambda _i,\lambda _j)\right],`$ (3.5)
$`\rho ^c(\lambda _1,\mathrm{},\lambda _k)`$ $`=`$ $`(1)^{k+1}{\displaystyle \underset{P}{}}K_n(\lambda _1,\lambda _2)K_n(\lambda _2,\lambda _3)\mathrm{}K_n(\lambda _k,\lambda _1),`$
where the sum is taken over the $`(k1)!`$ distinct cyclic permutation $`P`$ of the indices $`(1,2,\mathrm{},k)`$. The kernel and the polynomials are defined as
$`K_n(\lambda ,\mu )`$ $`=`$ $`{\displaystyle \frac{1}{n}}e^{\frac{n}{2}(\stackrel{~}{V}(\lambda )+\stackrel{~}{V}(\mu ))}{\displaystyle \underset{k=0}{\overset{n1}{}}}P_k(\lambda )P_k(\mu ),`$ (3.6)
$`\delta _{kl}`$ $`=`$ $`{\displaystyle 𝑑\lambda e^{n\stackrel{~}{V}(\lambda )}P_k(\lambda )P_l(\lambda )}.`$
The use of the orthogonal polynomial method is only possible for the canonical model because the measure $`\mathrm{exp}[n\text{Tr}\stackrel{~}{V}(M)]`$ inside the partition function eq. (3.1) factorizes in terms of the eigenvalues $`\lambda _i`$ of the hermitian matrix $`M`$. For the RTE no such property holds which forces us to seek for other methods. Here we will make use of homogeneity properties for monomial potentials.
Let us finally give the universal results for the canonical ensemble in the microscopic large-$`n`$ limit, where we restrict ourselves to the origin scaling limit due to the local translational invariance of the canonical ensemble. As mentioned in the beginning we measure eigenvalues in units of the mean level spacing which is $`D=1/(n\rho (0))`$ at the origin. Here $`\rho (0)`$ is the mean eigenvalue density at zero, taken in the macroscopic large-$`n`$ limit as given in eq. (2.18). We then define new variables $`z_i=\lambda _i/D`$ which are kept fixed in the large-$`n`$ limit. Since $`D0`$ as $`n\mathrm{}`$ the variables $`\lambda _i`$ have to go to zero as well. In this particular limit the microscopic correlators are defined as<sup>3</sup><sup>3</sup>3The factor $`n^k`$ appears as we had already defined the $`k`$-point correlator in eq. (3.3) to be normalized to unity. The appropriate unfolding procedure is usually defined for un-normalized correlators (see e.g. ), which provides us with the correct pre-factor $`1/\rho (0)^k`$.
$$\rho _S(z_1,\mathrm{},z_k)\underset{n\mathrm{}}{lim}(Dn)^k\rho (z_1D,\mathrm{},z_kD).$$
(3.7)
This limit, which is well behaved and finite, can be investigated by using the Darboux-Christoffel formula for the kernel eq. (3.6) and the asymptotic large-$`n`$ behavior of the polynomials $`P_k(\lambda )`$ to obtain
$$\underset{n\mathrm{}}{lim}DnK_n(z_1D,z_2D)=\frac{\mathrm{sin}(\pi (z_1z_2))}{\pi (z_1z_2)}.$$
(3.8)
This is the universal sine-law which is valid for all polynomial potentials $`\stackrel{~}{V}(\lambda )`$. Together with eqs. (3.5) and (3.7) it completely determines all $`k`$-point density correlators, connected and not-connected, where at coinciding arguments we have $`lim_n\mathrm{}DnK_n(zD,zD)=1`$,
$$\rho _S(z_1,\mathrm{},z_k)=\underset{1i,jk}{det}\left[\frac{\mathrm{sin}(\pi (z_iz_j))}{\pi (z_iz_j)}\right],$$
(3.9)
and similarly for $`\rho _S^c(z_1,\mathrm{},z_k)`$.
We conclude with the following remark. If we had taken the macroscopic large-$`n`$ limit instead, the connected correlators $`\rho ^c(z_1,\mathrm{},z_k)`$ had been of the order O$`(1/n^{2k2})`$ as one can see from eq. (3) together with the fact that the (connected) resolvents are of order O(1). However, when taking the microscopic limit keeping the $`z_i`$ fixed, the asymptotic kernel eq. (3.8) is of order O(1) and hence are the connected and not-connected microscopic $`k`$-point correlators from eq. (3.5). Consequently, the knowledge of both connected and not-connected correlators, is equivalent here since they can be obtained from each other by adding or subtracting $`l`$-point correlators $`(l<k)`$ of order O(1). In that sense the microscopic large-$`n`$ limit modifies the usual large-$`n`$ factorization of correlation functions.
### 3.1 Microscopic $`k`$-point correlation functions
As it has been mentioned already the correlation functions in the RTE can not be calculated using orthogonal polynomials because of the Dirac $`\delta `$-function in the measure. However, in the particular case of purely monomial potentials $`V(M)=M^{2p}`$, the evaluation of the connected $`k`$-point correlator in the microscopic limit is straightforward. We shall exploit some homogeneity properties of this case, which make it possible to relate ensemble averages in the monomial RTE to ensemble averages of the same quantities in the canonical ensemble<sup>4</sup><sup>4</sup>4This technique is a slight generalization of ref. , where only the macroscopic limit was investigated. with potential $`\stackrel{~}{V}(M)=gM^{2p}`$. Let us define a function $`F_k(M;\lambda )`$ of the matrix $`M`$ and of a set of parameters $`\lambda =(\lambda _1,\mathrm{},\lambda _k)`$, that satisfies the following homogeneity property under a rescaling of the matrix
$$F_k(tM;\lambda )=t^aF_k(M;t^b\lambda )\text{for}t,a,b\text{R}.$$
(3.10)
A simple example for such a function is the operator $`\text{Tr}\delta (\lambda M)/n`$ which has $`a=b=1`$. Also any such product is a homogeneous function, as inside the average of the $`k`$-point correlation function eq. (3.3), having $`a=k`$ and $`b=1`$. In appendix B we derive the following formula for such homogeneous functions $`F_k`$, which relates their canonical and RTE average
$$F_k(M;\lambda )_\delta =\frac{(gn^2)^{\frac{a}{2p}}\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}{A^{\frac{n^2}{p}2}}^1\left[\frac{F_k(M;\left(\frac{gn^2}{t}\right)^{\frac{b}{2p}}\lambda )}{t^{\frac{n^2+a}{2p}}}\right](A^2).$$
(3.11)
Here $`^1[h(t)](x)`$ is the inverse Laplace transform of a function $`h(t)`$, evaluated at the point $`x>0`$ (for an integral representation see eq. (B.5)). Eq. (3.11) holds for any finite $`n`$. If we choose the 1- or 2-point correlator from eq. (3.3) as an example we reproduce the finite-$`n`$ results of ref. for $`\rho _\delta (\lambda )`$ and $`\rho _\delta (\lambda ,\mu )`$ which were given in the case of a Gaussian potential $`p=1`$. If we choose for general $`k`$ to take $`F_k(M;\lambda )_\delta =\rho _\delta (\lambda _1,\mathrm{},\lambda _k)`$, which obviously fulfills the criterion (3.10), we obtain from eq. (3.11) the following expression for the $`k`$-point RTE correlator
$$\rho _\delta (\lambda _1,\mathrm{},\lambda _k)=\frac{(gn^2)^{\frac{k}{2p}}\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}{A^{\frac{n^2}{p}2}}^1[\frac{\rho (\left(\frac{gn^2}{t}\right)^{\frac{1}{2p}}\lambda _1,\mathrm{},\left(\frac{gn^2}{t}\right)^{\frac{1}{2p}}\lambda _k)}{t^{\frac{n^2k}{2p}}}](A^2)$$
(3.12)
Now we make use of the fact that at finite $`n`$ the correlations $`\rho (\lambda _1,\mathrm{},\lambda _k)`$ from eq. (3.5) can be written as a polynomial in all variables $`\lambda _i`$ times an exponential measure factor:
$$\rho (\lambda _1,\mathrm{},\lambda _k)e^{ng_{i=1}^k\lambda _i^{2p}}\underset{l_1,\mathrm{},ł_k=0}{\overset{2n2}{}}c_{\{l_1,\mathrm{},l_k\}}^{(n)}\lambda _1^{l_1}\mathrm{}\lambda _k^{l_k}.$$
(3.13)
Due to this fact we can actually perform the inverse Laplace transformation
$`\rho _\delta (\lambda _1,\mathrm{},\lambda _k)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}{(gn^2)^{\frac{k}{2p}}A^{\frac{n^2}{p}2}}}{\displaystyle \underset{l_1,\mathrm{},l_k=0}{\overset{2n2}{}}}{\displaystyle \frac{c_{\{l_1,\mathrm{},l_k\}}^{(n)}}{(gn^2)^{\frac{1}{2p}\mathrm{\Sigma }_il_i}}}\lambda _1^{l_1}\mathrm{}\lambda _k^{l_k}\times `$
$`\times ^1\left[t^{\frac{n^2k\mathrm{\Sigma }_il_i}{2p}}\right](A^2{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{k}{}}}\lambda _i^{2p})`$
$`=`$ $`{\displaystyle \frac{\theta \left(A^2\frac{1}{n}_{i=1}^k\lambda _i^{2p}\right)}{(gn^2)^{\frac{k}{2p}}A^{\frac{n^2}{p}2}}}{\displaystyle \underset{l_1,\mathrm{},l_k=0}{\overset{2n2}{}}}{\displaystyle \frac{c_{\{l_1,\mathrm{},l_k\}}^{(n)}\lambda _1^{l_1}\mathrm{}\lambda _k^{l_k}}{(gn^2)^{\frac{1}{2p}\mathrm{\Sigma }_il_i}}}\times `$
$`\times {\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}{\mathrm{\Gamma }\left(\frac{n^2k\mathrm{\Sigma }_il_i}{2p}\right)}}\left(A^2{\displaystyle \frac{1}{n}}{\displaystyle \underset{i=1}{\overset{k}{}}}\lambda _i^{2p}\right)^{\frac{n^2k\mathrm{\Sigma }_il_i}{2p}1},`$ (3.14)
where we have used the shift property $`^1[h(t)e^{\sigma t}](x)=^1[h(t)](x+\sigma )`$, the linearity of the inverse Laplace transform and eq. (B.7) of Appendix B. Eq. (3.14) is our first result, the finite-$`n`$ $`k`$-point correlation function for RTEs with monomial potential $`V(M)=M^{2p}`$ in terms of the corresponding canonical correlator at finite-$`n`$ with potential $`\stackrel{~}{V}(M)=gM^{2p}`$. In general, eqs. (3.13) and (3.14) are different from each other at finite-$`n`$. This remains true in the macroscopic large-$`n`$ limit, as we have seen in the previous section. In the remaining part we will show that in the microscopic large-$`n`$ limit, however, they happen to coincide.
In a first step, we determine the mean level spacing $`D_\delta =1/(n\rho _\delta (0))`$ in order to define the appropriate microscopic scaling limit. From eq. (3.14) at $`\lambda =0`$ one can read off $`\rho _\delta (0)`$, since then the sum over all $`l_i`$ collapses. We obtain
$$\rho _\delta (0)=\frac{c_{\{0\}}^{(n)}}{(gn^2A^2)^{\frac{1}{2p}}}\frac{\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}{\mathrm{\Gamma }\left(\frac{n^21}{2p}\right)}\frac{c_{\{0\}}}{(2pgA^2)^{\frac{1}{2p}}}$$
(3.15)
for its large-$`n`$ value, with $`c_{\{0\}}=\rho (0)`$ being the macroscopic large-$`n`$ limit of the canonical spectral density eq. (2.18) at the origin (which exists and is finite). In order to have the same mean level spacing as in the canonical ensemble we would have to set $`2pgA^2=1`$. This identification of coupling constants occurs also in the macroscopic large-$`n`$ limit in order to match the corresponding macroscopic spectral densities (see refs. ). However, since we measure all correlations in units of $`D`$ and $`D_\delta `$ respectively we do not need to identify $`D=D_\delta `$ since they drop out in the microscopic correlators anyway, as we will see below.
We will now take the microscopic limit analogue to eq. (3.7) with rescaling $`\lambda _i=z_iD_\delta `$ of our finite-$`n`$ relation (3.14). The large-$`n`$ limit of the different factors can be obtained as follows, starting with the $`\theta `$-function term in eq. (3.14)
$$\theta \left(A^2\frac{2pgA^2}{n^{2p+1}}\underset{i=1}{\overset{k}{}}\left(\frac{z_i}{c_{\{0\}}}\right)^{2p}\right)\theta (A^2)=1.$$
(3.16)
The remaining $`A`$-dependent terms yield
$`A^{\frac{n^2}{p}+2}\left(A^2{\displaystyle \frac{2pgA^2}{n^{2p+1}}}{\displaystyle \underset{i=1}{\overset{k}{}}}\left({\displaystyle \frac{z_i}{c_{\{0\}}}}\right)^{2p}\right)^{\frac{n^2k\mathrm{\Sigma }_il_i}{2p}1}=`$
$`=A^{\frac{1}{p}(k+\mathrm{\Sigma }_il_i)}\left(1\left[{\displaystyle \frac{n^2k\mathrm{\Sigma }_il_i}{2p}}1\right]{\displaystyle \frac{2pg}{n^{2p+1}}}{\displaystyle \underset{i=1}{\overset{k}{}}}\left({\displaystyle \frac{z_i}{c_{\{0\}}}}\right)^{2p}+\mathrm{}\right)`$
$`=A^{\frac{1}{p}(k+\mathrm{\Sigma }_il_i)}\left(1+\text{O}({\displaystyle \frac{1}{n^{2p1}}})\right),`$ (3.17)
which still has to be evaluated under the sum over $`l_i`$’s. Here we have used in the second step that the first factor inside the parenthesis is of O$`(n^2)`$ since
$$\underset{i=1}{\overset{k}{}}l_ik(2n2),$$
(3.18)
is at most of O$`(n)`$. Because of $`p1`$ the corrections are sub-leading. The factor containing $`\mathrm{\Gamma }`$-functions we evaluate together with the explicit factors of $`n`$ in eq. (3.14)
$`{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}{(gn^2)^{\frac{1}{2p}(k+\mathrm{\Sigma }_il_i)}\mathrm{\Gamma }\left(\frac{n^2}{2p}\frac{(k+\mathrm{\Sigma }_il_i)}{2p}\right)}}=`$ (3.19)
$`=(2pg)^{\frac{1}{2p}(k+\mathrm{\Sigma }_il_i)}\left(1{\displaystyle \frac{1}{4pn^2}}(k+{\displaystyle \underset{i=1}{\overset{k}{}}}l_i)(2p+k+{\displaystyle \underset{i=1}{\overset{k}{}}}l_i)+\mathrm{}\right).`$
It remains to be shown that the second term and thus higher terms in the expansion are sub-leading. A naive counting from eq. (3.18) suggests that this might not be the case. In order to use the microscopic results for the canonical correlators eq. (3.9) we put together our results obtained so far
$`\underset{n\mathrm{}}{lim}(D_\delta n)^k\rho _\delta (z_1D_\delta ,\mathrm{},z_kD_\delta )=`$
$`=\underset{n\mathrm{}}{lim}{\displaystyle \frac{(2pgA^2)^{\frac{k}{2p}}}{(c_{\{0\}})^k}}{\displaystyle \underset{l_1,\mathrm{},l_k=0}{\overset{2n2}{}}}c_{\{l_1,\mathrm{},l_k\}}^{(n)}\left({\displaystyle \frac{z_1(2pgA^2)^{\frac{1}{2p}}}{nc_{\{0\}}}}\right)^{l_1}\mathrm{}\left({\displaystyle \frac{z_k(2pgA^2)^{\frac{1}{2p}}}{nc_{\{0\}}}}\right)^{l_k}`$
$`\times (2pgA^2)^{\frac{1}{2p}(k+\mathrm{\Sigma }_il_i)}\left(1{\displaystyle \frac{1}{4pn^2}}(k+{\displaystyle \underset{i=1}{\overset{k}{}}}l_i)(2p+k+{\displaystyle \underset{i=1}{\overset{k}{}}}l_i)+\mathrm{}\right)`$
$`=\underset{n\mathrm{}}{lim}\left(1{\displaystyle \frac{1}{4pn^2}}({\displaystyle \underset{i=1}{\overset{k}{}}}_{z_i}z_i)(2p+{\displaystyle \underset{i=1}{\overset{k}{}}}_{z_i}z_i)+\mathrm{}\right)`$
$`\times {\displaystyle \frac{1}{(c_{\{0\}})^k}}{\displaystyle \underset{l_1,\mathrm{},l_k=0}{\overset{2n2}{}}}c_{\{l_1,\mathrm{},l_k\}}^{(n)}\left({\displaystyle \frac{z_1}{nc_{\{0\}}}}\right)^{l_1}\mathrm{}\left({\displaystyle \frac{z_k}{nc_{\{0\}}}}\right)^{l_k}.`$ (3.20)
Here, the pre-factors from eqs. (3.17) and (3.19) have canceled with the factors of $`(2pgA^2)^{\frac{1}{2p}}`$ from the unfolding. Now we know from the canonical ensemble eqs. (3.13), (3.7) and (3.9) that the limit of the sum over the $`l_i`$ exists and is finite:
$$\rho _S(z_1,\mathrm{},z_k)=\underset{n\mathrm{}}{lim}\frac{1}{(c_{\{0\}})^k}\underset{l_1,\mathrm{},l_k=0}{\overset{2n2}{}}c_{\{l_1,\mathrm{},l_k\}}^{(n)}\left(\frac{z_1}{nc_{\{0\}}}\right)^{l_1}\mathrm{}\left(\frac{z_k}{nc_{\{0\}}}\right)^{l_k}.$$
(3.21)
Hence the term in eq. (3.20) proportional to $`1/(4pn^2)`$ is indeed sub-leading and we have as a final result
$$\rho _{\delta ,S}(z_1,\mathrm{},z_k)=\rho _S(z_1,\mathrm{},z_k),$$
(3.22)
or more explicitly in terms of eq. (3.9). Since our derivation holds for the RTE with an arbitrary monomial potential $`V(M)=M^{2p}`$, we have not only derived all $`k`$point correlation functions but also proved their universality for the given class of potentials. Let us finally point out that the equivalence eq. (3.22) also holds for the corresponding connected $`k`$point correlation functions. As we have mentioned already at the end of the previous subsection, they are of the same order in the microscopic limit and they can be obtained from each other by adding or subtracting lower $`l`$-point correlators.
## 4 Conclusions
We have shown for RTEs as an example of constrained random matrix models that in one and the same model correlation functions may exhibit universal and non-universal behavior in different large-$`n`$ regimes. In particular, in the macroscopic large-$`n`$ limit all planar connected $`k`$-point resolvents are non-universal and a closed expression was given for the resolvent $`G_0(z,w)`$ for an arbitrary potential. Hence when switching from the canonical to the RTE or micro-canonical ensemble, the delta-function constraint destroys the macroscopic universality.
A different behavior appears in the study of correlations of eigenvalues at the scale of the mean level spacing $`1/n`$. Here we recover the sine-law of the canonical ensemble and prove its universality for the class of RTEs with monomial potential. This leads us to conjecture that microscopic universality holds also for more general RTEs. The result in the microscopic limit is not unexpected since a global constraint should not influence the local statistics of eigenvalues.
Acknowledgments: We wish to thank G. Cicuta and L. Molinari for their very enjoyable collaboration and many discussions on work done prior to this publication. Furthermore one of us (G.A.) wishes to thank the Physics Department of Parma for its hospitality extended to him on several occasions. The work of G.V. is supported in part by MURST within the project of “Theoretical Physics of fundamental Interactions”.
## Appendix A The functions $`\frac{dx}{dW}(p),\frac{dy}{dW}(p)`$ and $`\frac{di\overline{\alpha }}{dW}(p)`$
In this Appendix we apply the operator $`\frac{d}{dW}(p)`$ eq. (2.21) to the boundary conditions eqs. (2.17) and (2.19) and solve the linear set of equations for the quantities $`\frac{dx}{dW}(p),\frac{dy}{dW}(p)`$ and $`\frac{di\overline{\alpha }}{dW}(p)`$. Using the identity
$$\frac{}{W}(p)W^{}(\omega )=\frac{1}{(p\omega )^2},$$
(A.1)
we obtain
$`0`$ $`=`$ $`_p\varphi (p)+{\displaystyle \frac{1}{2}}\left(M_1{\displaystyle \frac{dx}{dW}}(p)+J_1{\displaystyle \frac{dy}{dW}}(p)\right)+{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}V^{}(\omega )\varphi (\omega ){\displaystyle \frac{di\overline{\alpha }}{dW}}(p)`$
$`0`$ $`=`$ $`_p(p\varphi (p))+{\displaystyle \frac{1}{2}}\left(xM_1{\displaystyle \frac{dx}{dW}}(p)+yJ_1{\displaystyle \frac{dy}{dW}}(p)\right)`$ (A.2)
$`+{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}\omega V^{}(\omega )\varphi (\omega ){\displaystyle \frac{di\overline{\alpha }}{dW}}(p),`$
from eq. (2.17) and
$`0`$ $`=`$ $`{\displaystyle \frac{C}{2}}\left(M_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{px}}\right)+{\displaystyle \frac{D}{2}}\left(J_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{py}}\right)`$ (A.3)
$`+{\displaystyle \frac{di\overline{\alpha }}{dW}}(p){\displaystyle _y^x}{\displaystyle \frac{d\lambda }{2\pi }}V(\lambda )\sqrt{(x\lambda )(\lambda y)}{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}{\displaystyle \frac{V^{}(\omega )}{\omega \lambda }}\varphi (\omega )`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle _𝒞}{\displaystyle \frac{d\omega }{2\pi i}}{\displaystyle \frac{V^{}(\omega )}{p\omega }}\sqrt{{\displaystyle \frac{(\omega x)(\omega y)}{(px)(py)}}},`$
from eq. (2.19) after some calculation. Here we have introduced
$`C`$ $``$ $`{\displaystyle _y^x}{\displaystyle \frac{d\lambda }{2\pi }}{\displaystyle \frac{V(\lambda )}{x\lambda }}\sqrt{(x\lambda )(\lambda y)},`$
$`D`$ $``$ $`{\displaystyle _y^x}{\displaystyle \frac{d\lambda }{2\pi }}{\displaystyle \frac{V(\lambda )}{y\lambda }}\sqrt{(x\lambda )(\lambda y)}.`$ (A.4)
We note that special care has to be taken when applying $`\frac{d}{dW}(p)`$ to $`(\lambda )`$ which is then no longer an analytic function in $`\lambda `$. Therefore the contour $`𝒞_{\mathrm{}}`$ has to be deformed in eq. (2.14) to contain only the pole and cut in the integrand, and not the new pole introduced by $`\frac{d}{dW}(p)`$ (see also Appendix A in ). The linear set of equations eq. (A.2) and eq. (A.3) simplifies considerably when setting $`W0`$ because of $`V_{eff}(\omega )=i\overline{\alpha }V(\omega )`$ in that case. It is this limit which we will need in order give a closed final expression for the planar 2-point resolvent of the pure delta-measure without the auxiliary potential in eq. (2.2). Using eq. (2.17) it follows
$`0`$ $`=`$ $`M_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{px}}+J_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{py}}`$
$`0`$ $`=`$ $`{\displaystyle \frac{x}{2}}\left(M_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{px}}\right)+{\displaystyle \frac{y}{2}}\left(J_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{py}}\right)+{\displaystyle \frac{2}{i\overline{\alpha }}}{\displaystyle \frac{di\overline{\alpha }}{dW}}(p)`$
$`0`$ $`=`$ $`{\displaystyle \frac{C}{2}}\left(M_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{px}}\right)+{\displaystyle \frac{D}{2}}\left(J_1{\displaystyle \frac{dx}{dW}}(p){\displaystyle \frac{\varphi (p)}{py}}\right)`$ (A.5)
$`+{\displaystyle \frac{1}{i\overline{\alpha }}}A^2{\displaystyle \frac{di\overline{\alpha }}{dW}}(p)+{\displaystyle \frac{1}{i\overline{\alpha }}}(G_0(p)\varphi (p)).`$
This can be easily solved for the desired quantities. We obtain
$`M_1{\displaystyle \frac{dx}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{\varphi (p)}{px}}+{\displaystyle \frac{4}{B(xy)}}(G_0(p)\varphi (p))`$
$`J_1{\displaystyle \frac{dy}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{\varphi (p)}{py}}{\displaystyle \frac{4}{B(xy)}}(G_0(p)\varphi (p))`$
$`{\displaystyle \frac{1}{i\overline{\alpha }}}{\displaystyle \frac{di\overline{\alpha }}{dW}}(p)`$ $`=`$ $`{\displaystyle \frac{1}{B}}(G_0(p)\varphi (p)),`$ (A.6)
as given in eq. (2.23) where the following abbreviation has been introduced
$$B=i\overline{\alpha }\left(A^2+2\frac{C+D}{xy}\right).$$
(A.7)
One can easily convince oneself that it equals the form given in eq. (2.2).
## Appendix B RTE via inverse Laplace transform
In this appendix we derive eq. (3.11) which expresses expectation values with respect to the delta-measure in terms of averages with respect to the canonical measure eq. (3.1) using the inverse Laplace transform. Let $`F_k(M;\lambda )`$ be a function of the Hermitian $`n\times n`$ matrix $`M`$ and of the set of parameters $`\lambda =(\lambda _1,\mathrm{},\lambda _k)`$, such that it satisfies the homogeneity property
$$F_k(tM;\lambda )=t^aF_k(M;t^b\lambda )$$
(B.1)
for some real $`t`$, $`a`$ and $`b`$. In other words $`F_k`$ is a homogeneous function of degree $`a`$ with respect to matrix elements $`M_{ij}`$ and of degree $`(b)`$ with respect to each parameter $`\lambda _i`$. An example for such a function is the operator inside the average of eq. (3.3). The matrix integral considered in this appendix is:
$$I[F_k]=𝒟M\delta \left(A^2\frac{1}{n}\text{Tr}[M^{2p}]\right)F_k(M;\lambda ),$$
(B.2)
with $`p`$ an integer number. It is proportional to the average $`F_k_\delta `$. Introducing the complex representation of the delta function $`2\pi \delta (x)=𝑑y\mathrm{exp}[ixy]`$ and then scaling all matrix elements by a factor $`(gn^2/(iy+0^+))^{1/2p}`$, eq. (B.2) reads:
$$I[F_k]=_{\mathrm{}}^+\mathrm{}\frac{dy}{2\pi }e^{iyA^2}\left(\frac{gn^2}{iy+0^+}\right)^{\frac{n^2}{2p}}𝒟Me^{ng\text{Tr}[M^{2p}]}F_k(\left(\frac{gn^2}{iy+0^+}\right)^{\frac{1}{2p}}M;\lambda )$$
(B.3)
where we have interchanged integrals. By using the homogeneity property (B.1) we can rewrite the matrix integral as a canonical ensemble average, up to the normalization factor $`𝒵`$ eq. (3.1):
$`I[F_k]`$ $`=`$ $`𝒵{\displaystyle _{\mathrm{}}^+\mathrm{}}{\displaystyle \frac{dy}{2\pi }}e^{iyA^2}\left({\displaystyle \frac{gn^2}{iy+0^+}}\right)^{\frac{n^2+a}{2p}}F_k(M;\left({\displaystyle \frac{gn^2}{iy+0^+}}\right)^{\frac{b}{2p}}\lambda )`$
$`=`$ $`𝒵(gn^2)^{\frac{n^2+a}{2p}}^1\left[{\displaystyle \frac{F_k(M;\left(\frac{gn^2}{t}\right)^{\frac{b}{2p}}\lambda )}{t^{\frac{n^2+a}{2p}}}}\right](A^2).`$ (B.4)
In the last step we used the complex representation of the inverse Laplace transform, i.e.
$$^1[h(t)](x)=\frac{1}{2\pi i}_{i\mathrm{}+0^+}^{+i\mathrm{}+0^+}𝑑te^{tx}h(t),x>0.$$
(B.5)
In order to obtain the correct normalization of the delta-average we evaluate eq. (B.4) in the case $`F=1`$ (with $`a=b=0`$):
$$I[1]=(gn^2)^{\frac{n^2}{2p}}\frac{A^{\frac{n^2}{p}2}}{\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)}𝒵,$$
(B.6)
where we used the formula
$$^1\left[\frac{1}{t^{\gamma +1}}\right](x)=\frac{x^\gamma }{\mathrm{\Gamma }(\gamma +1)}\theta (x),\text{Re}(\gamma )>1.$$
(B.7)
Finally, by putting together eqs. (B.4) and (B.6), the ensemble average $`F_k_\delta I[F_k]/I[1]`$ with respect to the measure $`\delta (A^2\text{Tr}[M^{2p}]/n)`$ reads:
$$F_k_\delta =\frac{(gn^2)^{\frac{a}{2p}}}{A^{\frac{n^2}{p}2}}\mathrm{\Gamma }\left(\frac{n^2}{2p}\right)^1\left[\frac{F_k(M;\left(\frac{gn^2}{t}\right)^{\frac{b}{2p}}\lambda )}{t^{\frac{n^2+a}{2p}}}\right](A^2),$$
(B.8)
which is just eq. (3.11). |
warning/0002/cond-mat0002164.html | ar5iv | text | # Spin and charge inhomogeneities in high-𝑇_𝑐 cuprates: Evidence from NMR and neutron scattering experiments
## I Introduction
Understanding the doping, frequency and temperature dependence of the magnetic response in the high-$`T_c`$ cuprates is one of the most challenging problems for the high-$`T_c`$ community. Insight into the nature of the strong antiferromagnetic fluctuations very likely holds the key to the unusual normal state properties of the cuprates , as seen in a variety of experimental techniques. Both nuclear magnetic resonance (NMR) and inelastic neutron scattering (INS) experiments are important tools to probe spin excitations in cuprates . While INS experiments provide insight into the momentum and frequency resolved imaginary part of the spin susceptibility, $`\chi (𝐪,\omega )`$, NMR experiments yield information on the momentum averaged real and imaginary part of $`\chi (𝐪,\omega )`$ in the zero frequency limit. However, in contrast to NMR experiments, INS measurements rely on the existence of large single crystals and often suffer from a rather limited experimental resolution. Therefore, these experimental techniques are complementary in the information provided on spin excitations in the cuprates.
One of the central questions in the cuprate superconductors is whether INS and NMR experiments can be simultaneously understood within a single theoretical scenario. Thus far it is not even been clear whether one can reach agreement between INS and NMR data as far as the order of magnitude of the spin susceptibility is concerned. This problem is caused, in part, by the fact that NMR is sensitive to all magnetic fluctuations, regardless of whether they are related to a pronounced momentum dependence of the spin susceptibility. In contrast, INS typically probes only those parts of the susceptibility which are strongly momentum dependent; the rest are typically attributed to the “background” of the signal, as might be caused by the scattering of neutrons on nuclei and phonons.
The situation has been further complicated by recent INS experiments on YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> which observe a crossover from a substantial incommensuration in the magnetic response at low frequencies, to a commensurate structure in $`\chi `$ at higher frequencies . The question thus arises: does this incommensurate order originate from spatially inhomogeneous correlations between the doped holes , which would preserve a locally commensurate magnetic response, or does it reflect homogeneous incommensuration, most likely due to Fermi surface effects.
In this communication we argue that since NMR experiments probe the local spin environment around a nucleus, they are able to distinguish between a locally commensurate and incommensurate magnetic response. Earlier theoretical studies of NMR experiments on YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> and YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> used a phenomenological form of $`\chi `$ :
$$\chi (𝐪,\omega )=\frac{\alpha \xi ^2}{1+\xi ^2(𝐪𝐐)^2i\omega /\omega _{sf}},$$
(1)
where $`\xi `$ is the magnetic correlation length in units of the lattice constant, $`a_0`$, $`\omega _{\mathrm{sf}}`$ an energy scale characterizing the diffusive spin excitations, $`\alpha `$ an overall temperature independent constant, and $`𝐐`$ is the position of the peak in momentum space which was assumed to be commensurate, i.e. $`𝐐=(\pi ,\pi )`$. Using this expression for $`\chi (𝐪,\omega )`$ a rather detailed quantitative understanding of various NMR data has been reached . In particular, from the analysis of the longitudinal spin lattice relaxation rate of the <sup>63</sup>Cu nuclei, $`1/^{63}T_1`$, and the spin spin relaxation time, $`1/T_{2\mathrm{G}}`$, scaling laws like $`\omega _{\mathrm{sf}}\xi ^z`$ with dynamical critical exponent, $`z`$, have been deduced. It was found that $`z1`$ between a lower crossover temperature, $`T_{}`$, and a higher one, $`T_{cr}`$, whereas $`z2`$ above $`T_{cr}`$. In the temperature range where $`z=1`$ scaling applies, it follows that $`T_1T/T_{2G}`$ is independent of temperature; a conjecture which was experimentally verified by Curro et al. for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> . Below the pseudogap temperature, $`T_{}`$, $`\omega _{\mathrm{sf}}`$ and $`\xi `$ decouple and a quasiparticle and spin pseudogap emerges. Moreover, it was recently shown that the temperature dependence of the incommensurate peak width, as determined in INS experiments on La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> , showed $`z=1`$ scaling over a wide range of temperatures, in agreement with the NMR findings for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. Because of its appearance in both NMR and INS results, we will therefore assume in the following analysis that $`z=1`$ is the proper scaling behavior for magnetically underdoped cuprates between $`T_{}`$ and $`T_{\mathrm{cr}}`$.
One of our central results is that the attempt to understand the NMR data for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> with homogeneous incommensuration is inconsistent with $`z=1`$ scaling. This implies that the local magnetic response, which is probed in NMR experiments, is commensurate. Should INS experiments show that this compound exhibits a (globally) incommensurate magnetic response similar to that seen in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, this would strongly support a dynamic charge and spin inhomogeneity (stripe) origin of the incommensuration. In Sec. V we discuss a spin and charge inhomogeneity scenario which reconciles the global incommensuration with the local commensuration in $`\chi `$.
In a second important result we carry out a quantitative comparison of the strength of the antiferromagnetic spin fluctuations, as measured by NMR and INS experiments. We find that agreement between the results obtained from these quite different experimental techniques, which explore not only a different frequency range, but different wavevector regimes, can be obtained within a factor of 2. Given the large uncertainties in determining the absolute value of $`\chi ^{\prime \prime }`$, we believe that this result demonstrates that a consistent description of INS and NMR data can be achieved within the framework of Eq.(1).
Although the form of $`\chi `$ in Eq.(1) was originally invented to understand the spin response at very low frequencies and above the superconducting transition temperature , it has proved interesting to investigate to what extent one can understand INS data at higher frequencies within the same framework. The results by Aeppli et al. support the picture of a unique incommensurate spin response from zero energy up to $`15\mathrm{meV}`$, the highest energy used in Ref.. Within the error bars of the experiment, only for momentum values away from the peak maximum do systematic deviations from Eq.(1) occur. Such deviations indicate the presence of lattice corrections to the continuum limit used in Eq.(1). Lattice corrections are expected to be extremely important for the local, momentum averaged, susceptibility:
$$\chi _{loc}^{\prime \prime }(\omega )=\frac{1}{4\pi ^2}_{BZ}d^2𝐪\chi ^{\prime \prime }(𝐪,\omega ),$$
(2)
since the phase space of momenta away from the peak maxima is considerable in two dimensions. For higher frequencies spin excitations away from the antiferromagnetic peak and energetically large compared to $`\omega _{\mathrm{sf}}`$ come into play and we therefore expect a poor description of $`\chi _{loc}^{\prime \prime }(\omega )`$ in terms of Eq.(1). This conclusion is independent of whether the peak position commensurate or incommensurate peaks.
Deviations from the universal continuum limit of $`\chi `$ are expected to be also of relevance for the relaxation rates measured in NMR experiments since these are weighted momentum averages of $`\chi (𝐪,\omega )`$. Thus, it is worthwhile to study whether one can find indications for lattice corrections from an analysis of NMR and INS data. Such corrections to $`\chi `$ are also of relevance for our understanding of the lifetime of so-called cold quasiparticles if one assumes that the lifetime of these quasiparticles is also dominated by scattering off spin fluctuations .
Our paper is organized as follows. In Sec. II we give a brief overview of the theoretical framework in which we analyze the INS and NMR data. In Sec. III A we analyze NMR data on YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> for both commensurate and incommensurate magnetic response and discuss the role of lattice corrections. In Secs. III B and III C we analyze NMR data on two YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> compounds. In Sec. IV we discuss the consistency between NMR and INS data, as well as the role of lattice corrections for the local susceptibility. In Sec. V we propose a spin and charge inhomogeneity scenario as a possible way to reconcile NMR and INS data. Finally, in Sec. VI we draw our conclusions.
## II Theoretical Overview
We briefly discuss the theoretical framework in which we discuss NMR and INS experiments. In order to analyze the low frequency NMR data, we use the dynamical spin susceptibility of Eq.(1), where for the commensurate case $`𝐐=(\pi ,\pi )`$. To allow for an incommensurate structure of $`\chi `$, we use Eq.(1) with $`𝐐=(1,1\pm \delta )\pi `$ and $`𝐐=(1\pm \delta ,1)\pi `$ and sum over all four peaks . The inclusion of the correct form of lattice corrections is rather difficult since it requires a microscopic model which is beyond the scope of this paper. In general we expect lattice corrections to appear in the form of an upper momentum cutoff, $`\mathrm{\Lambda }`$. In the following we choose a soft cutoff procedure for $`\delta 𝐪𝐪𝐐`$ in the denominator of Eq.(1):
$$\delta 𝐪^2\delta 𝐪^2\left(1+\frac{\delta 𝐪^2}{\mathrm{\Lambda }^2}\right).$$
(3)
We also expect a weak momentum dependence of $`\alpha `$ and $`\omega _{\mathrm{sf}}`$ once the continuum description breaks down; however, to keep the number of tunable parameters small we ignore these effects.
In NMR experiments, one measures the spin-lattice relaxation rate $`1/T_{1x}`$, with applied magnetic field in $`x`$-direction, and the spin-echo rate $`1/T_{2G}`$, which can be expressed in terms of the dynamical susceptibility as:
$`{\displaystyle \frac{1}{T_{1x}T}}`$ $`=`$ $`{\displaystyle \frac{k_B}{2\mathrm{}}}(\mathrm{}^2\gamma _n\gamma _e)^2{\displaystyle \frac{1}{N}}{\displaystyle \underset{q}{}}F_x(q)\underset{\omega 0}{lim}{\displaystyle \frac{\chi ^{\prime \prime }(q,\omega )}{\omega }},`$ (4)
$`\left({\displaystyle \frac{1}{T_{2G}}}\right)^2`$ $`=`$ $`{\displaystyle \frac{0.69}{128\mathrm{}^2}}(\mathrm{}^2\gamma _n\gamma _e)^4\{{\displaystyle \frac{1}{N}}{\displaystyle \underset{q}{}}\left[F_{ab}^{eff}(q)\chi ^{}(q,\omega )\right]^2`$ (6)
$`\left[{\displaystyle \underset{q}{}}F_{ab}^{eff}(q)\chi ^{}(q,\omega )\right]^2\},`$
where $`x=ab,c`$ describes the direction of the external magnetic field. The $`{}_{}{}^{63}Cu`$ form factors are given by
$`{}_{}{}^{63}F_{c}^{}(q)`$ $`=`$ $`\left[A_{ab}+2B\left(cos(q_x)+cos(q_y)\right)\right]^2,`$ (7)
$`{}_{}{}^{63}F_{ab}^{eff}(q)`$ $`=`$ $`\left[A_c+2B\left(cos(q_x)+cos(q_y)\right)\right]^2,`$ (8)
$`{}_{}{}^{63}F_{ab}^{}(q)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[^{63}F_{ab}^{eff}(q)+^{63}F_c(q)].`$ (9)
where $`A_{ab},A_c`$ and $`B`$ are the on-site and transferred hyperfine coupling constants, respectively .
It follows from Eq.(1) that the spin excitations in the normal state are completely described by three parameters, $`\alpha ,\xi `$, and $`\omega _{sf}`$. In order to extract these parameters from the experimental NMR data, we also need to obtain the three hyperfine coupling constants, $`A_{ab},A_c`$, and $`B`$; hence we have six unknown parameters in the above equations, and require six equations to determine them. So far we have two, Eqs.(4) and (6). An additional constraint arises from the temperature independence of the <sup>63</sup>Cu Knight shift in a magnetic field parallel to the $`c`$ axis in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> which yields
$$A_c+4B0.$$
(10)
A fourth constraint comes from the anisotropy of the <sup>63</sup>Cu spin-lattice relaxation rates , which for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> was measured to be
$${}_{}{}^{63}R=\frac{T_{1c}}{T_{1ab}}3.7\pm 0.1.$$
(11)
A fifth constraint involving the hyperfine coupling constants can be obtained by plotting the Knight shift $`{}_{}{}^{63}K_{ab}^{}`$ versus $`\chi _0(T)`$ , which yields
$$4B+A_{ab}200\frac{kOe}{\mu _B}.$$
(12)
A final constraint is obtained from the earlier conjecture that the antiferromagnetic correlation length at the crossover temperature $`T_{\mathrm{cr}}`$ is approximately two lattice constants . Recent microscopic calculations have confirmed this conjecture, showing that indeed $`\xi (T_{cr})`$ is of the order of a few lattice constants.
To the extent that experimental data for $`T_1`$ and $`T_{2G}`$ are available up to $`T_{cr}`$, one can fit the above set of six equations self-consistently to the data, and thus obtain not only the temperature independent parameters $`\alpha ,A_{ab},A_c`$, and $`B`$, but also the temperature dependence of $`\xi `$ and $`\omega _{sf}`$. It turns out, however, that this analysis is only possible for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>, since this is the single compound for which data up to sufficiently high temperatures have been obtained . Our corresponding results will be presented in section III A. In order to extract the relevant parameters from NMR data on the related, widely studied, YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub> compounds, and thus to make a comparison between NMR and INS experiments possible, we need to make two assumptions, which, at least partly, will be supported by the experimental data for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. First, we assume a relation between $`\omega _{sf}`$ and $`\xi `$, given by
$$\omega _{sf}=\widehat{c}\xi ^z,$$
(13)
where $`\widehat{c}`$ is a temperature independent constant. For YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>, where we can independently extract $`\omega _{sf}`$ and $`\xi `$, we will show that there is indeed a crossover from $`z=1`$ behavior at low temperature to $`z=2`$ behavior at higher temperatures. The second assumption concerns the temperature dependence of the magnetic correlation length, $`\xi (T)`$, which we take to be that obtained by one of us using a renormalization group (RG) approach : for $`z=1`$ the temperature dependence of the magnetic correlation length is given by
$$\xi ^2=\xi _0^2+bT^2,$$
(14)
a result which is in good agreement with INS experiments on La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> . Having all the necessary theoretical tools in place, we now address the questions raised in the introduction.
## III Analysis of NMR experiments
### A YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>
YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> is a benchmark system for our analysis, since Curro et al. have measured the relaxation times $`T_1`$ and $`T_{2G}`$ over a wide temperature range from 80K to 750K, and in particular between $`T_{}approx200`$ K and $`T_{cr}approx500`$ K. We are therefore able to extract the relevant parameters which we discussed above from a self-consistent fit of Eqs.(4)- (12) to the $`T_1`$ and $`T_{2G}`$ data. In doing so, we assume that the constraints given by Eqs.(10)-(12) for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> are also valid for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. Since we can determine $`\xi `$ and $`\omega _{sf}`$ independently, we are able to study the effect of both an incommensurate magnetic response, and of non-universal lattice corrections, on the scaling law $`\omega _{sf}=\widehat{c}/\xi ^z`$.
We first reconsider the case of a commensurate structure of $`\chi `$, and set $`\mathrm{\Lambda }=2\sqrt{\pi }`$, i.e. we choose a momentum cut-off which corresponds to the linear size of the BZ. We present the temperature dependence of both $`\omega _{sf}`$ and $`\xi ^1`$, which results from a self-consistent fit to the experimental data, in Fig. 1.
Between $`T_{}=200`$ K and $`T_{cr}=500`$ K, $`\omega _{sf}`$ and $`\xi ^1`$ clearly scale linearly with temperature. In order to study the extent to which $`\omega _{sf}`$ and $`\xi `$ obey the above scaling law, we plot in Fig. 2, $`\mathrm{ln}(\omega _{sf})`$ as a function of $`\mathrm{ln}(\xi )`$ (lower curve).
The dynamical scaling range is of course too limited to prove the existence of a scaling relation. However, assuming $`\omega _{\mathrm{sf}}\xi ^z`$, we find $`z1`$ below $`500\mathrm{K}`$, which we identify with $`T_{\mathrm{cr}}`$ and $`z2`$ above $`T_{\mathrm{cr}}`$. These results are consistent with an earlier analysis of the scaling behavior and the microscopic scenario of Ref. . To study the effect of stronger lattice corrections, we decreased $`\mathrm{\Lambda }`$ to $`\mathrm{\Lambda }=2\sqrt{\pi }/4`$ (upper curve in Fig. 2 which for clarity is offset). Our conclusions concerning the scaling relations remain unchanged. We can thus conclude that the dynamical scaling behavior of spin excitations is robust against even sizeable lattice corrections, as long as the structure of $`\chi ^{\prime \prime }`$ is commensurate.
We next examine the effect of incommensuration on the dynamical scaling. It was earlier argued that the magnetic response in YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> should be very similar to YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.8</sub> . Using the observed doping dependence of the incommensuration in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6+x</sub>, we estimate the incommensuration in YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> to be $`\delta =0.23`$. Setting $`\mathrm{\Lambda }=2\sqrt{\pi }`$ we plot in Fig. 3 $`\mathrm{ln}(\omega _{sf})`$ as a function of $`\mathrm{ln}(\xi )`$ (lower curve) and find, following the same argumentation as above, that the dynamical scaling exponent below $`T_{\mathrm{cr}}500`$ K is now $`z0.75`$. Increasing the strength of the lattice corrections by decreasing $`\mathrm{\Lambda }`$ we find that $`z`$ is increased. However, in order to obtain again $`z1`$ (the upper curve in Fig. 3), which one would expect from the INS data on La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, we need a very small momentum cut-off, $`\mathrm{\Lambda }2\sqrt{\pi }/15`$. Such a small cut-off of order $`O(1/\xi )`$ is inconsistent with the continuum theory which is the basis for a dynamical scaling approach. We therefore conclude that a spatially homogeneous incommensurate magnetic response is in contradiction with the available NMR data. The physical origin of this sensitivity of the magnetic response of a homogeneously incommensurate system, compared to a locally commensurate one, results from the fact that the former effectively decreases the spatial extent of the spin-spin correlations, thus increasing the role of large values of $`\delta 𝐪\mathrm{\Lambda }`$. It should also be noted that our arguments using lattice corrections to discriminate between these two scenarios works only for intermediate values of the correlation length. For very large $`\xi `$ the system should be insensitive to the cut-off procedure regardless of whether it is commensurate or incommensurate.
In what follows we assume that the low frequency magnetic response measured in an NMR experiment is indeed locally commensurate which implies that $`z=1`$ scaling prevails between $`T_{}`$ and $`T_{cr}`$. In Sec. V we present a possible theoretical scenario to resolve this apparent contradiction between the incommensuration seen in INS experiments and our conclusions based on the available NMR data. Finally, using the commensurate form of Eq.(1), and a momentum cut-off $`\mathrm{\Lambda }=2\sqrt{\pi }`$ we find the following parameters from the solution of the above self-consistent equations: $`\alpha =10eV^1`$, $`A_{ab}=20.2kOe/\mu _B`$, $`A_c=182.8kOe/\mu _B`$, and $`B=45.7kOe/\mu _B`$ for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>.
### B YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>
Since it was earlier argued that for the optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub>, $`T_{cr}125`$ K, only slightly above $`T_c`$, there is no significant temperature dependence of the relaxation rates between $`T_c`$ and $`T_{cr}`$. We will therefore perform our analysis of the NMR data only at $`T_{cr}`$. Assuming that $`\xi (T_{cr})=2`$, and setting $`\mathrm{\Lambda }=2\sqrt{\pi }`$, we utilize $`T_{2G}`$ data by Itoh et al. and Stern et al. , as well as $`T_{1c}`$ data by Hammel et al. to find the following hyperfine coupling constants $`A_{ab}=27.8kOe/\mu _B`$, $`A_c=175.2kOe/\mu _B`$, and $`B=43.9kOe/\mu _B`$. Note that these values agree well with the hyperfine coupling constants we extracted for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. With $`1/T_{2G}=10^4`$ s<sup>-1</sup> and $`T_{1c}T0.15`$ Ks at $`T_{cr}`$, we obtain $`\alpha =18.5`$ eV<sup>-1</sup>, $`\omega _{sf}19.8`$ meV and $`\widehat{c}39.6`$ meV. It turns out that the above parameter set is rather robust against changes in $`\mathrm{\Lambda }`$. When $`\mathrm{\Lambda }`$ is decreased from $`\mathrm{\Lambda }=2\sqrt{\pi }`$ to $`\mathrm{\Lambda }=2\sqrt{\pi }/4`$, $`A_c`$ and $`B`$ decrease by about 1.5%, $`A_{ab}`$ increases by about 8.5%, $`\alpha `$ increases by 11%, and $`\omega _{sf}`$ decreases by about 15%. A change of the momentum cut-off by a factor of 4, which decreases the area of integration in the BZ by a factor of 16, thus leads only to moderate changes in the parameter set. However, as we show in Sec. IV, the corresponding changes in the local susceptibility at high frequencies are much more dramatic.
### C YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub>
In what follows we show that, though only a limited set of NMR data on YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub> is available, we nevertheless reach the same conclusions regarding the scaling behavior as in YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. In particular, since no NMR data on YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub> are available above $`T=300`$ K, and since we do not know the exact value of $`T_{cr}`$, we cannot extract the parameter set from a self-consistent fit to the NMR data as we did in the case of YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. However, it is in this doping range that INS data are available and therefore a consistency check between the parameter sets extracted from NMR and INS measurements might be the most promising. It turns out that though we are not able to perform a fully self-consistent fit, we can still extract the parameter set if we make several additional assumptions which are supported by our results on YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. First, we assume that $`z=1`$ scaling is present between $`T_{}`$ and $`T_{cr}`$ and that we can describe the temperature dependence of $`\xi `$ and the relation between $`\omega _{sf}`$ and $`\xi `$ by Eqs.(13) and (14), respectively. Second, we assume that the hyperfine coupling constants are only weakly doping dependent, and that to good approximation we can use the same constants we extracted from the analysis of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub>. Third, we need an estimate for $`T_{cr}`$, which is unknown for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub>. However, since $`T_{cr}500`$ K for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>, and since $`T_{cr}`$ increases with decreasing doping, we assume $`T_{cr}550650`$ K for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub>. To demonstrate the effect of the uncertainty in the latter assumption, we calculate he parameters $`\alpha ,\omega _{sf}(T)`$ and $`\xi (T)`$ for $`T_{cr}=650`$ K and $`T_{cr}=550`$ K as well as for two different values of the momentum cut-off $`\mathrm{\Lambda }`$. The resulting parameter sets based on experimental data by Takigawa et al. are shown in Table I.
Note that the parameters in Table I are rather robust against large variations in $`T_{cr}`$ and/or $`\mathrm{\Lambda }`$ and only differ at the most by about 30%. In Fig. 4 we compare our theoretical results for the relaxation rates, $`1/T_1`$ and $`1/T_{2G}`$, with the experimental data by Takigawa et al. .
We find that the temperature dependence of the relaxation rates calculated for each of the four different parameter sets in Table I, is practically indistinguishable; we therefore only plot the results for the parameter set with $`T_{cr}=650`$ K and $`\mathrm{\Lambda }=2\sqrt{\pi }`$. This result is consistent with the robustness of the scaling behavior in the commensurate case against non-universal lattice corrections as found for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>. We find a consistent description of the experimental data for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub> using the assumption of a $`z=1`$ scaling relationship and local commensuration in the temperature regime $`T_{}230\mathrm{K}<T<300\mathrm{K}`$, and, as was the case for YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub>, assuming commensurate behavior leads to results which are practically independent of non-universal lattice corrections.
## IV Inelastic Neutron Scattering
In this section we address the following two questions:(i) What is the effect of lattice corrections on the local susceptibility, and (ii) can one describe the available INS data with the parameter set extracted in the previous section ?
As we discussed in the introduction we expect lattice corrections to lead to substantial corrections of $`\chi _{loc}^{\prime \prime }`$ at frequencies larger than $`\omega _{sf}`$. Above this energy scale we also expect that the typical energy, $`\mathrm{\Delta }_{\mathrm{sw}}`$, of propagating spin waves, which, at low frequencies are completely overdamped by particle hole excitations, comes into play, leading to a modified form of the spin susceptibility
$$\chi (𝐪,\omega )=\frac{\alpha \xi ^2}{1+\xi ^2(𝐪𝐐)^2i\omega /\omega _{sf}(\omega /\mathrm{\Delta }_{\mathrm{sw}})^2}.$$
(15)
For $`\omega <\mathrm{\Delta }_{\mathrm{sw}}^2/\omega _{\mathrm{sf}}`$ no sign of a propagating peak exists due to the diffusive character of the spin excitations. For $`\omega >\mathrm{\Delta }_{\mathrm{sw}}^2/\omega _{\mathrm{sf}}`$ and $`𝐪=𝐐`$ the consequence of a propagating mode is a pronounced pole in the excitation spectrum if $`\mathrm{\Delta }_{\mathrm{sw}}<\omega _{\mathrm{sf}}`$ and a soft upper cut off in energy if $`\mathrm{\Delta }_{\mathrm{sw}}>\omega _{\mathrm{sf}}`$. The form of $`\chi (𝐪,\omega )`$ in Eq.(15) can be shown to describe the propagating spin mode in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.5</sub> above $`T_c`$ .
Using the parameter sets (1) and (2) in Table I, extracted in the previous section from NMR experiments, we plot in Fig. 5 the local susceptibility of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub> for two different values of $`\mathrm{\Lambda }`$.
We clearly see that upon decreasing the momentum cut-off, the maximum in $`\chi _{loc}^{\prime \prime }`$ moves towards lower frequencies. Experimentally, however, the intensity at higher frequencies drops much faster than in Fig. 5, even for $`\mathrm{\Lambda }=2\sqrt{\pi }/4`$ . In order to explain this discrepancy, we note that at higher frequencies, the dominant contribution to $`\chi _{loc}^{\prime \prime }`$ comes from regions in momentum space which are far away from the peak position at $`(\pi ,\pi )`$. In these regions, $`\chi ^{\prime \prime }(𝐪,\omega )`$ is only weakly momentum dependent in which case its contribution might be easily attributed to the experimental background. In other words, we believe that the origin of the discrepancy lies in an underestimate of the experimental intensity at higher frequencies due to the problems one has in resolving it from the background .
In Fig. 6 we plot the calculated INS intensity, i.e., $`\chi ^{\prime \prime }(𝐐,\omega )`$ as a function of frequency at $`𝐐=(\pi ,\pi )`$. Since the calculated intensity is practically the same for all four parameter sets in Table I, we only present $`\chi ^{\prime \prime }(𝐐,\omega )`$ calculated with the second parameter set in Table I (solid line). In the same figure, we also include the experimental data by Fong et al. for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.7</sub> . For the temperature range of interest, this material is the closest match for YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.63</sub>. Calculating the overall INS intensity with the parameters extracted from NMR experiments in the previous section, we find that we underestimate the experimental INS intensity by about 56 %, or a factor of 2.3 (the dashed line shows the calculated INS intensity, multiplied by an overall factor of 2.3). Given the uncertainties in the experimental determination of the absolute scale of $`\chi ^{\prime \prime }`$, we believe that the above result represents reasonable agreement between the INS and NMR data and thus demonstrates consistency in the description of spin excitations based on these two experimental techniques.
## V Spin and Charge Inhomogeneities
We saw in Sec. III A that under the assumption of $`z=1`$ scaling, supported by recent INS experiments , a locally incommensurate magnetic response is inconsistent with the available experimental NMR data. We therefore concluded that NMR data support a locally commensurate magnetic structure. The question thus arises of how one can understand a locally commensurate, but globally incommensurate magnetic response as seen by INS? It has been suggested that charge stripes separated by an average distance $`l_0=2\pi /\delta `$ are the origin of the incommensurate magnetic response seen in INS experiments . For a large part of the sample this would leave the locally commensurate structure intact. Except for the Nd-doped 214 compounds and probably for the particular doping value $`x=\frac{1}{8}`$ in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>2</sub>, these stripes are believed to be dynamic rather than static in nature.
Considering a spatial charge variation:
$$\rho (r)=\rho _0+\delta \rho (r),$$
(16)
where $`\rho _0`$ is independent of $`r`$. The question arises whether the system is in a regime with $`\rho _0\delta \rho (r)`$, as suggested in Ref. , or rather characterized by rather small charge variations, i.e. $`\delta \rho (r)<\rho _0`$. The particular appeal of the former assumption of strong charge modulations is an immediate explanation of the doping dependence of the position of the incommensurate peak, $`\delta =2x`$. Stripes with one hole every second lattice site in a hole free antiferromagnetic environment move closer together upon doping, leading to the above $`x`$-dependence of $`\delta `$. Also, the exceptional behavior of various magnetic and transport properties for the doping value $`x=\frac{1}{8}`$ can be easily understood in terms of a stable commensurate arrangement of such stripes and the underlying lattice. Another argument in favor of this scenario are the recent results by Vojta and Sachdev who investigated the mean field behavior of a system with competing magnetic and long range Coulomb interactions and found areas with a strong charge modulation separated by regions with barely disturbed antiferromagnetic correlations.
Nevertheless, there are several conceptional problems with such a strong charge modulation, which lead us to present some arguments favoring a more moderate charge variation, $`\delta \rho (r)<\rho _0`$, under circumstances that the transverse mobility of the charge carriers with respect to the averaged position of the stripes is large. First, uncorrelated spatial stripes of width $`r_0`$, which fluctuate transversely over a distance $`d`$, will give rise to a broadening $`d/l_0^2`$ of the magnetic peaks observed by INS. The experimental INS results, however, provide strong support for a scenario in which the width of the magnetic peaks is determined by $`\xi ^1`$ . Thus, in order to observe separated incommensurate peaks, we need at least $`\xi d`$ to establish a well defined antiphase domain wall, implying $`dl_0`$. In other words, assuming uncorrelated stripes, the width $`d`$ over which stripes fluctuate must be very small to account for the existing INS data. Such ”stiff” but uncorrelated stripes seem to be in contradiction to the notion of stripes as a dynamical entity. Therefore, we arrive at the conclusion that spatial stripe fluctuations, if they exist, must be strongly correlated.
Second, at the characteristic temperatures where stripes with $`\rho _0\delta \rho (r)`$ appear, strong modifications of the resistivity and other transport properties are expected to occur. None or only moderate changes of this kind have been observed and no indication for a dimensional crossover from quasi one-dimensional to quasi two-dimensional dynamics seems to be present.
Third, for low frequencies, the spin excitations in doped cuprates are overdamped rather than propagating spin waves. The suppression of the spin damping upon opening of a quasiparticle gap in the superconducting as well as in the pseudogap state implies that the dominant source of the spin damping, $`\gamma `$, are particle-hole excitations. Assuming weak charge modulations, we then find $`\mathrm{Im}\chi ^1(\omega )\gamma \omega `$, in agreement with the results of INS experiments .
On the other hand, rigid stripes separated by one dimensional charge carriers, which are effectively bosonic in character, lead to strong deviations from the above frequency dependence of $`\mathrm{Im}\chi ^1(\omega )`$, in disagreement with INS experiments. Thus, $`\mathrm{Im}\chi ^1(\omega )`$ does not arise from spatially varying stripes even though they certainly affect the spin degrees of freedom and can cause the decay of these excitations.
Finally, NMR and other magnetic measurements on various cuprates which provide strong support for spatial inhomogeneities suggest mostly spatially varying spin degrees of freedom but not necessarily a strong modulation of the charge background.
The above points suggest that the inhomogeneities observed in the high-$`T_c`$ cuprates might be predominantly magnetic in origin with only moderate modulations of the charge density. This is not unrealistic since in strongly correlated systems small variations in the charge density can bring about substantial changes in the magnetic properties of the system. This scenario of a predominantly magnetic character of the inhomogeneities, leading to magnetic stripes and clusters, is schematically presented in Fig. 7.
Note, that the magnetic stripes, in which the electronic charge density is lower than that in the magnetic clusters, represent magnetic domain walls with a phase slip of $`\pi `$ in the ordering of the spins. Theoretically, we believe that this more homogeneous charge arrangement is a result of quantum fluctuations beyond the mean field level investigated in Ref. .
The formation of magnetic clusters and stripes requires competing interactions on different length scales. Besides the local Coulomb interaction, we need to identify a longer length scale interaction. It was argued earlier that in the cuprate superconductors an increase of the magnetic correlation length leads to strong vertex corrections, which in turn give rise to a gradient coupling between the fermionic and spin degrees of freedom . For a system with an ordered ground state, this takes the form of a dipole-dipole coupling which has been shown to cause inhomogeneities and domain formation in two dimensional magnetic systems. The length scale of this gradient coupling should roughly speaking be the magnetic correlation length, $`\xi `$. We thus need a minimum value of $`\xi `$ to (a) generate the gradient coupling, and to (b) observe the incommensuration. The temperature at which incommensuration should appear is roughly set by $`\xi (T)l_0=4`$, where the last equality arises from the momentum position of the incommensurate peaks. Since, as we argued earlier, $`\xi (T_{cr})2`$, $`T_{cr}`$ thus presents an upper bound for the formation of magnetic clusters. This anomalous coupling is of course enhanced once magnetic clusters are formed; the creation of magnetic inhomogeneities is a self-consistent process. The magnetic inhomogeneity scenario presented here is admittedly very qualitative and further microscopic investigations will be required to verify or disprove it.
## VI Conclusions
In this communication, we investigated whether it is possible to understand theoretically INS and NMR data within a single theoretical framework. Based on the observation of an incommensurate structure of the magnetic response by INS experiments , we considered whether this reflects a locally or globally incommensurate ordering. On analyzing NMR data on YBa<sub>2</sub>Cu<sub>4</sub>O<sub>8</sub> we found with the condition of $`z=1`$ scaling that a homogeneous incommensuration is, within the framework given by Eq.(1), inconsistent with the available experimental data. We thus concluded that the local magnetic structure as probed by NMR is likely commensurate.
We then investigated the effect of lattice corrections on both the parameter sets extracted from NMR data and on the local susceptibility measured in INS experiments. We found that while the resulting corrections to the NMR parameters are only weak, the most pronounced effect of such corrections appears in the local susceptibility determined in INS measurements at frequencies above $`\omega _{sf}`$. Specifically, we found that decreasing the momentum cut-off leads to a suppression of $`\chi _{loc}^{\prime \prime }`$ at higher frequencies, thus improving the agreement with the experimental data. We argued that the remaining discrepancies can be explained by an experimental underestimate of the INS intensity at higher frequencies. Furthermore, we quantitatively compared INS and NMR data and found agreement within a factor of 2. Given the large uncertainties in resolving $`\chi ^{\prime \prime }`$ from the large background and in determining the absolute intensity of $`\chi ^{\prime \prime }`$, we believe that this result demonstrates that a consistent description of INS and NMR data can be obtained using the expression for $`\chi `$ given in Eq.(1).
Finally, we discussed a spin and charge inhomogeneity scenario to reconcile the local commensurations, as seen in NMR experiments, and the incommensurate peaks, seen by INS. This scenario, even though similar in spirit to earlier charge stripe pictures, is based on a moderate to weak modulation of the charge density, causing a pronounced inhomogeneity in the magnetic properties and leading to magnetically coupled clusters separated by weakly correlated stripes.
We would like to thank A.V. Chubukov, P. Dai, B. Keimer, T. Mason, H. Mook, R. Stern, C.P. Slichter and B. Stojkovic for valuable discussions. This work has been supported in part by the Science and Technology Center for Superconductivity through NSF-grant DMR91-20000 (D.K.M.), and by DOE at Los Alamos (D.P.). |
warning/0002/astro-ph0002334.html | ar5iv | text | # Distance-Redshift in Inhomogeneous FLRW
## 1 INTRODUCTION
As limits on the global cosmological parameters $`\mathrm{\Omega }_m`$ and $`\mathrm{\Lambda }`$ have been refined, Schmidt et al. (1998) and Perlmutter et al. (1999), the optical inadequacy of the standard distance-redshift relation ($`D`$-$`z`$) of FLRW has become more apparent. The problem was first recognized long ago by Zel’dovich (1964), Bertotti (1966), and Kantowski (1969) but the lack of relevant data limited its significance. Even though the average mass density parameter $`\mathrm{\Omega }_m`$ (along with $`H_0`$ and $`\mathrm{\Lambda }`$) determines the large scale dynamic behavior of the pressure free universe, knowledge of the actual mass inhomogeneity is necessary to accurately determine these parameters from most observations. Most observations determine $`\mathrm{\Omega }_m`$ and $`\mathrm{\Lambda }`$ by (indirectly) comparing theoretical $`D`$-$`z`$ curves to observed data. However, $`D`$-$`z`$ depends on more than the average mass density. It can depend significantly on details of how the mass is distributed, i.e., on how inhomogeneous the mass is on the scale of the widths of the observing beams. If some significant fraction ($`\rho _I/\rho _01`$) of the total mass density is in the form of inhomogeneities and is excluded from the lines of sight to the distant objects observed, a modified, i.e., a partially filled-beam $`D`$-$`z`$ is required.
The necessity of taking into account the effect of inhomogeneities on observations is relatively easy to understand. Homogeneous matter inside an observing beam of light gravitationally focuses the beam much differently than does an equal-mass clump of externally lensing matter. The simplest correction for this gravity-light effect requires the introduction of another parameter $`\nu ,0\nu 2,`$ which gives the fraction $`\rho _I/\rho _0=\nu (\nu +1)/6`$ of the mass density of the universe removed from the observing beams as inhomogeneities. Using $`\nu `$ rather than $`\rho _I/\rho _0`$ or some other parameter is dictated by the mathematics of special functions. A reduced mass density in an observing beam causes it to diverge relative to a standard FLRW beam. For an observed object in such a universe to have the standard FLRW angular size it would thus have to be moved to a smaller $`z`$; i.e., objects will appear less bright than in the standard FLRW universe. A reasonable application of this model to SNe Ia observations takes $`\rho _I`$ as the galactic contribution to the total mass density $`\rho _0`$ and the remaining contribution as a smooth intergalactic medium. Galaxies are easily excluded from SNe Ia foregrounds by selection (intended or not) and if galaxy mass roughly follows light, including their mass in $`\rho _I`$ is appropriate. In the partially filled-beam model where the additional parameter $`\nu 0`$ has been introduced, only lensing by mass clumps external to the beam has been neglected. To compare individual observations to $`D`$-$`z`$ of this model requires only an occasional lensing correction; however, comparison with the standard FLRW $`D`$-$`z`$ ($`\nu =0`$) model requires a defocusing correction for the partially empty-beam of every observation, as well as the occasional lensing correction. If only weak and transparent lensing occurs (to the $`z_{max}`$ being observed) the standard FLRW $`D`$-$`z`$ ($`\nu =0`$) should give the mean $`D`$-$`z`$ curve. Wang (1999) argues that by using flux-averaging the mean can be accurately obtained. Kantowski (1998a) and Kantowski (1998b) claims that determining cosmological parameters from data compared with the partially filled Hubble curves given here is likely to be easier. Beyond selection effects, unknown lensing probabilities can be highly non-Gaussian and should make the mean more difficult to observationally determine, i.e., should require more data if a given accuracy of the cosmic parameters is to be obtained, Bertotti (1966); Holz & Wald (1998); Holz (1998). The down side for partially filled-beam models is that you must select against lensing and must determine the additional parameter $`\nu `$.
In Sec. 2 we outline the procedure required to obtain $`D`$-$`z`$ for partially filled-beam FLRW observations and how the result simplifies for the three special cases of $`\nu =`$ 0, 1, and 2. In Sec. 3 we give the new results for these three special cases. Some concluding remarks are given in Sec. 4 and in the Appendix we discuss our Fortran implementation of these results.
## 2 The Luminosity Distance-redshift Relation
For models being discussed here (and for most cosmological models), angular or apparent size distance is related to luminosity distance by $`D_<(z)=D_{\mathrm{}}(z)/(1+z)^2`$. Hence we need to give only one or the other, and we have chosen to give luminosity distances. The $`D_{\mathrm{}}(z)`$ which accounts for a partially depleted mass density in the observing beam but neglects lensing by external masses is found by integrating the second order differential equation for the cross sectional area $`A(z)`$ of an observing beam from source ($`z=z_s`$) to observer ($`z=0`$), see Kantowski (1998a) for some history of this equation:
$`(1+z)^3\sqrt{1+\mathrm{\Omega }_mz+\mathrm{\Omega }_\mathrm{\Lambda }[(1+z)^21]}\times `$
$`{\displaystyle \frac{d}{dz}}(1+z)^3\sqrt{1+\mathrm{\Omega }_mz+\mathrm{\Omega }_\mathrm{\Lambda }[(1+z)^21]}{\displaystyle \frac{d}{dz}}\sqrt{A(z)}`$
$`+{\displaystyle \frac{(3+\nu )(2\nu )}{4}}\mathrm{\Omega }_m(1+z)^5\sqrt{A(z)}=0.`$ (1)
The required boundary conditions are:
$`\sqrt{A}|_s`$ $`=`$ $`0,`$
$`{\displaystyle \frac{d\sqrt{A|_s}}{dz}}`$ $`=`$ $`\sqrt{\delta \mathrm{\Omega }}{\displaystyle \frac{c}{H_s(1+z_s)}},`$ (2)
where $`\delta \mathrm{\Omega }`$ is the solid angle of the beam at the source and the FLRW value of the Hubble parameter at $`z_s`$ is related to the current value $`H_0`$ at $`z=0`$ by:
$$H_s=H_0(1+z_s)\sqrt{1+\mathrm{\Omega }_mz_s+\mathrm{\Omega }_\mathrm{\Lambda }[(1+z_s)^21]}.$$
(3)
The luminosity distance is then simply related to the area $`A|_0`$ of the beam at the observer by:
$$D_{\mathrm{}}^2\frac{A|_0}{\delta \mathrm{\Omega }}(1+z_s)^2.$$
(4)
Equation (1) can be put into the form of a Lamé equation and its solution has been given in terms of Heun functions in Kantowski (1998a). Solutions can also be given in terms of Lamé functions but neither Heun nor Lamé functions are currently available in standard computer libraries. Consequently, such expressions are not particularly useful for comparison with data, at this time. For the special case where $`\mathrm{\Lambda }=0`$ the Lamé functions reduce to associated Legendre functions and these expressions are useful. Other special cases also exist as is pointed out in Kantowski (1998a).
In the next section we give useful expressions for $`D_{\mathrm{}}`$ for three special cases where $`\mathrm{\Lambda }`$ is arbitrary but where the filling parameter $`\nu `$ is restricted to values 0, 1, and 2. For these three cases we can write $`D_{\mathrm{}}`$ as an elliptic integral and hence we can give $`D_{\mathrm{}}`$ in terms of the three fundamental incomplete Legendre elliptic integrals $`F(\varphi ,\mathrm{k}),E(\varphi ,\mathrm{k}),`$ and $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$. These functions are universally available and these new expressions significantly speed up the evaluation of $`D_{\mathrm{}}`$ (see the Appendix). Distance-redshift for $`\mathrm{\Omega }_0=1`$ can be given in terms of hypergeometric functions, see (21) and (53), or associated Legendre functions, see (22) and (54); however, we also give $`D_{\mathrm{}}`$ as more complicated expressions involving Legendre elliptic integrals, (23) and (55), because these expressions evaluate more rapidly using currently available Fortran routines.
It is not at all clear that the solution of (1) can be written as elliptic integrals for the special cases of $`\nu =`$ 0, 1 and 2. However, the steps required to arrive at this conclusion can be found in Whittaker & Watson (1927) under integral functions for Lamé and Matthew equations (see especially Sec. 19.53). The authors have carried out the conversion directly for all three cases; however, the $`\nu =0`$ and $`2`$ conversions can be reached by simpler means. The integral for $`\nu =0`$, the standard FLRW filled-beam case, is given in (5) and is well known. The $`\nu =2`$ (empty-beam) integral given in (46) is easy to obtain because the coefficient of $`\sqrt{A}`$ vanishes in (1). The first integral is trivial and the second is elliptic resulting in (46). For $`\nu =1`$, the 66% filled-beam model, the integral is given in (30); however, no simple way of getting this from (1) seems to exist.
In Sec. 3. we outline results for all big bang models in the first quadrant of the $`\mathrm{\Omega }_m`$– $`\mathrm{\Omega }_\mathrm{\Lambda }`$ plane (see Fig. 1), hoping to facilitate their usage. Luminosity distances for the three large open domains are given in subsections A, and for the boundaries of these domains in subsections B.
## 3 Luminosity Distances as Legendre Elliptic Integrals
I. $`\nu =0`$, Completely Filled-Beam Observations (Standard FLRW)
A. Three Open Big Bang Domains
Kaufman & Schucking (1971) and Kaufman (1971) gave magnitude-redshift relations for standard pressure-free FLRW models as inverse Weierstrass functions and more recently Feige (1992) gave comoving distances and light travel times for these models using Legendre elliptic integrals. In this section we give simpler and more useful results which are directly comparable with Edwards (1972) who used Jacobi elliptic functions. The well known and often used integral form for luminosity distance in standard FLRW is:
$$D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =0;z)=\frac{c}{H_0}\frac{1+z}{\sqrt{|1\mathrm{\Omega }_0|}}S_\kappa \left[\sqrt{|1\mathrm{\Omega }_0|}_0^z\frac{dz}{\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}\right]$$
(5)
which we integrate using Byrd & Friedman (1971) to obtain,
$$D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =0;z)=\frac{c}{H_0}\frac{1+z}{\sqrt{|1\mathrm{\Omega }_0|}}S_\kappa [g\{F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})\}],$$
(6)
or equivalently using an addition formula for $`F(\varphi ,\mathrm{k})`$, i.e., $`F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})`$ = $`F(\mathrm{\Delta }\varphi _z,\mathrm{k})`$ we get:
$$D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =0;z)=\frac{c}{H_0}\frac{1+z}{\sqrt{|1\mathrm{\Omega }_0|}}S_\kappa [gF(\mathrm{\Delta }\varphi _z,\mathrm{k})].$$
(7)
The parameter $`\kappa `$ ($`\mathrm{\Omega }_01`$)/$`|\mathrm{\Omega }_01|`$ is determined by the sign of the 3-curvature and $`S_\kappa []`$ is one of two functions:
$$S_\kappa []=\{\begin{array}{ccc}\mathrm{sinh}[]\hfill & :& \kappa =1,\hfill \\ \mathrm{sin}[]\hfill & :& \kappa =+1.\hfill \end{array}$$
Constants $`g`$ and $`\mathrm{k}`$ depend on the cosmic parameters $`\mathrm{\Omega }_m\&\mathrm{\Omega }_\mathrm{\Lambda }`$, and $`F(\varphi ,\mathrm{k})`$ is the incomplete Legendre elliptic integral of the first kind.<sup>1</sup><sup>1</sup>1 $`F(\varphi ,\mathrm{k})_0^\varphi 1/\sqrt{1\mathrm{k}^2\mathrm{sin}^2\varphi }𝑑\varphi `$ The constants $`g`$ and $`\mathrm{k}`$ depend on $`\mathrm{\Omega }_m\&\mathrm{\Omega }_\mathrm{\Lambda }`$ only through a combination called $`b`$ defined by:
$$b(27/2)\frac{\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }}{(1\mathrm{\Omega }_0)^3},\mathrm{}b\mathrm{},$$
(8)
$`b<0\kappa =1,`$
$`b>0\kappa =+1.`$
The functions $`\varphi _z`$ and $`\mathrm{\Delta }\varphi _z`$ depend on the redshift $`z`$ and the cosmic parameters $`\mathrm{\Omega }_m\&\mathrm{\Omega }_\mathrm{\Lambda }`$ (not just on the combination $`b`$). Domains for the various $`b`$ values in the $`\mathrm{\Omega }_m`$– $`\mathrm{\Omega }_\mathrm{\Lambda }`$ plane are shown in Fig. 1.
1. For the two open domains defined by $`b<0`$ and $`2<b`$, quantities $`g,\mathrm{k},`$ $`\varphi _z`$, and $`\mathrm{\Delta }\varphi _z`$ are conveniently written in terms of intermediate constants $`v_\kappa ,y_1`$ and $`A`$ defined by:
$$v_\kappa \left[\kappa (b1)+\sqrt{b(b2)}\right]^{1/3},v_\kappa 1.$$
(9)
$$y_1\frac{1+\kappa (v_\kappa +v_\kappa ^1)}{3},$$
(10)
$$A=A(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })\sqrt{y_1(3y_1+2)}=\sqrt{\frac{v_\kappa ^2+v_\kappa ^2+1}{3}}1.$$
(11)
Parameters $`g`$ and $`\mathrm{k}`$ are then given by:
$$g=g(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })=1/\sqrt{A(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })}=\left[\frac{3}{v_\kappa ^2+v_\kappa ^2+1}\right]^{1/4}1,$$
(12)
and
$$\mathrm{k}^2=\mathrm{k}^2(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })=\frac{2A+\kappa (1+3y_1)}{4A}=\left[\frac{1}{2}+\frac{1}{4}g^2(v_\kappa +v_\kappa ^1)\right]1.$$
(13)
Functions $`\varphi _z`$, and $`\mathrm{\Delta }\varphi _z`$ are given by:
$$\varphi _z=\varphi (\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda };z)=\mathrm{cos}^1\left[\frac{(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+\kappa y_1A}{(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+\kappa y_1+A}\right],$$
(14)
and
$$\mathrm{\Delta }\varphi _z=\mathrm{\Delta }\varphi (\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda };z)=2\mathrm{tan}^1\left[\frac{z\sqrt{A}\sqrt{|1\mathrm{\Omega }_0|}\sqrt{1+z[1(1\mathrm{\Omega }_0)\mathrm{\Omega }_m^1y_1]^1}}{1+z[1(1\mathrm{\Omega }_0)\mathrm{\Omega }_m^1y_1]^1+\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}\right].$$
(15)
2. For the domain $`0<b<2`$ ($`\kappa =1`$) three intermediate parameter $`y_1,y_2`$ and $`y_3`$ are convenient to use, although none are really necessary. In this domain of $`b`$, intermediate parameters $`y_1,y_2`$ and $`y_2`$ are related to the cosmic parameters $`\mathrm{\Omega }_m\&\mathrm{\Omega }_\mathrm{\Lambda }`$ through $`b`$ by:
$`y_1`$ $``$ $`{\displaystyle \frac{1}{3}}\left(1+\mathrm{cos}\left[{\displaystyle \frac{\mathrm{cos}^1(1b)}{3}}\right]+\sqrt{3}\mathrm{sin}\left[{\displaystyle \frac{\mathrm{cos}^1(1b)}{3}}\right]\right),0y_11/3,`$
$`y_2`$ $``$ $`{\displaystyle \frac{1}{3}}\left(12\mathrm{cos}\left[{\displaystyle \frac{\mathrm{cos}^1(1b)}{3}}\right]\right),1y_22/3,`$
$`y_3`$ $``$ $`{\displaystyle \frac{1}{3}}\left(1+\mathrm{cos}\left[{\displaystyle \frac{\mathrm{cos}^1(1b)}{3}}\right]\sqrt{3}\mathrm{sin}\left[{\displaystyle \frac{\mathrm{cos}^1(1b)}{3}}\right]\right),2/3y_30.`$ (16)
The following expressions are valid only in the lower right part of the $`\mathrm{\Omega }_m`$– $`\mathrm{\Omega }_\mathrm{\Lambda }`$ plane. In the upper left domain where $`b`$ also satisfies $`0b2`$, expressions can be given, but there a big bang doesn’t occur. The parameters $`g`$ and $`\mathrm{k}`$ and functions $`\varphi _z`$ and $`\mathrm{\Delta }\varphi _z`$ needed to evaluate (6) and (7) are:
$$g=g(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })\frac{2}{\sqrt{y_1y_2}},$$
(17)
$$\mathrm{k}^2=\mathrm{k}^2(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda })\frac{y_1y_3}{y_1y_2}1,$$
(18)
$$\varphi _z=\varphi (\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda };z)=\mathrm{sin}^1\sqrt{\frac{y_1y_2}{(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+y_1}},$$
(19)
$`\mathrm{\Delta }\varphi _z`$ $`=`$ $`\mathrm{\Delta }\varphi (\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda };z)`$ (20)
$`=`$ $`2\mathrm{tan}^1\left[{\displaystyle \frac{\sqrt{y_1y_2}\left[\sqrt{y_3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}\sqrt{y_3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}\right]}{\sqrt{[y_1\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)][y_2(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]}+\sqrt{[y_1y_2]}}}\right],`$
where $`y_1y_2`$ means repeat the previous term with $`y_1`$ and $`y_2`$ exchanged.
B. Boundaries
1. $`\mathrm{\Omega }_0\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$
For the spatially flat model ($`b\pm \mathrm{})`$ a much simpler expression involving hypergeometric functions results:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =0;z)={\displaystyle \frac{c}{H_0}}(1+z){\displaystyle _0^z}{\displaystyle \frac{dz}{\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}}`$
$`={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{2(1+z)}{\mathrm{\Omega }_m^{1/3}}}\left[_2F_1\right({\displaystyle \frac{1}{6}},{\displaystyle \frac{2}{3}};{\displaystyle \frac{7}{6}};1\mathrm{\Omega }_m)`$
$`\left({\displaystyle \frac{1}{[1+\mathrm{\Omega }_mz(3+3z+z^2)]^{1/6}}}\right){}_{2}{}^{}F_{1}^{}({\displaystyle \frac{1}{6}},{\displaystyle \frac{2}{3}};{\displaystyle \frac{7}{6}};{\displaystyle \frac{1\mathrm{\Omega }_m}{1+\mathrm{\Omega }_mz(3+3z+z^2)}})].`$ (21)
When $`\mathrm{\Omega }_m1`$ (21) can be expressed as associated Legendre functions,
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =0;z)={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{2^{1/6}\mathrm{\Gamma }\left(1/6\right)(1+z)}{3[\mathrm{\Omega }_m^5(1\mathrm{\Omega }_m)]^{1/12}}}`$
$`\times \left[\mathrm{P}_{1/6}^{1/6}\left({\displaystyle \frac{1}{\sqrt{\mathrm{\Omega }_m}}}\right){\displaystyle \frac{1}{(1+z)^{(1/4)}}}\mathrm{P}_{1/6}^{1/6}\left(\sqrt{{\displaystyle \frac{1+\mathrm{\Omega }_mz(3+3z+z^2)}{\mathrm{\Omega }_m(1+z)^3}}}\right)\right].`$ (22)
If (22) is given in terms of Legendre elliptic integrals the result is more complicated:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =0;z)`$ $`={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{1+z}{(3)^{1/4}\sqrt{\mathrm{\Omega }_m}(\mathrm{\Omega }_m^11)^{1/6}}}\left[\{F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})\}\right],`$ (23)
$`={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{1+z}{(3)^{1/4}\sqrt{\mathrm{\Omega }_m}(\mathrm{\Omega }_m^11)^{1/6}}}\left[F(\mathrm{\Delta }\varphi _z,\mathrm{k})\right],`$
where
$$\mathrm{k}^2=\left[\frac{1}{2}+\frac{\sqrt{3}}{4}\right],$$
(24)
$$\varphi _z=\varphi (\mathrm{\Omega }_m;z)=\mathrm{cos}^1\left[\frac{1+z+(1\sqrt{3})(\mathrm{\Omega }_m^11)^{1/3}}{1+z+(1+\sqrt{3})(\mathrm{\Omega }_m^11)^{1/3}}\right],$$
(25)
and
$$\mathrm{\Delta }\varphi _z=\mathrm{\Delta }\varphi (\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda };z)2\mathrm{tan}^1\left[\frac{z\sqrt{\sqrt{3}\mathrm{\Omega }_m(1/\mathrm{\Omega }_m1)^{1/3}}\sqrt{1+z[1+(1/\mathrm{\Omega }_m1)^{1/3}]^1}}{1+z[1+(1/\mathrm{\Omega }_m1)^{1/3}]^1+\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}\right].$$
(26)
2. $`b=2`$
This value of $`b`$ can be identified with “critical” values of the cosmic parameters, Felten & Isaacman (1986). We give a result good only for the lower $`b=2`$ curve, see (44). These models start with a big bang and expand to the the finite Einstein radius at $`t=\mathrm{}`$, see A3(vii-b) in the appendix of McVittie (1965):
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m),\nu =0;z)={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{1+z}{\sqrt{|1\mathrm{\Omega }_0|}}}`$
$`\times \mathrm{sin}\left\{\mathrm{ln}\left({\displaystyle \frac{\left[\sqrt{1/3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}+1\right]\left[\sqrt{1/3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}1\right]}{\left[\sqrt{1/3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}1\right]\left[\sqrt{1/3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}+1\right]}}\right)\right\}.`$ (27)
3. $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$
This result is due to Mattig (1958), we include it for completeness:
$$D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=0,\nu =0;z)=\frac{2c}{H_0\mathrm{\Omega }_m^2}\{\mathrm{\Omega }_mz+(\mathrm{\Omega }_m2)(\sqrt{1+\mathrm{\Omega }_mz}1)\}.$$
(28)
4. $`\mathrm{\Omega }_m=0`$
These are massless big bang models, $`\mathrm{\Omega }_\mathrm{\Lambda }<1`$, discussed by Robertson (1933):
$$D_{\mathrm{}}(\mathrm{\Omega }_m=0,\mathrm{\Omega }_\mathrm{\Lambda },\nu =0;z)=\frac{c(1+z)}{H_0\mathrm{\Omega }_\mathrm{\Lambda }}\left\{1+z\sqrt{\mathrm{\Omega }_\mathrm{\Lambda }+(1+z)^2(1\mathrm{\Omega }_\mathrm{\Lambda })}\right\}.$$
(29)
II. $`\nu =1`$, 66% Filled-Beam Observations
A. Four Open Big Bang Domains
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =1;z)=`$
$`{\displaystyle \frac{c}{H_0}}2(1+z)\mathrm{Sign}[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]\sqrt{\left|{\displaystyle \frac{\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\left[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]}{(1\mathrm{\Omega }_0)[36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3]}}\right|}\times `$
$`S_{(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },z)}[\sqrt{|(1\mathrm{\Omega }_0)[36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3]|}\times `$
$`P{\displaystyle _0^z}{\displaystyle \frac{dz}{2\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}}],`$ (30)
where
$$S_{(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },z)}[]=\{\begin{array}{ccc}\mathrm{cosh}[]\hfill & :& b<0\&\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\left[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]<0,\hfill \\ \mathrm{sinh}[]\hfill & :& b<0\&\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\left[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]>0,\hfill \\ \mathrm{sin}[]\hfill & :& 0<b<486,\hfill \\ \mathrm{sinh}[]\hfill & :& 486<b.\hfill \end{array}$$
Only the principal value of the integral (P) is needed and unlike the $`\nu =0`$ case, this integral takes on different forms when evaluated using Legendre elliptic integrals, depending on the value of the parameter $`b`$. Parts of the analytic result (31) sometimes diverge even though the total expression remains finite. For example when $`b=486`$, i.e., when $`\sqrt{36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3}=0`$ or equivalently $`y_1=3`$, a limit must be taken. The resulting $`D_{\mathrm{}}`$ on this new boundary can be found in II.B.5 below. This new boundary splits the one open domain $`2<b<\mathrm{}`$ into two parts, see Fig. 2. Consequently, the $`\mathrm{\Omega }_m`$– $`\mathrm{\Omega }_\mathrm{\Lambda }`$ plane is more complicated for $`\nu =1`$ than for either $`\nu =0`$ or $`\nu =2`$. See A1 below for additional trouble points that occur.
1. For the three open domains defined by $`b<0,2<b<486`$, and $`486<b`$ the luminosity distance $`D_{\mathrm{}}`$ takes the form:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =1;z)=`$
$`{\displaystyle \frac{c}{H_0}}2(1+z)\mathrm{Sign}[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]\sqrt{\left|{\displaystyle \frac{\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\left[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]}{(1\mathrm{\Omega }_0)[36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3]}}\right|}\times `$
$`S_{(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },z)}[{\displaystyle \frac{\kappa \sqrt{|36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3|}}{2\sqrt{A}[A+\kappa (y_13)]}}\{[F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})]`$
$`+{\displaystyle \frac{A\kappa (y_13)}{2\kappa (y_13)}}[\mathrm{P}\mathrm{\Pi }(\varphi _z,\widehat{\alpha }^2,\mathrm{k})\mathrm{P}\mathrm{\Pi }(\varphi _0,\widehat{\alpha }^2,\mathrm{k})]\}+f_b],`$ (31)
where $`y_1,A,\mathrm{k},`$ and $`\varphi _z`$ are defined in (10)-(14) and the additional constant $`\widehat{\alpha }^2`$ is:
$$\widehat{\alpha }^2\frac{(A+\kappa (y_13))^2}{4A\kappa (y_13)}.$$
(32)
$`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$ is the incomplete Legendre elliptic integral of the third kind<sup>2</sup><sup>2</sup>2 $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})_0^\varphi 1/\left[(1\alpha ^2\mathrm{sin}^2\varphi )\sqrt{1\mathrm{k}^2\mathrm{sin}^2\varphi }\right]𝑑\varphi `$. In arriving at the results for the two-thirds filled beam model we discovered that equation 361.54 of Byrd & Friedman (1971) has the two square-root terms interchanged for the case $`\alpha ^2/(\alpha ^21)>\mathrm{k}^2`$. and P $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$ is the principal part of that integral. The function $`f_b`$ is one of,
$$f_b=\{\begin{array}{ccc}\frac{1}{4}\mathrm{ln}\left|\{[1+h(z)][1h(0)]\}/\{[1+h(0)][1h(z)]\}\right|\hfill & :& b<0\mathrm{or}486<b,\hfill \\ \frac{1}{2}\left[\mathrm{tan}^1h(z)\mathrm{tan}^1h(0)\right]\hfill & :& 2<b<486,\hfill \end{array}$$
where $`h(z)`$ is defined by:
$$h(z)\frac{\sqrt{|36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3|}\sqrt{(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+\kappa y_1}}{(3y_1)\sqrt{\left[(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|\kappa (1+y_1)/2\right]^2(1+y_1)(13y_1)/4}}.$$
(33)
Some care has to be taken when using these expressions. Divergences in the function $`f_b`$ necessarily occur and cancel divergences in $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$. Divergences in $`f_b`$ also occur which add to divergences in $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$ and cancel zeros in the multiplicative factor $`\sqrt{\left|\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\left[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]\right|}`$ of (31). Redshift independent divergences occur when $`\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)=3`$ and when $`\mathrm{\Omega }_m(3y_1)/(1\mathrm{\Omega }_0)=y_1(2y_1+5)`$. These points are plotted in Figure 2. Redshift dependent divergences occur at $`(1+z)=3(1\mathrm{\Omega }_0)/\mathrm{\Omega }_m`$ and at $`(1+z)\mathrm{\Omega }_m(3y_1)/(1\mathrm{\Omega }_0)=y_1(5+2y_1)`$. These points appear in the $`\mathrm{\Omega }_m`$– $`\mathrm{\Omega }_\mathrm{\Lambda }`$ plane respectively to the left of the $`\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)=3`$ line and between the $`\mathrm{\Omega }_m(3y_1)/(1\mathrm{\Omega }_0)=y_1(5+2y_1)`$ and $`b=486`$ curves.
Computer evaluation of (31) can be speeded up by reducing the number of Legendre elliptic integrals that must be evaluated. As in (7) we can use the addition formula for $`F(\varphi ,\mathrm{k})`$, i.e., $`F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})=F(\mathrm{\Delta }\varphi _z,\mathrm{k})`$ and an addition formula for $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$,<sup>3</sup><sup>3</sup>3This equation is 116.03 of Byrd & Friedman (1971), corrected for two sign errors.
$$\mathrm{\Pi }(\varphi _z,\alpha ^2,\mathrm{k})\mathrm{\Pi }(\varphi _0,\alpha ^2,\mathrm{k})=\mathrm{\Pi }(\mathrm{\Delta }\varphi _z,\alpha ^2,\mathrm{k})+\frac{1}{2}\sqrt{\frac{\alpha ^2}{(\alpha ^21)(\alpha ^2\mathrm{k}^2)}}\mathrm{log}\left(\frac{1+\xi }{1\xi }\right),$$
(34)
where
$$\xi \frac{\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{sin}\mathrm{\Delta }\varphi _z\sqrt{\alpha ^2(\alpha ^21)(\alpha ^2\mathrm{k}^2)}}{1\alpha ^2\mathrm{sin}^2\mathrm{\Delta }\varphi _z\alpha ^2\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{cos}\mathrm{\Delta }\varphi _z\sqrt{1\mathrm{k}^2\mathrm{sin}^2\mathrm{\Delta }\varphi _z}},$$
(35)
to cut the number of elliptic functions from four to two. We were not able to simplify this expression enough to justify inclusion of a rewritten version of (31). However, it was used in our Fortran implementation (see Appendix).
2. For the open domain defined by $`0<b<2`$ the luminosity distance $`D_{\mathrm{}}`$ has a somewhat simpler form:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =1;z)={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{2(1+z)}{\sqrt{|1\mathrm{\Omega }_0|}}}\sqrt{{\displaystyle \frac{\left[3\mathrm{\Omega }_m(1+z)/(1\mathrm{\Omega }_0)\right]\left[3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]}{\left[36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3\right]}}}\times `$
$`\mathrm{sin}[{\displaystyle \frac{\sqrt{36+\mathrm{\Omega }_m^2\mathrm{\Omega }_\mathrm{\Lambda }/(1\mathrm{\Omega }_0)^3}}{(3y_1)\sqrt{y_1y_2}}}\{[F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})]`$
$`+[\mathrm{\Pi }(\varphi _z,{\displaystyle \frac{y_13}{y_1y_2}},\mathrm{k})\mathrm{\Pi }(\varphi _0,{\displaystyle \frac{y_13}{y_1y_2}},\mathrm{k})]\}].`$ (36)
The constants $`y_1,y_2`$ and $`\mathrm{k}`$, and the function $`\varphi _z`$ are as defined in I.A.2 above \[see (16)-(19)\]. Just as in the previous case, the number of Legendre elliptic functions in (36) can be reduced from four to two by using the appropriate addition formulas. For $`F(\varphi ,\mathrm{k})`$ the formula is always the same, see (6) and (7), but because $`\alpha ^2`$ is negative (34) changes to:
$$\mathrm{\Pi }(\varphi _z,\alpha ^2,\mathrm{k})\mathrm{\Pi }(\varphi _0,\alpha ^2,\mathrm{k})=\mathrm{\Pi }(\mathrm{\Delta }\varphi _z,\alpha ^2,\mathrm{k})\frac{1}{2}\sqrt{\frac{\alpha ^2}{(1\alpha ^2)(\alpha ^2\mathrm{k}^2)}}\mathrm{tan}^1\left(\xi \right),$$
(37)
where
$$\xi \frac{\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{sin}\mathrm{\Delta }\varphi _z\sqrt{\alpha ^2(1\alpha ^2)(\alpha ^2\mathrm{k}^2)}}{1\alpha ^2\mathrm{sin}^2\mathrm{\Delta }\varphi _z\alpha ^2\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{cos}\mathrm{\Delta }\varphi _z\sqrt{1\mathrm{k}^2\mathrm{sin}^2\mathrm{\Delta }\varphi _z}}.$$
(38)
This is 116.02 of Byrd & Friedman (1971) with one sign error corrected.
B. Boundaries
1. $`\mathrm{\Omega }_0\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$
For these models $`b\pm \mathrm{}`$ and a much simpler expression results:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =1;z)`$
$`={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{2(1+z)^{3/2}}{\sqrt{1\mathrm{\Omega }_m}}}\mathrm{sinh}\left[{\displaystyle \frac{\sqrt{1\mathrm{\Omega }_m}}{2}}{\displaystyle _0^z}{\displaystyle \frac{dz}{(1+z)\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}}\right],`$
$`={\displaystyle \frac{c}{H_0\sqrt{1\mathrm{\Omega }_m}}}(1+z)^2[\left({\displaystyle \frac{1+\sqrt{1\mathrm{\Omega }_m}}{\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}+\sqrt{1\mathrm{\Omega }_m}}}\right)^{1/3}`$
$`\left({\displaystyle \frac{1\sqrt{1\mathrm{\Omega }_m}}{\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}\sqrt{1\mathrm{\Omega }_m}}}\right)^{1/3}].`$ (39)
This result can be given in terms of Legendre elliptic integrals $`F(\varphi ,\mathrm{k})`$ and $`\mathrm{\Pi }(\varphi ,\alpha ^2,\mathrm{k})`$; however, the authors can think of no useful purpose in doing so.
2. $`b=2`$
See the description for the $`\nu =0`$ case in section I.B.2 including (44) for this “critical” value of $`b`$:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m),\nu =1;z)={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{(1+z)\sqrt{3/2}}{\sqrt{|1\mathrm{\Omega }_0|}(11)}}\{`$
$`\sqrt{8}\left[\sqrt{13{\displaystyle \frac{\mathrm{\Omega }_m}{1\mathrm{\Omega }_0}}}\sqrt{13{\displaystyle \frac{\mathrm{\Omega }_m(1+z)}{1\mathrm{\Omega }_0}}}\right]\mathrm{cos}\left({\displaystyle \frac{4}{\sqrt{6}}}\mathrm{log}(h_z)\right)`$
$`+[8+\sqrt{\left(13{\displaystyle \frac{\mathrm{\Omega }_m}{1\mathrm{\Omega }_0}}\right)\left(13{\displaystyle \frac{(1+z)\mathrm{\Omega }_m}{1\mathrm{\Omega }_0}}\right)}]\mathrm{sin}\left({\displaystyle \frac{4}{\sqrt{6}}}\mathrm{log}(h_z)\right)\},`$ (40)
where $`h_z`$ is defined by:
$$h_z=\left(\frac{1+\sqrt{1/3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}}{1+\sqrt{1/3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}}\right)\sqrt{\frac{2/3+(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}{2/3+\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}}.$$
(41)
3. $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$
This result was first given by Dyer & Roeder (1973),
$$D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=0,\nu =1;z)=\frac{c}{H_0}\frac{4}{3\mathrm{\Omega }_m^2}[(\frac{3}{2}\mathrm{\Omega }_m1+\frac{1}{2}\mathrm{\Omega }_mz)\sqrt{1+\mathrm{\Omega }_mz}(\frac{3}{2}\mathrm{\Omega }_m1)].$$
(42)
4. $`\mathrm{\Omega }_m=0`$
This result is exactly the same as the $`\nu =0`$ result (29). If there is no mass in the universe then removing 33% of no mass from the beam changes nothing.
5. $`b=486`$
This result is equivalent to the $`b486`$ limit of (31) but is simpler to use. Because $`\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m)`$ is double valued for $`b=`$ constant $`2`$ , two expressions must be given to draw the $`b=486`$ curve, see Fig. 2. For the upper part of the curve:
$$\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m)=1\mathrm{\Omega }_m+3\sqrt{2/b}\mathrm{\Omega }_m\mathrm{cosh}\left[\frac{\mathrm{cosh}^1\left[\sqrt{b/2}(\mathrm{\Omega }_m^11)\right]}{3}\right],$$
(43)
where $`0\mathrm{\Omega }_m1/(1\sqrt{2/b})`$. In this expression hyperbolic cosine analytically becomes cosine for $`\mathrm{\Omega }_m1/(1+\sqrt{2/b})`$. For the lower part of the curve:
$$\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m)=1\mathrm{\Omega }_m+3\sqrt{2/b}\mathrm{\Omega }_m\mathrm{cos}\left[\frac{\mathrm{cos}^1\left[\sqrt{b/2}(1\mathrm{\Omega }_m^1)\right]+\pi }{3}\right],$$
(44)
where $`1\mathrm{\Omega }_m1/(1\sqrt{2/b})`$. The simplified result is:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m),\nu =1;z)={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{(1+z)}{\sqrt{|1\mathrm{\Omega }_0|}(33)^{(3/4)}}}\sqrt{\left[3{\displaystyle \frac{\mathrm{\Omega }_m(1+z)}{(1\mathrm{\Omega }_0)}}\right]\left[3{\displaystyle \frac{\mathrm{\Omega }_m}{(1\mathrm{\Omega }_0)}}\right]}\times `$
$`\{F(\varphi _0,{\displaystyle \frac{\sqrt{33}+5}{2\sqrt{33}}})F(\varphi _z,{\displaystyle \frac{\sqrt{33}+5}{2\sqrt{33}}})2[E(\varphi _0,{\displaystyle \frac{\sqrt{33}+5}{2\sqrt{33}}})E(\varphi _z,{\displaystyle \frac{\sqrt{33}+5}{2\sqrt{33}}})]`$
$`+2(33)^{(1/4)}[{\displaystyle \frac{\sqrt{8+\left[2+\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]^2}}{\sqrt{3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}\left[3+\sqrt{33}\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]}}`$
$`{\displaystyle \frac{\sqrt{8+\left[2+(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]^2}}{\sqrt{3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}\left[3+\sqrt{33}(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)\right]}}]\}.`$ (45)
The arguments of the elliptic functions, $`\varphi _z`$ and $`\varphi _0`$, can be calculated from (14) using $`y_1=3`$ and $`A=\sqrt{33}`$. To reduce the number of elliptic functions needed to evaluate (45), addition formulas for $`F(\varphi ,\mathrm{k})`$ and $`E(\varphi ,\mathrm{k})`$ can be used \[see (7), (48), and (49)\]. The value of $`\mathrm{\Delta }\varphi _z`$ is given by (15).
III. $`\nu =2`$, Empty-Beam Observations
A. Three Open Big Bang Domains
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =2;z)`$ $`=`$ $`{\displaystyle \frac{c}{H_0}}(1+z)^2{\displaystyle _0^z}{\displaystyle \frac{dz}{(1+z)^2\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}}.`$ (46)
Like the $`\nu =1`$ case this integral takes on different forms when evaluated in terms of Legendre elliptic integrals, depending on the value of the parameter $`b`$.
1. For the two open domains defined by $`b<0`$ and $`2<b`$ the luminosity distance $`D_{\mathrm{}}`$ takes the form:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =2;z)`$ $`=`$ $`{\displaystyle \frac{c}{H_0}}{\displaystyle \frac{(1+z)^2}{\mathrm{\Omega }_\mathrm{\Lambda }}}\{`$ (47)
$`(A+\kappa y_1)\left[{\displaystyle \frac{\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}{(1+z)[(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+A+\kappa y_1]}}{\displaystyle \frac{1}{\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+A+\kappa y_1}}\right]`$
$`{\displaystyle \frac{(A\kappa y_1)\sqrt{|1\mathrm{\Omega }_0|}}{2\sqrt{A}}}[F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})]+\sqrt{A}\sqrt{|1\mathrm{\Omega }_0|}[E(\varphi _z,\mathrm{k})E(\varphi _0,\mathrm{k})]\}`$
where $`y_1,A,\mathrm{k},`$ and $`\varphi _z`$ are defined in (10)-(14).<sup>4</sup><sup>4</sup>4 $`E(\varphi ,\mathrm{k})_0^\varphi \sqrt{1\mathrm{k}^2\mathrm{sin}^2\varphi }𝑑\varphi `$ Just as with the result for the $`\nu =1`$ case, i.e., (31), the number of Legendre elliptic integrals required to evaluate (47) can be reduced from four to two by using addition formulas 116.01 of Byrd & Friedman (1971). The addition formula for $`E(\varphi ,\mathrm{k})`$ is:
$$E(\varphi _z,\mathrm{k})E(\varphi _0,\mathrm{k})=E(\mathrm{\Delta }\varphi _z,\mathrm{k})\mathrm{k}^2\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{sin}\mathrm{\Delta }\varphi _z.$$
(48)
For this case
$`\mathrm{k}^2\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{sin}\mathrm{\Delta }\varphi _z={\displaystyle \frac{2\left[2A+\kappa (1+3y_1)\right]}{\left[(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)y_1\kappa A\right]}}`$
$`\times {\displaystyle \frac{\sqrt{\left[(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)y_1\right]\left[\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)y_1\right]}}{\left[\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)y_1\kappa A\right]\left[\mathrm{tan}(\mathrm{\Delta }\varphi _z/2)+1/\mathrm{tan}(\mathrm{\Delta }\varphi _z/2)\right]}},`$ (49)
where an expression for $`\mathrm{tan}(\mathrm{\Delta }\varphi _z/2)`$ is given by (15).
2. For the domain $`0<b<2`$ the luminosity distance $`D_{\mathrm{}}`$ takes the form:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda },\nu =2;z)=`$
$`{\displaystyle \frac{c}{H_0}}{\displaystyle \frac{(1+z)^2}{\mathrm{\Omega }_\mathrm{\Lambda }}}\{y_3[{\displaystyle \frac{\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}{(1+z)[(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+y_3]}}{\displaystyle \frac{1}{\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+y_3}}]`$
$`{\displaystyle \frac{y_2\sqrt{|1\mathrm{\Omega }_0|}}{\sqrt{y_1y_2}}}[F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})]\sqrt{y_1y_2}\sqrt{|1\mathrm{\Omega }_0|}[E(\varphi _z,\mathrm{k})E(\varphi _0,\mathrm{k})]\},`$ (50)
where the constants $`y_1,y_2,y_3`$ and $`\mathrm{k}`$ are defined in (16)-(18) but the function $`\varphi _z`$ is now defined as
$`\varphi _z=\varphi (\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda };z)`$ $`=`$ $`\mathrm{sin}^1\sqrt{{\displaystyle \frac{(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+y_2}{(1+z)\mathrm{\Omega }_m/|1\mathrm{\Omega }_0|+y_3}}}.`$ (51)
For this case the value of $`\mathrm{\Delta }\varphi _z`$ needed to reduce the number of elliptic integrals is the NEGATIVE of that given by (20) for the $`\nu =0`$ case. When the addition formula (48) is used, an additional term is contributed to (50) which can be evaluated using,
$`\mathrm{k}^2\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{sin}\mathrm{\Delta }\varphi _z={\displaystyle \frac{\mathrm{\Omega }_m(y_1y_3)|1\mathrm{\Omega }_0|^{(3/2)}(y_1y_2)^{(1/2)}}{[(1+z)\mathrm{\Omega }_m^2/(1\mathrm{\Omega }_0)^2(2+z)y_1\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)2y_1(1+y_1)]}}`$
$`\times \left\{{\displaystyle \frac{[y_2\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]\sqrt{(1+z)^2(1+\mathrm{\Omega }_mz)z(z+2)\mathrm{\Omega }_\mathrm{\Lambda }}}{[y_3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]}}{\displaystyle \frac{[y_2(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]}{[y_3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)]}}\right\}.`$ (52)
B. Boundaries
1. $`\mathrm{\Omega }_0\mathrm{\Omega }_m+\mathrm{\Omega }_\mathrm{\Lambda }=1`$
This case is the $`b\pm \mathrm{}`$ limit of (46) and a simpler expression containing hypergeometric functions results:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =2;z)`$
$`={\displaystyle \frac{c}{H_0}}(1+z)^2\{1{\displaystyle \frac{1}{(1+z)\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}}`$
$`+{\displaystyle \frac{3}{5}}\mathrm{\Omega }_m^{1/3}[\left({\displaystyle \frac{1}{[1+\mathrm{\Omega }_mz(3+3z+z^2)]^{5/6}}}\right){}_{2}{}^{}F_{1}^{}({\displaystyle \frac{5}{6}},{\displaystyle \frac{1}{3}};{\displaystyle \frac{11}{6}};{\displaystyle \frac{1\mathrm{\Omega }_m}{1+\mathrm{\Omega }_mz(3+3z+z^2)}})`$
$`{}_{2}{}^{}F_{1}^{}({\displaystyle \frac{5}{6}},{\displaystyle \frac{1}{3}};{\displaystyle \frac{11}{6}};1\mathrm{\Omega }_m)]\}.`$ (53)
When $`\mathrm{\Omega }_m1`$, (53) can be expressed in terms of associated Legendre functions as,
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =2;z)=`$
$`{\displaystyle \frac{c}{H_0}}(1+z)^2\{1{\displaystyle \frac{1}{(1+z)\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}}+{\displaystyle \frac{\mathrm{\Gamma }\left(5/6\right)}{2^{1/6}}}\left[{\displaystyle \frac{\mathrm{\Omega }_m}{1\mathrm{\Omega }_m}}\right]^{5/12}`$
$`\times [{\displaystyle \frac{(1+z)^{1/4}}{\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}}\mathrm{P}_{1/6}^{5/6}\left(\sqrt{{\displaystyle \frac{1+\mathrm{\Omega }_mz(3+3z+z^2)}{\mathrm{\Omega }_m(1+z)^3}}}\right)\mathrm{P}_{1/6}^{5/6}\left({\displaystyle \frac{1}{\sqrt{\mathrm{\Omega }_m}}}\right)]\}.`$ (54)
When $`\mathrm{\Omega }_m1`$, (53) can also be expressed in terms of Legendre elliptic integrals as,
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=1\mathrm{\Omega }_m,\nu =2;z)`$ $`=`$ $`{\displaystyle \frac{c}{H_0}}{\displaystyle \frac{(1+z)^2}{1\mathrm{\Omega }_m}}\{`$ (55)
$`(\sqrt{3}+1)\left(\mathrm{\Omega }_m^11\right)^{1/3}\left[{\displaystyle \frac{\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}}{(1+z)[1+z+(\sqrt{3}+1)\left(\mathrm{\Omega }_m^11\right)^{1/3}]}}{\displaystyle \frac{1}{1+(\sqrt{3}+1)\left(\mathrm{\Omega }_m^11\right)^{1/3}}}\right]`$
$`{\displaystyle \frac{1}{(\sqrt{3}+1)(3)^{1/4}}}\sqrt{\mathrm{\Omega }_m}\left(\mathrm{\Omega }_m^11\right)^{1/6}\left[F(\varphi _z,\mathrm{k})F(\varphi _0,\mathrm{k})\right]`$
$`+(3)^{1/4}\sqrt{\mathrm{\Omega }_m}(\mathrm{\Omega }_m^11)^{1/6}[E(\varphi _z,\mathrm{k})E(\varphi _0,\mathrm{k})]\},`$
where the constant $`\mathrm{k}`$ is given by (24) and the functions $`\varphi _z`$ and $`\mathrm{\Delta }\varphi _z`$ are given respectively by (25) and (26). For this case the additional term needed to use the addition formula (48) in (55) is:
$`\mathrm{k}^2\mathrm{sin}\varphi _z\mathrm{sin}\varphi _0\mathrm{sin}\mathrm{\Delta }\varphi _z`$
$`={\displaystyle \frac{z2(3)^{3/4}\left(2+\sqrt{3}\right)\sqrt{1\mathrm{\Omega }_m}}{\left[1+z+\left(1+\sqrt{3}\right)\left(\mathrm{\Omega }_m^11\right)^{1/3}\right]\left[1+\left(1+\sqrt{3}\right)\left(\mathrm{\Omega }_m^11\right)^{1/3}\right]}}`$
$`\times {\displaystyle \frac{\left\{z+\left[1+\left(\mathrm{\Omega }_m^11\right)^{1/3}\right]\left[1+\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}\right]\right\}}{\left\{2+3z\mathrm{\Omega }_m+z^2\mathrm{\Omega }_m\left[1+\left(1+\sqrt{3}\right)\left(\mathrm{\Omega }_m^11\right)^{1/3}\right]+2\sqrt{1+\mathrm{\Omega }_mz(3+3z+z^2)}\right\}}}.`$ (56)
2. $`b=2`$
See the description for the $`\nu =0`$ case in section I.B.2 including (44) for this “critical” value of $`b`$:
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }(\mathrm{\Omega }_m),\nu =2;z)`$
$`={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{9\mathrm{\Omega }_m(1+z)^2}{2|1\mathrm{\Omega }_0|^{3/2}}}\{{\displaystyle \frac{1}{(1+z)}}\sqrt{{\displaystyle \frac{1}{3}}{\displaystyle \frac{(1+z)\mathrm{\Omega }_m}{1\mathrm{\Omega }_0}}}\sqrt{{\displaystyle \frac{1}{3}}{\displaystyle \frac{\mathrm{\Omega }_m}{1\mathrm{\Omega }_0}}}`$
$`+{\displaystyle \frac{\mathrm{\Omega }_m}{1\mathrm{\Omega }_0}}\mathrm{log}\left[{\displaystyle \frac{1+\sqrt{1/3(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}}{1+\sqrt{1/3\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}}}\sqrt{{\displaystyle \frac{2/3+\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}{2/3+(1+z)\mathrm{\Omega }_m/(1\mathrm{\Omega }_0)}}}\right]\}.`$ (57)
3. $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$
This result was first given by Dyer & Roeder (1972),
$`D_{\mathrm{}}(\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }=0,\nu =2;z)`$
$`={\displaystyle \frac{c}{H_0}}{\displaystyle \frac{\mathrm{\Omega }_m(1+z)^2}{4(1\mathrm{\Omega }_m)^{3/2}}}[{\displaystyle \frac{3\mathrm{\Omega }_m}{2(1\mathrm{\Omega }_m)}}\mathrm{ln}\left\{\left({\displaystyle \frac{1+\sqrt{1\mathrm{\Omega }_m}}{1\sqrt{1\mathrm{\Omega }_m}}}\right)\left({\displaystyle \frac{\sqrt{1+\mathrm{\Omega }_mz}\sqrt{1\mathrm{\Omega }_m}}{\sqrt{1+\mathrm{\Omega }_mz}+\sqrt{1\mathrm{\Omega }_m}}}\right)\right\}`$
$`+{\displaystyle \frac{3}{\sqrt{1\mathrm{\Omega }_m}}}({\displaystyle \frac{\sqrt{1+\mathrm{\Omega }_mz}}{1+z}}1)+{\displaystyle \frac{2\sqrt{1\mathrm{\Omega }_m}}{\mathrm{\Omega }_m}}(1{\displaystyle \frac{\sqrt{1+\mathrm{\Omega }_mz}}{(1+z)^2}})],`$ (58)
and can be rewritten using the identity
$$\mathrm{sinh}^1\sqrt{\frac{1\mathrm{\Omega }_m}{\mathrm{\Omega }_m(1+z)}}=\frac{1}{2}\mathrm{ln}\left(\frac{\sqrt{1+\mathrm{\Omega }_mz}+\sqrt{1\mathrm{\Omega }_m}}{\sqrt{1+\mathrm{\Omega }_mz}\sqrt{1\mathrm{\Omega }_m}}\right).$$
(59)
When $`\mathrm{\Omega }_m>1`$ equation (58) is analytically continued using $`\sqrt{1\mathrm{\Omega }_m}\pm i\sqrt{\mathrm{\Omega }_m1}`$, which simplifies by using, $`\mathrm{sinh}^1(ix)=i\mathrm{sin}^1(x)`$ to give a form containing only real variables. The $`\mathrm{\Omega }_m=1`$ result for all $`\nu `$ was given by Dashevskii & Slysh (1966):
$$D_{\mathrm{}}(\mathrm{\Omega }_m=1,\mathrm{\Omega }_\mathrm{\Lambda }=0,\nu ;z)=\frac{c}{H_0}\frac{1}{(\nu +\frac{1}{2})}\left[(1+z)^{(\frac{\nu }{2}+1)}(1+z)^{(\frac{\nu }{2}+\frac{1}{2})}\right].$$
(60)
4. $`\mathrm{\Omega }_m=0`$
This result is exactly the same as the $`\nu =0`$ and $`\nu =1`$ result (29). If there is no mass in the universe then removing 100% of no mass from the beam removes nothing.
## 4 Conclusions
We have given useful forms for the luminosity distance in three currently relevant cosmologies. They are all dynamically FLRW cosmologies in the large but differ in how gravitating matter effects optical observations. The models are labeled by an additional parameter $`\nu `$ ($`\nu `$ = 0, 1, and 2) beyond the familiar $`H_0,\mathrm{\Omega }_m,`$ and $`\mathrm{\Lambda }`$. The $`\nu =0`$ model is standard FLRW where all matter is homogeneous and transparent on the scale of the observing beam widths. This model is called the ‘filled-beam’ model. The $`\nu =2`$ model assumes the opposite; all matter is inhomogeneous and excluded from the observing beams. This extreme case is called the ‘empty-beam’ model. The $`\nu =1`$ model assumes that 1/3 of the mass density of the universe is excluded from observing beams and hence it is the ‘two-thirds filled-beam’ model. These three cases were singled out because their distance-redshift relations can be given in terms of incomplete elliptic integrals; functions which are universally available in computer libraries and very efficiently evaluated.<sup>5</sup><sup>5</sup>5 The results appearing in Section 3 have been coded and are posted at http://www.nhn.ou.edu$``$thomas/z2dl.html. This code is discussed in the Appendix and compared to the numerical integration times of Kayser et al. (1997). For the $`\nu =1`$ and 2 cases, somewhat simpler expressions than what we have given exist, but only for complex arguments of the elliptic integrals. We chose to give expressions whose arguments are real and which can be rapidly evaluated. Results are available for all $`0\nu 2`$ but only in terms of the less familiar and unavailable Heun functions, Kantowski (1998a). We have extended the flat space, $`\mathrm{\Omega }_0=1`$, results given here to arbitrary filling parameter $`\nu `$. These new results will be available shortly. Related results have been independantly found by Damianski et al. (2000). A calculation similar to the $`\nu =1`$ case given here is that of the age of the Universe as a function of redshift and can be found in Thomas & Kantowski (2000).
R. Kantowski wishes to thank VP for Research, E. Smith, for funds to support J.K. Kao’s visit to OU during the summer of 1998 when the first elliptic integral results were obtained. R. C. Thomas thanks P. Helbig for discussions of his code, see Kayser et al. (1997), and E. Baron for benchmarking discussions.
## Appendix A Appendix
One expected practical use of the results given in this paper is to speedup distance evaluations for the $`\nu =`$ 0,1,2 partially filled beam FLRW models. We have implemented and made publicly available a Fortran 90 version of this work called Z2DL (see http://www.nhn.ou.edu$``$thomas/z2dl.html for Z2DL with documentation and extensive CPU-time benchmark results). Z2DL uses Carlson elliptic integrals (see Press et al. (1994) and references therein) and results in a fast distance calculator. We have benchmarked Z2DL by comparing it with the commonly used and fast numerical integration routine ANGSIZ (see Kayser et al. (1997)). For a given ($`\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }`$), the total CPU-time required to convert $`5\times 10^5`$ redshifts (equally spaced between z=0 and z=5) to luminosity distance using Z2DL and ANGSIZ separately were recorded. By calculating the ratio of ANGSIZ CPU-time to Z2DL CPU-time on a grid of points in ($`\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }`$) we have generated three speedup surfaces, one for each value of $`\nu =`$ 0,1,2 (see Fig. 3 for the $`\nu =0`$ surface). The results for all three comparisons are given as contour plots at the web site. Using an IBM AIX 375 MHz Power III approximately 7 hours was required to generate each ($`\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }`$) grid of 30 x 30 points (minus models without a big bang).
For the purpose of a clearer presentation, we omitted speedup points along the $`\mathrm{\Omega }_m=0`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ lines. Along these boundaries speedup factors are greater than 100. The large open domains of the $`\mathrm{\Omega }_m`$-$`\mathrm{\Omega }_\mathrm{\Lambda }`$ plane, i.e., subsection ‘A’ cases, constitute the majority of models in the grid and also those with the least impressive speedup. However, even for these cases, the improvement is substantial: typically 17-20 for $`\nu =0`$ (standard filled beam FLRW), 6-8 for $`\nu =1`$ (66% filled beam FLRW), and 11-13 for $`\nu =2`$ (empty beam FLRW).
To gauge the level of agreement between distances computed by ANGSIZ and Z2DL, a finer grid of ($`\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }`$) with 3000 x 3000 points (between 0 and 3 in both directions, also excluding models without a big bang) was used. For each ($`\mathrm{\Omega }_m,\mathrm{\Omega }_\mathrm{\Lambda }`$), both routines were used to compute luminosity distance for z=1. Most often the results agree to within one part in $`10^6`$. Cases where disagreements greater than one part in $`10^3`$ occur are near the upper b=2 line (see Fig. 1). We found that ANGSIZ was giving less accurate distances near this boundary of non-big bang models as ANGSIZ documentation explains. |
warning/0002/hep-th0002008.html | ar5iv | text | # References
hep-th/0002008
$`IIB`$ supergravity and various aspects of light-cone formalism in AdS space-time <sup>1</sup><sup>1</sup>1 Based on talk given at International Workshop ”Supersymmetries and Quantum Symmetries”, Dubna, Russia, July 27-31, 1999
R.R. Metsaev
Department of Theoretical Physics, P.N. Lebedev Physical Institute, Leninsky prospect 53, 117924, Moscow, Russia
Abstract
Light-cone gauge manifestly supersymmetric formulation of type IIB 10-dimensional supergravity in $`AdS_5\times S^5`$ background is discussed. The formulation is given entirely in terms of light-cone scalar superfield, allowing us to treat all component fields on an equal footing. Discrete energy spectrum of field propagating in AdS space is explained within the framework of light-cone approach. Light-cone gauge formulation of self-dual fields propagating in AdS space is developed. An conjectured interrelation of higher spin massless fields theory in AdS space-time and superstring theory in Minkowski space is discussed.
Light-cone gauge formulation of $`IIB`$ supergravity in $`AdS_5\times S^5`$ background. In recent years, due to the conjectured duality between the string theory and $`𝒩=4`$, $`4d`$ SYM theory (for review see ) there has been rapidly growing interest in strings propagating in AdS space. Inspired by this conjecture the Green-Schwarz formulation of strings propagating in $`AdS_5\times S^5`$ was suggested in (for further developments see -). Despite considerable efforts these strings have not yet been quantized<sup>2</sup><sup>2</sup>2Some related interesting discussions are in . Alternative approaches can be found in ,. Twistor-like formulations are discussed in .. As is well known, quantization of GS superstrings propagating in flat space is straightforward only in the light-cone gauge. Because the string theories are approximated at low energies by supergravity theories it seems reasonable that we should first study a light-cone gauge formulation of supergravity theories in AdS spacetime. Light-cone description of type $`IIB`$ supergravity in $`AdS_5\times S^5`$ background can teach us a lot. For example it can be used to construct charges of global symmetries so that certain of them have the same form for both supergravity and superstring. Keeping in mind extremely important applications to string theory let us now restrict our discussion to $`IIB`$ supergravity in $`AdS_5\times S^5`$.
Our method is conceptually very close to the one used in (see also ) to find the light-cone form of $`IIB`$ supergravity in flat space and is based essentially on a light-cone gauge description of field dynamics developed recently in <sup>3</sup><sup>3</sup>3A discussion of $`IIB`$ supergravity at the level of gauge invariant equations of motion and actions can be found in and respectively.. To discuss light-cone formulation we use the following parametrization of $`AdS_5\times S^5`$ space
$$ds^2=\frac{1}{z^2}(dt^2+dx_1^2+dx_2^2+dz^2+dx_4^2)+\frac{1}{4}dy_{ij}dy_{ij}^{},$$
introduce light-cone variables $`x^\pm (x^4\pm x^0)/\sqrt{2}`$, $`x(x^1+\mathrm{i}x^2)/\sqrt{2}`$, $`\overline{x}x^{}`$, where $`x^0t`$ and treat $`x^+`$ as evolution parameter. Here and below we set the radii of both $`AdS_5`$ and $`S^5`$ equal to unity<sup>4</sup><sup>4</sup>4The $`S^5`$ coordinates $`y^{ij}`$ are subject to the constraints $`y^{ik}y_{kj}=\delta _j^i`$, $`y_{ij}=\frac{1}{2}ϵ_{ijkn}y^{kn}`$, $`y_{ij}^{}=y^{ij}`$ where indices take the values $`i,j,k,n=1,2,3,4`$ The $`ϵ^{ijkn}=\pm 1`$ is the Levi-Civita tensor of $`su(4)`$.. The coordinates $`y_{ij}`$ are related to the standard $`so(6)`$ cartesian coordinates $`y^M`$, $`M=1,\mathrm{},6`$, which satisfy the constraint $`y^My^M=1`$ through the formula $`y_{ij}=\rho _{ij}^My^M`$, where $`\rho _{ij}^M`$ are the Clebsh Gordan coefficients of $`su(4)`$ algebra .
Our goal is to find a realization of $`psu(2,2|4)`$ superalgebra on the space of IIB supergravity fields propagating in $`AdS_5\times S^5`$. To do that we use light-cone superspace formalism. First, we introduce light-cone superspace which by definition is based on position $`AdS_5\times S^5`$ coordinates $`x^\pm `$, $`x`$, $`\overline{x}`$, $`z`$, $`y^{ij}`$ and Grassmann position coordinates $`\theta ^i`$ and $`\chi ^i`$ which transform in fundamental $`\mathrm{𝟒}`$ irreps of $`su(4)`$ algebra. Second, on this superspace we introduce scalar superfield $`\mathrm{\Phi }(x^\pm ,x,\overline{x},z,y^{ij},\theta ^i,\chi ^i)`$. In the remainder of this paper we find it convenient to Fourier transform to momentum space for all coordinates except for $`x^+`$, $`z`$ and $`S^5`$ coordinates $`y^{ij}`$. This implies using $`p^+`$, $`\overline{p}`$, $`p`$, $`\lambda _i`$, $`\tau _i`$ instead of $`x^{}`$, $`x`$, $`\overline{x}`$, $`\theta ^i`$, $`\chi ^i`$ respectively. The $`\lambda _i`$ and $`\tau _i`$ transform in $`\overline{\mathrm{𝟒}}`$ irreps of $`su(4)`$. Thus we consider the superfield $`\mathrm{\Phi }(x^+,p^+,p,\overline{p},z,y^{ij},\lambda _i,\tau _i)`$<sup>5</sup><sup>5</sup>5An expansion of $`\mathrm{\Phi }`$ in powers of Grassmann momenta $`\lambda _i`$ and $`\tau _i`$ can be found in . which by definition satisfies the following reality constraint
$$\mathrm{\Phi }(x^+,p,z,y,\lambda ,\tau )=p^{+4}d^4\lambda ^{}d^4\tau ^{}e^{(\lambda _i\lambda _i^{}+\tau _i\tau _i^{})/p^+}(\mathrm{\Phi }(x^+,p,z,y,\lambda ,\tau ))^{}.$$
In the light-cone formalism the $`psu(2,2|4)`$ superalgebra has the generators
$$P^+,P,\overline{P},J^{+x},J^{+\overline{x}},K^+,K,\overline{K},Q^{+i},Q_i^+,S^{+i},S_i^+,D,J^+,J^{x\overline{x}},$$
(1)
which we refer to as kinematical generators and
$$P^{},J^x,J^{\overline{x}},K^{},Q^i,Q_i^{},S^i,S_i^{}$$
(2)
which we refer to as dynamical generators. The kinematical generators have positive or zero $`J^+`$ charges, while dynamical generators have negative $`J^+`$ charges<sup>6</sup><sup>6</sup>6For $`x^+=0`$ the kinematical generators are quadratic in the physical field $`\mathrm{\Phi }`$, while the dynamical generators receive corrections in interaction theory. Here we deal with free fields.. At a quadratic level both kinematical and dynamical generators have the following representation in terms of the physical light-cone superfield
$$\widehat{G}=𝑑p^+d^2p𝑑z𝑑S^5d^4\lambda d^4\tau p^+\mathrm{\Phi }(x^+,p,z,y,\lambda ,\tau )G\mathrm{\Phi }(x^+,p,z,y,\lambda ,\tau ),$$
where $`G`$ are the differential operators acting on $`\mathrm{\Phi }`$. Thus we should find representation of $`psu(2,2|4)`$ in terms of differential operators acting on light-cone superfield $`\mathrm{\Phi }`$. To simplify expressions let us write down the generators for $`x^+=0`$. The kinematical generators are then given by
$$P=p,\overline{P}=\overline{p},P^+=p^+,J^{+x}=_pp^+,J^{+\overline{x}}=\overline{}_pp^+,$$
(3)
$$Q^{+i}=p^+\theta ^i,Q_i^+=\lambda _i,S_i^+=\frac{1}{\sqrt{2}}z\tau _i\lambda _i_p,S^{+i}=\frac{1}{\sqrt{2}}zp^+\chi ^i+p^+\theta ^i\overline{}_p,$$
(4)
where $`_p/\overline{p}`$, $`\overline{}_p/p`$. Dynamical generator $`P^{}`$ is given by
$$P^{}=\frac{p\overline{p}}{p^+}+\frac{_z^2}{2p^+}\frac{1}{2z^2p^+}A,AX\frac{1}{4},$$
(5)
where
$$Xl^i{}_{j}{}^{2}+4\tau l\chi +(\chi \tau 2)^2,l^i{}_{j}{}^{2}l^i{}_{j}{}^{}l_{}^{j}{}_{i}{}^{},\tau l\chi \tau _il^i{}_{j}{}^{}\chi _{}^{i},\chi \tau \chi ^i\tau _i.$$
(6)
The $`l^i_j`$ is the orbital part of $`su(4)`$ angular momentum $`J^i_j`$. Explicit expressions for the remaining kinematical and dynamical generators may be found in <sup>7</sup><sup>7</sup>7For $`\lambda _i`$, $`\tau _i`$ and $`\tau ^i`$, $`\chi ^i`$ we adopt the following anticommutation rules $`\{\theta ^i,\lambda _j\}=\delta _j^i`$, $`\{\chi ^i,\tau _j\}=\delta _j^i`$.. Making use of the expression for $`P^{}`$ (5) we can immediately write down the light-cone gauge action<sup>8</sup><sup>8</sup>8Since the action is invariant with respect to the symmetries generated by $`psu(2,2|4)`$ superalgebra, the formalism we discuss is sometimes referred to as an off shell light-cone formulation .
$$S_{l.c.}=𝑑x^+𝑑p^+d^2p𝑑z𝑑S^5d^4\lambda d^4\tau p^+\mathrm{\Phi }(x^+,p,z,y,\lambda ,\tau )(\mathrm{i}^{}+P^{})\mathrm{\Phi }(x^+,p,z,y,\lambda ,\tau ).$$
Following we shall call the operator $`A`$ the $`AdS`$ mass operator. A few comments are in order.
(i) The operator $`A`$ is equal to zero only for massless representations realized as irreducible representations of conformal algebra , which for the case of $`AdS_5`$ space is the $`so(5,2)`$ algebra. Because the operator $`X`$ (6) has eigenvalues equal to squared integers in the whole spectrum of compactification of $`IIB`$ supergravity on $`AdS_5`$ (see ) the operator $`A`$ is never equal to zero. This implies that the scalar fields as well as all remaining fields of compactification of $`IIB`$ supergravity do not satisfy conformally invariant equations of motion.
(ii) The coordinate $`\theta ^i`$ (or $`\lambda _i`$) constitutes odd part of light-cone superspace appropriate to superfield description of light-cone gauge $`N=4`$, $`4d`$ (or $`N=1`$, $`10d`$) SYM theory. The translation generators, Lorentz boosts and Poincaré supercharges given in (3),(4) take the same form as in SYM theory. From this we conclude that one can expect that superstring dynamics in $`AdS_5\times S^5`$ could be presented as dynamics of free left and right movers appropriate to description of open superstrings supplemented by nonlinear dynamics of remaining degrees of freedom.
(iii) The generator $`P^{}`$ involves terms up to fourth order in $`\chi ^i`$ and $`\tau _i`$. The first thing to note is that terms of fourth order in $`\chi ^i`$ and $`\tau _i`$ appear trough the number operator $`\chi \tau `$. The second point is that terms of fourth order can be excluded by introducing new superfield<sup>9</sup><sup>9</sup>9We prefer to use superfield $`\mathrm{\Phi }`$ instead of $`\mathrm{\Phi }^{new}`$ because the $`\mathrm{\Phi }`$ has conventional canonical dimension. $`\mathrm{\Phi }=z^{\chi \tau 2}\mathrm{\Phi }^{new}`$. On the space of $`\mathrm{\Phi }^{new}`$ the hamiltonian $`P^{}`$ takes the form
$$P^{}=\frac{p\overline{p}}{p^+}+\frac{_z^2}{2p^+}+\frac{\chi \tau 2}{zp^+}_z\frac{1}{2z^2p^+}(l^i{}_{j}{}^{2}+4\tau l\chi +\chi \tau \frac{9}{4})$$
i.e. the terms of fourth order in $`\chi ^i`$ and $`\tau _i`$ are absent. This suggests (but does not prove) light-cone gauge string action which does not involve higher than second order terms in anticommuting variables<sup>10</sup><sup>10</sup>10Such action has been found recently in Ref...
Discrete energy spectrum of field in AdS space-time. As is well known energy spectrum of field propagating in AdS space takes discrete values (see for review ). This phenomenon is not immediately visible in light-cone formulation. In this section we would like to explore how the discrete energy spectrum is obtained within the framework of the light-cone formulation.
The AdS translation generators $`\widehat{P}^a`$ are expressible as<sup>11</sup><sup>11</sup>11In this section, in contrast to , we use hermitean $`\widehat{P}^a`$, $`P^a`$, $`K^a`$ which are related with the anithermitean ones of as follows $`\widehat{P}_{herm}^a=\mathrm{i}\widehat{P}_{antherm}^a`$, $`P_{herm}^a=\mathrm{i}P_{antiherm}^a`$, $`K_{herm}^a=+\mathrm{i}K_{antiherm}^a`$. $`\widehat{P}^a=(1/2)P^a+K^a`$. The $`\widehat{P}^0`$ is the energy operator. In light-cone frame we have<sup>12</sup><sup>12</sup>12We use parametrization of $`AdS_d`$ space in which $`ds^2=(dx^{02}+dx_I^2+dx_{d1}^2)/z^2`$. Light-cone coordinates in $`\pm `$ directions are defined as $`x^\pm =(x^{d1}\pm x^0)/\sqrt{2}`$. From now on we adopt the following conventions: $`I,J=1,\mathrm{},d2`$; $`i,j,k,l=1,\mathrm{},d3`$. $`^I=_I/x^I`$, $`^+=_{}/x^{}`$, $`_{p^+}/p^+`$, $`zx^{d2}`$. In momentum representation $`^+`$ takes the form $`^+=\mathrm{i}p^+`$.
$$\widehat{P}^0=\frac{1}{2\sqrt{2}}(P^++2K^+P^{}2K^{}).$$
Representation of generators $`P^a`$ and $`K^a`$ on the space of physical fields has been found in . For the case of totally symmetric fields this representation takes the form (for $`x^+=0`$)
$`P^+=p^+,K^+={\displaystyle \frac{1}{2}}x_I^2p^+,x_I^2x^Ix^I`$
$`P^{}={\displaystyle \frac{_I^2}{2p^+}}{\displaystyle \frac{1}{2z^2p^+}}({\displaystyle \frac{1}{2}}M_{ij}^2+{\displaystyle \frac{(d4)(d6)}{4}}),_I^2^I^I`$
$`K^{}={\displaystyle \frac{1}{2}}x_I^2P^{}_{p^+}D+{\displaystyle \frac{1}{2p^+}}l^{IJ}M^{IJ}{\displaystyle \frac{x^I}{2zp^+}}\{M^{zJ},M^{JI}\},`$
$`D=_{p^+}p^++x^I^I+{\displaystyle \frac{d2}{2}},M^{IJ}\alpha ^I\overline{\alpha }^J\alpha ^J\overline{\alpha }^I.`$ (7)
where $`l^{IJ}x^I^Jx^J^I`$ and we adopt a convention: $`\{a,b\}ab+ba`$. The above generators act on physical field whose components are collected in Fock vector $`|\varphi `$:
$$|\varphi =\varphi ^{I_1\mathrm{}I_s}\alpha ^{I_1}\mathrm{}\alpha ^{I_s}|0,\overline{\alpha }^I|0=0,[\overline{\alpha }^I,\alpha ^J]=\delta ^{IJ},$$
(8)
where physical traceless tensor field $`\varphi ^{I_1\mathrm{}I_s}`$ depends on $`x^+,x^i,z,p^+`$. Using these expressions we get the following representation for energy operator (for $`x^+=0`$)
$$2\sqrt{2}\widehat{P}^0=(1+x_I^2)(p^+P^{})+2_{p^+}D\frac{1}{p^+}l^{IJ}M^{IJ}+\frac{x^I}{zp^+}\{M^{zJ},M^{JI}\}.$$
(9)
It is convenient to make the following transformation of wave function
$$\varphi =U\stackrel{~}{\varphi },U(p^+)^{\frac{1}{2}(x^I^I+\frac{d5}{2})}$$
and use the formula
$$_{p^+}D\varphi =U\left(p^+_{p^+}^2\frac{1}{2}_{p^+}+\frac{1}{4p^+}(x_I^2_J^2\frac{1}{2}l_{IJ}^2+\frac{(d3)(d5)}{4})\right)\stackrel{~}{\varphi }.$$
(10)
Now to simplify our presentation let us restrict ourselves to the case of four dimensional ($`d=4`$) AdS space. For this case because the indices $`i,j`$ which label $`d3`$ directions take one value $`i,j=1`$ the spin operator $`M^{ij}`$ is equal to zero. Now taking into account that $`M_{IJ}^2=2M_{z1}^2`$ and exploiting (10) in (9) we get the following representation of $`\widehat{P}^0`$ in $`\stackrel{~}{\varphi }`$
$$2\sqrt{2}\widehat{P}^0=H_1+H_2,$$
(11)
where
$`H_1{\displaystyle \frac{1}{2}}_I^2+x_I^2,`$ (12)
$`H_22p^+_{p^+}^2_{p^+}{\displaystyle \frac{1}{2p^+}}\left(2({\displaystyle \frac{1}{2}}l^{IJ}+M^{IJ})^2+{\displaystyle \frac{1}{4}}\right)+p^+.`$ (13)
Because the hamiltonians $`H_1`$, $`H_2`$ commute with each other we can diagonalize them simultaneously. Let us find their eigenvalues. To this end we introduce complex coordinates $`x`$, $`\overline{x}`$ instead of $`x^1,z`$: $`x(x^1+\mathrm{i}z)/\sqrt{2}`$, $`\overline{x}x^{}`$. Next we decompose the wave function $`|\stackrel{~}{\varphi }`$ as follows
$$|\stackrel{~}{\varphi }=|\varphi _{+s}+|\varphi _s,$$
where the $`|\varphi _{\pm s}`$ are eigenvectors of $`M^{x\overline{x}}`$: $`M^{x\overline{x}}|\varphi _{\pm s}=\pm s|\varphi _{\pm s}`$. The $`|\varphi _{\pm s}`$ themselves can be decomposed into eigenvectors of operator $`l^{x\overline{x}}`$
$$|\varphi _{\pm s}=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}e^{\mathrm{i}m\phi }|\varphi _{\pm s,m},l^{x\overline{x}}e^{\mathrm{i}m\phi }=me^{\mathrm{i}m\phi },\phi \text{arg}(x^1+\mathrm{i}z).$$
On space of $`|\varphi _{\pm s,m}`$ the hamiltonians $`H_1`$, $`H_2`$ take the form
$`H_1={\displaystyle \frac{1}{2}}(_r^2+{\displaystyle \frac{1}{r}}_r)+{\displaystyle \frac{m^2}{2r^2}}+r^2,`$ (14)
$`H_2=2p^+_{p^+}^2_{p^+}+{\displaystyle \frac{1}{2p^+}}\left(\kappa ^2{\displaystyle \frac{1}{4}}\right)+p^+,`$ (15)
where $`r`$ is a radial variable $`r|x^1+\mathrm{i}z|`$ and $`\kappa m\pm 2s`$. Because $`H_1`$ and $`H_2`$ commute with each other we can decompose the wave function as follows
$$|\varphi _{\pm s,m}=\varphi _{\pm s,m}^{(1)}(r)|\varphi _{\pm s,m}^{(2)}(p^+),$$
where $`\varphi _{\pm s,m}^{(1)}(r)`$ and $`|\varphi _{\pm s,m}^{(2)}(p^+)`$ are eigenvectors of $`H_1`$ and $`H_2`$ respectively. Introducing instead of $`p^+`$ a new variable $`y`$, by relation $`y^2=p^+`$ and rescaling wave function
$$\varphi _{\pm s,m}^{(1)}(r)=r^{1/2}\stackrel{~}{\varphi }_{\pm s,m}^{(1)}(r),$$
we get the following hamiltonians
$`H_1={\displaystyle \frac{1}{2}}_r^2+{\displaystyle \frac{1}{2r^2}}(m^2{\displaystyle \frac{1}{4}})+{\displaystyle \frac{\omega _0^2}{2}}r^2,`$ (16)
$`H_2={\displaystyle \frac{1}{2}}_y^2+{\displaystyle \frac{1}{2y^2}}(\kappa ^2{\displaystyle \frac{1}{4}})+{\displaystyle \frac{\omega _0^2}{2}}y^2,`$ (17)
where $`\omega _0\sqrt{2}`$. The eigenvalues of the hamiltonians (16) and (17) responsible for square integrable eigenvectors are well known and are given by
$$E_1=(2n_1+|m|+1)\omega _0,E_2=(2n_1+|m\pm 2s|+1)\omega _0,$$
where $`n_1,n_2=0,1,\mathrm{}`$. Taking into account the relation (11) we get the following eigenvalues of $`\widehat{P}^0`$
$$E=n_1+n_2+|\frac{m}{2}\pm s|+|\frac{m}{2}|+1$$
(18)
A few comments are in order. (i) Our energy spectrum (as usual) is discrete. As compared to standard energy spectrum (see ) which depends on two integers our energy spectrum (18) depends on three integers $`n_1`$, $`n_2`$, $`m`$. Such the difference is well know from analysis of quantum mechanical energy spectrum of the standard three dimensional oscillator and is related to a choice of coordinates in which energy spectrum is evaluated. Normally, (see ) one uses global spherical coordinates while we use Poincaré coordinates. (ii) As usual for massless spin $`s`$ particle in $`AdS_4`$ our energy spectrum is bounded from below by $`E_{min}=s+1`$. (iii) There is 2-fold degeneracy of lowest energy $`E_{min}`$ related to the two helicity states $`|\varphi _{\pm s}`$. (iv) for each helicity state lowest energy has $`(2s+1)`$-fold degeneracy, i.e. for $`n_1=n_2=0`$ there exist $`2s+1`$ values of $`m`$ for which $`E`$ given in (18) is equal to $`E_{min}`$. This degeneracy explains well known $`so(3)`$ symmetry of lowest energy state.
Self-dual fields in AdS space-time. In this section we would like to discuss self-dual fields in AdS space. In our knowledge they have not been discussed previously in literature. As in Minkowski space the strengths of AdS (anti) self-dual fields satisfy self-duality constraint
$$F_\pm ^{\mu _1\mathrm{}\mu _{d/2}}=\pm \frac{1}{(d/2)!}\frac{ϵ^{\mu _1\mathrm{}\mu _{d/2}\nu _1\mathrm{}\nu _{d/2}}}{\sqrt{|g|}}F_{\pm \nu _1\mathrm{}\nu _{d/2}}$$
(19)
Because this constraint it is impossible to construct an appropriate Lorentz covariant action without introducing auxiliary fields<sup>13</sup><sup>13</sup>13Lorentz covariant formulations including auxiliary fields are discussed in ,.. The action for self-dual fields can be most easily understood within the framework of light-cone gauge formulation. For definiteness let us restrict ourselves to the case of six dimensional AdS space and to second rank antisymmetric tensor field. First of all we would like to discuss formulation of self-dual fields in Minkowski space which is most convenient for generalization to AdS space. In light-cone gauge physical degrees of freedom of self-dual field in Minkowski space are described by field $`\varphi ^{IJ}`$ which satisfies the self-duality constraint $`\varphi ^{IJ}=(1/2)ϵ^{IJKL}\varphi ^{KL}`$. The field $`\varphi ^{IJ}`$, which is the $`so(4)`$ tensor, can be decomposed into $`so(3)`$ tensors $`\varphi ^{ij}`$ and $`\varphi ^i`$, $`\varphi ^i\varphi ^{zi}`$. The self-duality constraint tells us then that $`\varphi ^{ij}`$ is expressible in terms of $`\varphi ^i`$: $`\varphi ^{ij}=ϵ^{ijk}\varphi ^k`$. Thus, if one wishes, in Minkowski space the self-dual field can be described by unconstrained field $`\varphi ^i`$ (or $`\varphi ^{ij}`$) <sup>14</sup><sup>14</sup>14It is worth mentioning that formulation suggested in (and its generalization to curved space given in ) is extremely appropriate to description of self-dual fields in AdS space. The formulation given in breaks Lorentz invariance $`so(d1,1)`$ to $`so(d2,1)`$ invariance and this does not fit with manifest Lorentz symmetry of Minkowski space metric. On the other hand in AdS space it is the $`so(d2,1)`$ symmetry that is manifest symmetry of $`AdS_d`$ space metric considered in Poincaré coordinates (equivalently, the $`so(d2,1)`$ is a manifest symmetry of AdS algebra considered in conformal algebra notation).. Note that this form of description breaks manifest $`so(4)`$ (which is $`so(d2)`$ for $`d=6`$) invariance to manifest $`so(3)`$ (which is $`so(d3)`$ for $`d=6`$) invariance. On the other hand light-cone gauge formalism in $`AdS_d`$ space respects only manifest $`so(d3)`$ (see ). Therefore the description based on unconstrained field $`\varphi ^i`$ is most appropriate to be generalized to AdS space. Thus our aim is to find realization of AdS algebra generators on the space of unconstrained field $`\varphi ^i`$. The realization of AdS algebra on the space of physical fields found in (, formulas (4.1)-(4.8) and (4.17)-(4.19)) is formulated in terms of spin operator $`M^{IJ}`$ and operators $`A`$, $`B`$. The form of spin operator is fixed by representation of $`so(4)`$ algebra we interested in. To proceed we introduce creation operator $`\alpha ^i`$ and construct Fock space vector
$$|\varphi \varphi ^i\alpha ^i|0,\overline{\alpha }^i|0=0,[\overline{\alpha }^i,\alpha ^j]=\delta ^{ij}.$$
(20)
For this form of realization the spin operators take the form
$$M^{ij}=\alpha ^i\overline{\alpha }^j\alpha ^j\overline{\alpha }^i,M^{zi}=\frac{1}{2}ϵ^{ijk}M^{jk}.$$
(21)
Note that it is second relation in (21) that tells that we deal with self-dual representation of $`so(4)`$ algebra. As to above mentioned operators $`A`$ and $`B`$ they are fixed by defining equations found in ()
$`2\{M^{zi},A\}[[M^{zi},A],A]=0,`$ (22)
$`[M^{zi},[M^{zj},A]]+\{M^{iL},M^{Lj}\}=2\delta ^{ij}B.`$ (23)
where operator $`A`$ is invariant under $`so(d3)`$ spin rotations, i.e. $`[A,M^{ij}]=0`$. From $`[A,M^{ij}]=0`$, the second relation in (21) and equations (22) we find the equation $`\{A,M^{zi}\}=0`$ which implies that $`A=0`$. By using (21) it is easy to get then the following relation
$$\{M^{iL},M^{Lj}\}=2\delta ^{ij}M_{zl}^2,M_{zl}^2M^{zl}M^{zl}.$$
From this relation and (23) we conclude that $`B=M_{zi}^2`$. It is easy to check that in the case under consideration the $`M_{zi}^2`$ is diagonalized: $`M_{zi}^2|\varphi =2|\varphi `$. Due to (21) the $`M_{zi}^2`$ commutes with $`M^{IJ}`$ and therefore we can put $`B=2`$. To summarize we have found
$$A=0,B=2.$$
(24)
Thus we found all entries of light-cone formulation. By substituting the expressions for spin operators (21) and operators $`A`$ and $`B`$ (24) in expressions (4.1)-(4.8) and (4.17)-(4.19) of Ref. we find realization of AdS algebra generators on space of AdS self-dual field $`|\varphi `$ (20). This provides complete description of self-dual field $`|\varphi `$ (20) in $`AdS_6`$ space.
Generalization to the case of arbitrary spin $`s`$ self-dual fields in $`AdS_6`$ is straightforward. In this case we introduce
$$|\varphi =\varphi ^{i_1\mathrm{}i_s}\alpha ^{i_1}\mathrm{}\alpha ^{i_s}|0,$$
where $`\varphi ^{i_1\mathrm{}i_s}`$ is totally symmetric traceless $`so(3)`$ tensor field. The spin operators take the form given in (21) while for operators $`A`$ and $`B`$ we get
$$A=0,B=s(s+1).$$
The light-cone gauge action for self-dual field in AdS takes then the form
$$S=d^dx^+\varphi |(^{}+P^{})|\varphi ,P^{}=\frac{_I^2}{2^+},$$
which coincides with the one in Minkowski space. We think that this coincidence can be traced to the conformal invariance of self-dual fields . Note that to describe anti self-dual fields one needs to replace the second relation in (21) by relation $`M^{zi}=(1/2)ϵ^{ijk}M^{jk}`$. All remaining formulas above given do not change their form.
Superstring theory and AdS higher spin massless fields theory. Conjecture. One of major motivations for our investigation of higher spin massless fields theory in the AdS space is to seek a possible relation between this theory and string theory. It seems rather attractive to conjecture that string theory can be interpreted as resulting from some kind of a spontaneous breakdown of higher spin symmetries. In it has been conjectured that superstrings could be considered as the ones living at the boundary of $`11`$-dimensional AdS space while their unbroken (symmetric) phase is realized as the theory of higher spin massless fields living in this $`AdS_{11}`$ space. AdS theories are symmetric with respect to isometry algebra of $`AdS_d`$ space which is $`so(d1,2)`$. This algebra is not realized however as symmetry algebra of string S-matrix . This implies that $`so(d1,2)`$ symmetry should be (spontaneously) broken. Here we would like to demonstrate that if one restricts attention to totally symmetric fields and make some mild assumptions about the (spontaneously) broken form of AdS theory hamiltonian $`P^{}`$ then some interesting and non-trivial test of our conjecture can be carried out. To this end we consider the $`P^{}`$ for AdS totally symmetric fields (8). As was demonstrated in one has
$$P^{}=\frac{_I^2}{2^+}+\frac{1}{2z^2^+}A,A=\frac{1}{2}M_{ij}^2+\frac{(d4)(d6)}{4}.$$
(25)
where $`M^{ij}`$ is given in (7). This $`P^{}`$ can be rewritten as
$$P^{}=\frac{_i^2+M^2}{2^+},M^2_z^2+\frac{1}{z^2}A.$$
(26)
The operator $`M^2`$ in (26) can be interpreted as mass operator for a field propagating in $`(d1)`$ dimensional Minkowski spacetime while the $`P^{}`$ (26) can be considered as hamiltonian of this field. The operator $`M^2`$ given in (26) has continuous spectrum. This implies that AdS theories which are not supplemented by (spontaneous) symmetry breaking lead to boundary theories which have continuous mass spectrum. In order to get theory with discrete mass spectrum one needs to break AdS symmetry. Now we have to make assumption about term which breaks AdS symmetry. Our suggestion is to consider the following ansatz for (spontaneously) broken $`P_{s.b}^{}`$
$$P_{s.b}^{}=P^{}+\frac{\omega ^2z^2\overline{\omega }}{2^+},\omega \frac{1}{2\alpha ^{}},\overline{\omega }\frac{d1}{2\alpha ^{}},$$
(27)
where $`P^{}`$ is given in (25) and $`\alpha ^{}`$ is the universal Regge slope parameter<sup>15</sup><sup>15</sup>15Here we discuss open AdS and string theories. For the case of closed theories the $`\omega `$ and $`\overline{\omega }`$ in (27) should be multiplied by factor 2.. For this $`P_{s.b}^{}`$ the mass operator takes the form
$$M_{s.b}^2=_z^2+\frac{1}{z^2}A+\omega ^2z^2\overline{\omega }.$$
(28)
This $`M_{s.b}^2`$, in contrast to $`M^2`$ given in (26), has discrete spectrum. Let us evaluate this spectrum. To this end we decompose the field $`|\varphi `$, which transforms in irreducible representation of $`so(d2)`$ algebra, into irreducible representations of $`so(d3)`$ subalgebra $`|\varphi _s^{}`$
$$|\varphi =\underset{s^{}=0}{\overset{s}{}}|\varphi _s^{}.$$
(29)
Because of relation $`M_{ij}^2|\varphi _s^{}=2s^{}(s^{}+d5)|\varphi _s^{}`$ the operator $`M_{s.b}^2`$ for $`|\varphi _s^{}`$ takes the form
$$M_{s.b}^2=_z^2+\frac{1}{z^2}(\nu ^2\frac{1}{4})+\omega ^2z^2\overline{\omega },\nu s^{}+\frac{d5}{2}.$$
(30)
The spectrum of this operator is well known and is given by
$$m_{s.b}^2=2\omega (\nu +2n+1)\overline{\omega },n=0,1,\mathrm{}.$$
(31)
From this it is seen that for leading term in decomposition (29), i.e. for $`s^{}=s`$, and for $`n=0`$ the mass spectrum is given by
$$m_{s.b}^2=(s1)/\alpha ^{}$$
(32)
and this coincides exactly with the mass spectrum of massless and massive string states belonging to leading Regge trajectory. Note that overall normalization factor as well as additive constant in r.h.s of relation (32) are result of the form of conjectured term in (27) which breaks (spontaneously) AdS symmetry. What is extremely important is that the dependence on $`s`$ in (32) is essentially fixed by specific form of AdS mass operator $`A`$ given in (25).
By now due to it is known that to construct self-consistent interaction of higher spin massless fields in $`AdS_4`$ it is necessary to introduce, among other things, a infinite chain of massless totally symmetric fields which consists of every spin just once . In higher dimensions the totally symmetric fields should be supplemented by mixed symmetry fields. What is important however in higher dimensional AdS space it is also necessary to introduce the same infinite chain of massless totally symmetric fields, i.e. the one which consists of every spin just once. It is this chain of totally symmetric fields that can be found on leading Regge trajectory of string theory.
To summarize we have demonstrated that if (spontaneously) broken hamiltonian $`P^{}`$ takes the form given in (27) then the leading components of AdS massless totally symmetric arbitrary spin $`s`$ states $`|\varphi _s`$ (29) become massive string states belonging to leading Regge trajectory. |
warning/0002/astro-ph0002468.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Near-infrared luminosity provides a good measure of a galaxy’s mass, over a wide range of Hubble types , redshifts, and star formation histories ,. With the availability of sensitive, large-format infrared array cameras on large telescopes, it is now practical to obtain infrared galaxy samples reaching below $`L^{}`$ at $`z>1`$ over areas large enough to encompass hundreds of such galaxies. In late 2001, the Space Infrared Telescope Facility (SIRTF ) will provide $`\mu `$Jy-level sensitivity in the mid-infrared, enabling rest-frame $`2\mu `$m-selected samples reaching $`L^{}`$ at $`z>3`$ to be obtained .
Kauffmann & Charlot have recently proposed that the fraction of $`z>1`$ galaxies in deep infrared $`K`$-selected samples provides a powerful means of discriminating between pure luminosity evolution (PLE) and hierarchical (CDM) scenarios for massive galaxy formation and evolution. They argue that the paucity of $`z>1`$ galaxies in the Hawaii $`K`$-band samples , already provides strong evidence against the PLE scenario and is consistent with $`\mathrm{\Omega }=1`$ CDM models. Similar arguments have been made based on the absence of red objects in surveys covering several square arcminutes to $`K22`$ , primarily the KPNO 4m Infrared Imager observations of the Hubble Deep Field (HDF IRIM) obtained by Dickinson et al.
The HDF spans a volume small enough that it would be expected to contain only a few dozen L\* galaxies with $`1<z<2`$ . In combination with the strong clustering seen in Lyman break galaxies over substantially larger fields , this suggests that global conclusions drawn from samples like the HDF should be treated with some caution.
Using the KPNO 4m, Elston, Eisenhardt, and Stanford (hereafter EES) have recently completed a substantially larger $`K`$-selected survey, whose properties are summarized and compared to the Hawaii and HDF IRIM surveys in Table 1. The EES survey is divided into 4 regions around the sky, providing some indication of field-to-field variations caused by clustering. Figure 1 shows a color–color diagram for the survey.
Table 1. IR Field Survey Characteristics.
| Survey | Bands | $`K(10\sigma )`$ | Area | N<sub>gal</sub> | N/A |
| --- | --- | --- | --- | --- | --- |
| | | mag | sq arcmin | | #/sq arcmin |
| Hawaii | $`B,I,K`$ | 19.3 | 26 | 122 | 5 |
| EES | $`B,R,I,Z,J,K`$ | 20 | 124 | 1683 | 14 |
| | | 19 | | 720 | 6 |
| HDF-IRIM | $`J,H,K`$ | 21.2 | 7 | 149 | 21 |
| | | 20 | | 76 | 11 |
| | | 19 | | 44 | 6 |
## 2 Extremely Red Objects
The surface density of ”extremely red objects” (ERO’s) in the EES sample is of interest. Graham & Dey and Cimatti define ERO’s as objects with $`RK>6`$ and show that at least one such source (Hu & Ridgway 10) is a dusty galaxy at $`z=1.44`$ with detectable submm emission, implying L$`{}_{bol}{}^{}>10^{12}`$$`\mathrm{L}_{}`$and a star formation rate of several hundred $`\mathrm{M}_{}`$per year. If the recently identified population of field sources with a surface density $`1`$ per square arcmin at comparable sub–mm fluxes are similar in nature, they would dominate the global star formation rate. SIRTF will be able to characterize this population from 3.6 to 160 $`\mu `$m with relative ease.
We find 0.7 sources per square arcminute with $`RK>6`$ and $`K<20`$. At this meeting, Barger defined very red galaxies by $`IK>4`$, finding 16 such objects to $`K=20.1`$ in a 62 square arcmin survey centered on the HDF (i.e. 0.26 per square arcmin). We find an average of 2.5 sources per square arcmin (ranging from 2.1 to 3.4 among the four EES regions) with $`IK>4`$ and $`K<20`$. The reason for this discrepancy is uncertain, but we suspect that at least part of the answer is that the HDF simply has an unusually low abundance of red galaxies.
## 3 $`JK`$ as a Lower Limit to the $`z>1`$ fraction
Next we consider a different measure of the red population: $`JK>1.9`$. Figure 1 shows that $`JK`$ is primarily sensitive to redshift, at least for earlier type galaxies. A galaxy with the spectral energy distribution of a present day elliptical galaxy would have $`JK1.9`$ at $`z=1`$. Passive or active evolution will tend to make colors bluer, so the fraction of galaxies with $`JK>1.9`$ is a reasonably reliable lower limit to the fraction of galaxies with $`z>1`$. Figure 2 demonstrates that this assertion holds up under spectrosscopic scrutiny: while there are indeed galaxies with $`JK<1.9`$ and $`z>1`$, only one object has $`JK>1.9`$ and $`z<1`$.
In Table 2 we list the lower limit to the fraction of galaxies with $`z>1`$ calculated from the $`JK>1.9`$ criterion for the EES and HDF IRIM samples, together with lower and upper limits determined from spectroscopy of the Hawaii and HDF IRIM samples, and predictions from Kauffmann and Charlot . The spectropscopic lower limits assume that none of the objects in the sample with unknown redshifts lie at $`z>1`$, while the upper limits assume that all unknown redshifts are at $`z>1`$.
While the lower limits from the $`JK`$ method for the EES sample are not as high as the fractions predicted for the PLE scenario , they are consistent with PLE, and not particularly supportive of the hierarchical model predictions. Thus we consider the PLE scenario still viable, at least based on the probable redshift distribution of faint $`K`$–selected galaxy samples. We have obtained hundreds more spectra for the EES sample with LRIS and Keck in the fall of 1998, so we expect to determine the redshift distribution of the sample with higher confidence in the near future. We also look forward to repeating the test using SIRTF to obtain a rest-frame $`K`$-selected sample out to $`z3`$.
Table 2. Fraction of $`z>1`$ Galaxies in IR Field Samples.
| $`K`$ | K&C 98 | | Hawaii | EES | HDF IRIM | |
| --- | --- | --- | --- | --- | --- | --- |
| (mag) | PLE | Hier | Spec | $`JK>1.9`$ | Spec | $`JK>1.9`$ |
| 16–18 | 28% | 0% | 2–11% | $`>15\%`$ | 14% | $`>0\%`$ |
| 18–19 | 54% | 3% | 10–17% | $`>23\%`$ | 5–50% | $`>5\%`$ |
| 19–20 | | | | $`>28\%`$ | 15-52% | $`>9\%`$ |
| 20–21 | }75% | }20% | | | 17–70% | $`>3\%`$ |
| 21–21.5 | | | | | 12–88% | $`>17\%`$ |
## 4 How Representative is the HDF?
The lower limit to the fraction of $`K`$–selected $`z>1`$ galaxies in the HDF determined via the $`JK>1.9`$ criterion is very low, and reasonably consistent with the hierarchical model predictions of (although the spectroscopically determined limits for this fraction are much less conclusive.) Red objects appear to be uncommon in and around the HDF - a fact used by Zepf and Barger to argue against a significant population of passively evolving elliptical galaxies at high redshift. This may be due to clustering effects, since early type galaxies are strongly clustered at the present epoch.
We have examined the variation in surface density of $`JK>1.9`$ and $`K<20`$ objects within the EES survey. Although the mean surface density for the EES sample is 3.4 such objects per square arcmin, this value ranges from 1.0 to 6.7 in 16 EES subfields the size of the HDF. The value in the HDF itself is 0.6, reinforcing the impression that the red population HDF is unusually sparse. Results from the HDF-South should help settle this question. The total surface density of the HDF for all colors for $`K<20`$ is also low, but within the EES range: 11.8 per square arcmin, whereas the EES subfields range from 11.8 (two subfields) to 25.8. This latter field contains a z=0.58 Rosat cluster, and the next highest density is 16.9.
The alert reader may have noticed the clump of $`z=1.27`$ redshifts in Figure 2. This cluster in the Lynx EES field was identified by Stanford et al . Thus it is fair to ask how representative is the EES sample, or at least the Lynx portion of it (which also includes the $`z=0.58`$ cluster)? The surface density of $`JK>1.9`$, $`K<20`$ objects in Lynx is 3.6 per square arcmin, and 16 per square arcmin for all $`JK`$, vs. 3.4 and 13.5 for these values respectively for the entire EES sample. The EES sample comoving volume out to $`z2`$ is a few $`10^5`$ Mpc<sup>3</sup>. If the present number density of clusters ($`10^5`$ Mpc<sup>-3</sup>) does not evolve rapidly, the presence of these clusters is not surprising.
Acknowledgements. The EES and HDF IRIM surveys would not have been possible without generous allocations from the KPNO time allocation committee. We thank the organizers for putting together an extremely stimulating and enjoyable meeting, capped by a truly spectacular conference dinner in the Chateau de Chambord. Portions of the research described in this paper were carried out by the Jet Propulsion Laboratory, California Institute of Technology, under a contract with NASA. |
warning/0002/gr-qc0002010.html | ar5iv | text | # Detection of gravitational waves from inspiraling compact binaries using a network of interferometric detectors*footnote **footnote *Based on talk given at GWDAW-99, Rome, in December 1999.
## I Introduction
Inspiraling compact binaries form prime candidates for detection by earth-based interferometric gravitational-wave (GW) detectors owing to the well understood waveform (chirp) emitted by them. Searching for chirps using a network of such detectors is gaining importance due to (a) its superior sensitivity vis a vis that of a constituent detector and (b) improving feasibility for a real-time computational search. Here, we formulate the problem of how to optimally detect the Newtonian chirp using a network of arbitrarily orientated and arbitrarily located detectors. This extends a similar study in Ref. of a network with coincident detectors.
We use the maximum likelihood method for optimizing the detection problem. A single likelihood ratio (LR) is deduced for the entire network. A super-threshold value for the maximized likelihood ratio (MLR) implies a detection. The MLR is obtained by maximizing the LR over the eight parameters that determine the Newtonian chirp: the distance $`r`$ to the binary, the inital phase $`\delta `$ of the waveform, the polarization angle $`\psi `$, the inclination angle $`ϵ`$ of the binary orbit, the time of arrival $`t_a`$ at a fiducial detector (fide), the source-direction angles $`\{\varphi ,\theta \}`$, and the chirp time $`\xi `$. In principle, this can always be done numerically using a grid in the eight dimensional parameter space. In practice, however, such a strategy is computationally unfeasible and wasteful. We show that maximization of the LR over four parameters, $`\{r,\delta ,\psi ,ϵ\}`$, can be performed analytically using the symmetries in detector responses. This allows us to scan this parameter subspace continuously. Further, the Fast Fourier Transform (FFT) can be used to maximize LR over $`t_a`$, as in the case of a single detector. Such an analytic maximization and the FFT allow us to save substantially on computational costs. Numerical maximization is required over the remaining parameters, $`\{\varphi ,\theta ,\xi \}`$, which we discuss in a future work. Here, we follow the convention laid out in Ref. .
## II The signal
There are four distinctly different reference frames of interest, associated with the source, wave, fide, and a representative detector in the network. Physical quantities in these frames are related by orthogonal transformations, $`𝒪_\mathrm{k}`$, which are defined in terms of three sets of Euler angles that specify the orientation of one frame with respect to another. Let $`𝗑`$ be an arbitrary three-dimensional real vector. Then, $`𝗑_{\mathrm{wave}}=𝒪_1(\psi ,ϵ,0)𝗑_{\mathrm{source}}`$, $`𝗑_{\mathrm{fide}}=𝒪_2(\varphi ,\theta ,0)𝗑_{\mathrm{wave}}`$, and $`𝗑_{\mathrm{detector}}=𝒪_3^1(\alpha ,\beta ,\gamma )𝗑_{\mathrm{fide}}`$, Here, the source axes have been chosen in accordance with Ref. .
The wave tensor $`w_{ij}`$ associated with any source can be expanded in terms of the STF-2 tensors $`𝒴_{2n}^{ij}`$ in an arbitrary frame as :
$$w^{ij}(t)=\sqrt{\frac{2\pi }{15}}\left[\left(h_+(t)ih_\times (t)\right)T_2{}_{}{}^{n}𝒴_{2n}^{ij}+\left(h_+(t)+ih_\times (t)\right)T_2{}_{}{}^{n}𝒴_{2n}^{ij}\right],$$
(1)
where $`i`$, $`j`$ denote spatial indices, and $`h_+`$ and $`h_\times `$ are the two GW polarizations in the transverse-traceless gauge, as measured in some given frame. The expansion coefficients $`T_{\pm 2}^n`$ are the Gel’fand functions, which depend on the Euler angles through which one must rotate that frame into the frame in which $`w^{ij}`$ is being analyzed. The above form suggests the definitions, $`e_L^{ij}=\sqrt{8\pi /15}T_2{}_{}{}^{n}𝒴_{2n}^{ij}`$ and $`e_R^{ij}=\sqrt{8\pi /15}T_2{}_{}{}^{n}𝒴_{2n}^{ij}`$, for the left- and right-circular polarization tensors, respectively. They obey, $`e_L^{ij}=e_R^{ij}`$, $`e_{L,R}^{ij}e_{L,Rij}^{}=1`$, and $`e_{L,R}^{ij}e_{R,Lij}^{}=0`$, in any frame. Thus,
$$w^{ij}(t)=\mathrm{Re}\left[\left(h_+(t)+ih_\times (t)\right)e_R^{ij}\right]2\kappa \mathrm{Re}\left[R(t)e_R^{ij}\right],$$
(2)
where $`\kappa =\sqrt{\xi }/r`$ (up to a normalization factor) and $`R(h_+(t)+ih_\times (t))/(2\kappa )`$. For a chirp, we define $`R`$ in the source frame. Then $`h_{+,\times }`$ are GW amplitudes for a face-on binary (i.e., for $`ϵ=0`$), and $`R`$ depends only on $`\{\delta ,t_a,\xi \}`$.
The response amplitude (i.e., the signal) in the $`I`$-th detector is the scalar product $`s^I=w^{ij}d_{ij}^I`$, which depends on projections of $`e_{L,R}^{ij}`$ onto the $`I`$-th detector tensor, $`d_{ij}^I`$. One such projection defines the extended beam-pattern function:
$$F^I=e_L^{ij}d_{ij}^IT_2{}_{}{}^{p}(\psi ,ϵ,0)D_p^I,p=\pm 2,$$
(3)
which corresponds to the left-circular polarization. Above,
$$D_p^I\sqrt{\frac{8\pi }{15}}T_p{}_{}{}^{n}(\varphi ,\theta ,0)d_{ij}^I𝒴_{2n}^{ij}=ig^IT_p{}_{}{}^{n}(\varphi ,\theta ,0)(T_{2n}^IT_{2n}^I),$$
(4)
where $`T_{\pm 2n}^I=T_{\pm 2n}(\alpha _I,\beta _I,\gamma _I)`$ and $`d_{ij}^I=g^I(n_{1i}^In_{1j}^In_{2i}^In_{2j}^I)`$, with $`𝐧_{1,2}^I`$ being unit vectors along the two arms of the $`I`$-th interferometer, respectively. Also, $`g^I`$ is the detector’s noise power spectral density. Then,
$$s^I(t)=2\kappa \mathrm{Re}\left(F^IR^I\right)2\kappa \mathrm{Re}\left(F^IS^Ie^{i\delta }\right)$$
(5)
where $`R^I`$ is defined via Eq. (2) and $`S^I`$ is independent of $`\delta `$.
## III The optimal network statistic
Under the Neyman-Pearson decision criterion, the optimal network statistic is the network LR, $`\lambda `$. If the noise in each detector is additive and independent of the noise in any other detector in the network, then $`\lambda `$ reduces to a product of the individual detector LR’s. Further, for Gaussian noise, the logarithmic likelihood ratio (LLR), $`\mathrm{ln}\lambda `$, simplifies to the following sum of LLR’s of $`N`$ individual detectors :
$$\mathrm{ln}\lambda =\underset{I=1}{\overset{N}{}}s_I,x_I_I\frac{1}{2}\underset{I=1}{\overset{N}{}}s_I,s_I_I=𝖻\underset{I=1}{\overset{N}{}}z_I,x_I_I\frac{1}{2}𝖻^2,$$
(6)
where $`𝖻2\kappa (_{I=1}^NF_I^2)^{1/2}`$ and $`z_I=s_I/𝖻`$. Above, $`r`$ appears only in $`𝖻`$.
Maximizing $`\mathrm{ln}\lambda `$ with respect to $`𝖻`$ and $`\delta `$ gives, $`\mathrm{ln}\lambda |_{\widehat{𝖻},\widehat{\delta }}=\left|_{I=1}^NQ_IC_I^{}\right|^2/2`$, where $`Q_I2\kappa F_I/𝖻`$ and $`C_I^{}S_I,x_I_I`$. This shows that the network vector $`𝖲`$, with $`S_I`$’s as its components, is the matched network-filter. Also, $`\mathrm{ln}\lambda |_{\widehat{𝖻},\widehat{\delta }}`$ is a function of six parameters, namely, $`\{\psi ,ϵ,t_a,\varphi ,\theta ,\xi \}`$. To extend these results to the case where the detectors are arbitrarily located, note that the dependence of $`\mathrm{ln}\lambda |_{\widehat{𝖻},\widehat{\delta }}`$ on $`\{\psi ,ϵ\}`$ can be isolated. This is because the network vector $`𝖰`$, with $`Q^I`$’s as its components, is:
$$𝖰=𝖥^1\left(T_2^2(\psi ,ϵ,0)𝖣_2+T_2^2(\psi ,ϵ,0)𝖣_2\right)Q^2\widehat{𝖣}_2+Q^2\widehat{𝖣}_2,$$
(7)
where $`𝖣_p`$ define network vectors with $`D_p^I`$ as their components; $`\widehat{𝖣}_p`$ are their normalized counterparts. Thus, $`Q_2=\widehat{𝖣}_2𝖰=Q^{+2}+Q^2\widehat{𝖣}_2\widehat{𝖣}_2`$. Hence, $`\{\widehat{𝖣}_2,\widehat{𝖣}_2\}`$ define a two-dimensional complex plane, $`𝒫`$ (the helicity space, a subspace of $`𝒞^N`$), on which a metric $`g_{pq}`$ can be defined. Then, $`Q_p=g_{pq}Q^q`$ with $`p,q=\pm 2`$. The $`N`$-dimensional correlation vector $`𝖢`$, in general, lies outside $`𝒫`$. However, $`𝖰`$ lies totally in $`𝒫`$. Thus, the statistic reduces to
$$\mathrm{ln}\lambda |_{\widehat{𝖻},\widehat{\delta }}=\left|𝖰𝖢^{}\right|^2/2=\left|𝖰𝖢_𝒫^{}\right|^2/2,$$
(8)
where $`𝖢_𝒫`$ is the projection of $`𝖢`$ onto the helicity space.
Maximization of $`\mathrm{ln}\lambda |_{\widehat{𝖻},\widehat{\delta }}`$ over $`\{\psi ,ϵ\}`$ is achieved by aligning $`𝖰`$ along $`𝖢_𝒫`$. This requires that $`Q^{\pm 2}=C_𝒫^{\pm 2}/𝖢_𝒫_𝒫`$, which implies: $`Q^{+2}/Q^2=C_𝒫^{+2}/C_𝒫^2`$. Since the RHS above can take any value in the complex plane, we need to prove that $`Q^{+2}/Q^2`$ obeys the same property. To this end, note that,
$$\frac{Q^{+2}}{Q^2}=\frac{T_2^2(\psi ,ϵ,0)}{T_{+2}^2(\psi ,ϵ,0)}=\left(\frac{1\mathrm{cos}ϵ}{1+\mathrm{cos}ϵ}\right)^2\mathrm{exp}(4i\psi ),$$
(9)
which can indeed attain any value on the Argand plane. The values of $`\psi `$ and $`ϵ`$ that maximize the statistic are, $`\widehat{\psi }=\mathrm{arg}(x)/4`$ and $`\widehat{ϵ}=\mathrm{cos}^1[(1\sqrt{x})/(1+\sqrt{x})]`$ where $`xC_𝒫^{+2}/C_𝒫^2`$. Thus, the statistic maximized over these four parameters is,
$$\mathrm{ln}\lambda |_{\widehat{𝖻},\widehat{\delta },\widehat{\psi },\widehat{ϵ}}=𝖢_𝒫^2/2.$$
(10)
Let $`\widehat{V}^\pm `$ denote a pair of orthonormal, complex basis vectors on $`𝒫`$. Then,
$$𝖢_𝒫^2=C^+^2+C^{}^2=(c_0^+)^2+(c_{\pi /2}^+)^2+(c_0^{})^2+(c_{\pi /2}^{})^2,$$
(11)
where $`C^\pm =𝖢_𝒫\widehat{V}^\pm =c_0^\pm +ic_{\pi /2}^\pm `$. It can be verified that the network statistic is, therefore, a sum of the squares of four Gaussian random variables with constant variance. This simplifies the computation of thresholds and detection probabilities.
A network filter is constructed as follows: For a given $`\xi `$ compute the Newtonian chirp for $`t_a=0`$. Then, for a given direction $`\{\varphi ,\theta \}`$, use the appropriate time-delays with respect to fide to time-displace the chirp at each detector. This collection of time-displaced chirps constitute the network filter. Also, $`\widehat{t}_a`$ is obtained by shifting the network filter ‘rigidly’ on the time axis, which can be done efficiently using FFT. The bank of filters on $`\{\varphi ,\theta ,\xi \}`$ can be obtained by correlating two neighboring normalized filters in the usual way. This work is now in progress.
## IV Conclusion
We have given here a formulation to optimally detect the Newtonian chirp with a network of detectors in which the noise is additive, Gaussian, and uncorrelated between detectors. We have shown how this can be done efficiently by analytically maximizing the LR over four parameters and using FFT to maximize over the time-of-arrival. In a future work we hope to address key issues such as required computational power for such a search and also estimate errors in parameter values.
###### Acknowledgements.
SB acknowledges support from PPARC grant PPA/G/S/1997/00276. AP acknowledges support from a CSIR grant and from AEI, Potsdam, for a 3-month visit. |
warning/0002/astro-ph0002107.html | ar5iv | text | # Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum
## 1 Introduction
The intrinsic spectral energy distribution (SED) of active galactic nuclei (AGN), which extends from the radio up to $`\gamma `$-rays, cannot be directly observed from the Lyman limit and up to several hundred eV due to Galactic and intrinsic absorption. However, the extreme-UV (EUV) and soft X-ray continuum can be investigated indirectly by the infrared coronal line emission. These lines are emitted by collisionally excited forbidden fine-structure transitions of highly ionized atoms, whose ionization potentials extend well beyond the Lyman limit up to hundreds of eV. Unlike the strong permitted lines of these ions, which are also emitted in the obscured EUV, the reddening-insensitive forbidden IR coronal lines and semi-forbidden optical coronal lines can be observed. Therefore, when photoionization is the main ionization mechanism, the coronal lines can provide information on the intrinsic obscured SED and the accretion process that powers the AGN. This information can be extracted by photoionization models of the NLR.
The coronal lines are collisionally suppressed in the dense broad line region (BLR) close to the continuum source and are efficiently emitted only from the more rarefied gas in the narrow line region (NLR), hundreds of pc away from the center. It is well established that large quantities of gas attenuate the continuum emission in many AGN. These gas clouds are detected by narrow UV absorption lines (e.g. Kriss et al. 1992a ) or by X-ray absorption features and emission lines (e.g. George et al. George98 (1998)). Although their exact location along the line of sight is unknown, there are reasons to believe that in some cases they may be inside the NLR. In particular, the warm absorbers that block the X-ray continuum appear to cover a large fraction of the continuum source (George et al. George98 (1998)). This raises the possibility that in some AGN the ionizing SED, which is traced by the coronal lines, is not the intrinsic one produced by the accretion process, but rather one that is filtered by intervening absorbers inside the NLR. It has been proposed that such absorbers are common in Seyfert galaxies, and are responsible for the observed correlations between the soft X-ray slope and the narrow emission line spectra of Seyfert 1.5 galaxies (Kraemer, Ruiz & Crenshaw KTCG99 (1999)).
This study of the Seyfert 2 galaxy $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ is part of the *ISO-SWS* program on bright galactic nuclei. Previous studies in this program include the reconstruction of the SED of the Seyfert 2 Circinus galaxy (Moorwood et al. Moorwood96 (1996); Alexander et al. Alexander99 (1999)) and of the Seyfert 1 Galaxy NGC 4151 (Alexander et al. Alexander99 (1999)). In both cases we found evidence of a “Big Blue Bump” signature of a thin accretion disk (Shakura & Sunyaev SS73 (1973)). However, in the case of NGC 4151 this structure is masked by a deep absorption trough of an absorber situated between the BLR and the NLR, which filters the light that photoionizes the NLR.
In this paper we apply our SED reconstruction method to $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$, one of the closest, brightest and most extensively studied Seyfert 2 galaxies, which is considered a prototype of this AGN class. The first detection of broad permitted emission lines in the polarized light of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ (Antonucci & Miller AM85 (1985)) provided a major argument for the Seyfert 1 and 2 unification scheme (Antonucci Antonucci93 (1993)). This scheme postulates that the two Seyfert types have both broad and narrow line regions and an obscuring torus that lies between the two. When the torus is face on and the BLR is directly observed, the AGN is classified as a Seyfert 1. When the BLR is obscured by the torus, the AGN is classified as a Seyfert 2, and the BLR can be observed only indirectly in scattered polarized light. The factors that determine the accretion properties and the ionizing SED of AGN are currently unknown. However, the Seyfert unification picture implies that all possibly relevant factors being equal, such as luminosity, host galaxy type or redshift, the intrinsic SED of both AGN types should be similar. It is therefore of interest to complement our previous study of the nearby luminous Seyfert 1 galaxy NGC 4151 with a corresponding study of a nearby luminous Seyfert 2 galaxy with a similar host galaxy type, such as $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$. The *ISO-SWS* observations of *$`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$* are presented in a companion paper (Lutz et al. Lutz00 (2000)) and are used there to derive the gas density and to place constraints on the structure and dynamics of the NLR.
This paper is organized as follows. In §2 we summarize the physical properties of the nucleus of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ that are needed for constructing the photoionization models and constraining their results. In §3 we present the emission line flux compilation that we use in our modeling. In §4 we briefly discuss the construction and fitting of the NLR photoionization models. We present the results in §5 and discuss them in §6.
## 2 The physical properties of NGC 1068
$`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ is a barred spiral galaxy at $`z=0.0036`$ (distance $`D=16.6`$Mpc for $`H_0=65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$) with magnitude $`m_B=9.17`$ (e.g. Lipovetsky, Neizvestny & Neizvestnaya LNN88 (1988)). Observations of the nucleus of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ and models of the nuclear line emission indicate that the gas in the nucleus forms a complex system, which is composed of various spatial and dynamical components. These are excited by several physical mechanisms, including photoionization by the nuclear continuum, photoionization by hot stars, and possibly also by shocks and energetic particles from a radio jet. In order to isolate the effects of the nuclear continuum and to construct photoionization models of the NLR it is necessary to understand the morphology and content of the galactic nucleus. We present here a brief overview of the properties of the nucleus that are relevant to this work.
### 2.1 The galactic nucleus
The most prominent morphological feature in the nucleus of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ is the asymmetric bi-polar pattern of both the radio and the optical line emission. The radio emission extends over $`15\mathrm{}`$ and has a sharply defined northern lobe and a weaker diffuse southern lobe (Wilson & Ulvestad WU83 (1983); Muxlow et al. Muxlow96 (1996)). The southern lobe is both smaller and redder, which is consistent with the picture that the large northern lobe is observed above the galactic disk, generally facing the observer, and the southern lobe is seen through the galactic disk (Unger et al. ULPA92 (1992); Macchetto et al. Macchetto94 (1994)). Images of the NLR in low excitation lines show mainly the northern cone (Cecil, Bland & Tully CBT90 (1978); Unger et al. ULPA92 (1992); Evans et al. Evans91 (1991)). The precise value of the opening angle associated with the emission maps depends on the way the edge of the cone is defined and on the assumed position of the nucleus. The location of the nucleus can be determined to within $`0.05\mathrm{}`$ by the center of symmetry of the UV polarization pattern (Capetti et al. Capetti95 (1995), Kishimoto Kishimoto99 (1999)). The position of the nucleus does not appear to coincide with the maximum of the continuum emission, which implies that the nucleus is heavily obscured even in the infrared. This is consistent with an obscuring torus of column density in excess of $`10^{24}`$ cm<sup>-2</sup>, as is inferred from X-ray (Marshall et al. Marshal93 (1993)) and CO observations (Tacconi et al. Tacconi94 (1994)). The positioning of the nucleus makes it possible to estimate the opening angle of the radiation cone at $`70\mathrm{°}`$ up to $`100\mathrm{°}`$. The ionization cone appears to be only partially filled.
Marconi et al. (MWMO96 (1996)) find that the coronal line emission peaks $`0.5\mathrm{}`$ NE of the nucleus and extends up to $`4\mathrm{}`$. They note that the blueshift of the emission line profile centroid increases, and the FWHM of the profile decreases with $`E_{\mathrm{ion}}`$ (the ionization energy required to produce the emitting ion from the preceding ionization stage). They interpret this as evidence that the emission lines are emitted from outflowing material. The high ionization lines are emitted in the inner light cone, where the velocity field is relatively coherent, while the lower ionization lines are emitted from slower, more extended areas with different velocities. The NLR $`[\mathrm{O}\text{iii]}\lambda 5007`$ emission extends over the few inner arcseconds (Evans et al. Evans91 (1991); Unger et al. ULPA92 (1992); Dietrich & Wagner DW98 (1998)). The extended emission line region (EELR) $`[\mathrm{O}\text{iii]}\lambda 5007`$ emission extends over more than $`10\mathrm{}`$ (Unger et al. ULPA92 (1992)).
The unresolved BLR is seen only in scattered polarized light and has FWHM of $`3000`$ km s<sup>-1</sup> (Miller et al MGM91 (1990)). As is seen in other Seyfert 2 galaxies (Wilson Wilson88 (1988); Capetti et al. Capetti96 (1996)), the morphology of the NLR emission maps is correlated with that of the radio structure (Wilson & Ulvestad WU83 (1983); Capetti et al. CAM97 (1997)). In particular, a 3-dimensional reconstruction of the positions of individual clouds based on polarimetric measurements also indicates that the northern NLR cone is directed towards the observer, and the southern part away from the observer (Kishimoto Kishimoto99 (1999)). The close correspondence between the NLR and the jet suggests that the jet outflow sweeps and compresses the ambient gas and thereby increases the line emissivity. The NLR has a very complex structure (Macchetto et al. Macchetto94 (1994)) and displays large scale clumpiness. It is composed of many line emitting cloud complexes (Alloin et al. Alloin83 (1983); Meaburn & Pedlar MP1986 (1986); Evans et al. Evans91 (1991); Dietrich & Wagner DW98 (1998)). The individual components have FWHM ranging from $`200`$ km s<sup>-1</sup> to $`700`$ km s<sup>-1</sup>, and extend over $`2500`$ km s<sup>-1</sup> in velocity space, resulting in an integrated $`[\mathrm{O}\text{iii]}\lambda 5007`$ profile with FWHM of $`1150`$ km s<sup>-1</sup> (Dietrich & Wagner DW98 (1998)). The velocity field of the NLR clouds with the lowest FWHM is consistent with rotation around the nucleus. The bulk velocities of intermediate FWHM clouds are clearly split relative to the symmetry axis of the radio jet, which suggests that they are associated with the interaction between the radio jet and the NLR gas. The highest FWHM clouds are associated with highly polarized emitting structures (Capetti et al Capetti95 (1995)), and could therefore be a reflected image of an inner obscured region. The EELR appears to follow the galactic rotation (Unger et al. ULPA92 (1992)).
The large *ISO* apertures ($`14\mathrm{}\times 20\mathrm{}`$ to $`20\mathrm{}\times 33\mathrm{}`$) includes both the outflowing inner NLR and the rotating EELR.
### 2.2 The line emitting gas
There is evidence that at least three different ionization mechanisms are at work in the NLR. The very high ionization states are probably due to the the central continuum source. Marconi et al. (MWMO96 (1996)) find that the infrared coronal line ratios point to photoionization as the main excitation mechanism of the coronal gas, and that both collisional excitation or photoionization by very hot stars can be ruled out. At lower ionization states, the jet / ISM interaction can provide internal sources of excitation in addition to the external central continuum, for example by fast shocks or cosmic rays. The morphological connection between the radio and line emission suggests that the jet outflow shapes the NLR. Estimates of high gas temperatures (Kriss et al. 1992b ) and the existence of dense but highly ionized clouds near the nucleus on both sides of the jet axis suggest that the jet may also play a role in photoionizing the clouds (Capetti et al CAM97 (1997); Axon et al. Axon98 (1998); Dietrich & Wagner DW98 (1998)). Hot stars are a third ionization mechanism. Unresolved UV continuum point sources in the inner $`7\mathrm{}\times 7\mathrm{}`$, which are not observed in $`[\mathrm{O}\text{iii]}\lambda 5007`$, could be OB associations (Macchetto et al. Macchetto94 (1994)). Hot stars are certainly a component in the ring-like structure that surrounds the nucleus. An ellipse of H ii emission delineates the NLR at an average radius of $`13\mathrm{}`$ (Cecil, Bland & Tully CBT90 (1978); Bruhweiler et al. BTA91 (1991)) and starburst knots and CO emission encircle the nucleus and the NLR at an average radius of $`18\mathrm{}`$ (Planesas, Scoville & Myers Planesas91 (1991)). However, the overall similarity in the profiles of the mid-infrared high-excitation lines, which cannot be excited by hot stars, and the profiles of the intermediate excitation lines ($`E_{\mathrm{ion}}>30\mathrm{eV}`$), which could be excited by hot stars, strongly suggests that gas excited by hot stars within the large *ISO* aperture does not contribute more than $`20\%`$ to these NLR line fluxes (Lutz et al. Lutz00 (2000)). An additional complication due to stars is contamination of the observed emission line spectrum with stellar absorption features. This appears in the difference spectrum of the $`30\mathrm{}`$ and $`18\mathrm{}`$ *HUT* apertures, which shows both a reddened early-type stellar continuum and stellar absorption features (Kriss et al. 1992b ).
As is discussed by Lutz et al. (Lutz00 (2000)), the *ISO-SWS* line ratios indicate that the mid-IR lines are emitted from gas with a hydrogen density of $`2000\mathrm{cm}^3`$. The NLR appears to contain also higher density gas. The density of individual knots in the high ionization core is estimated at $`10^4`$ to $`3\times 10^4`$ cm<sup>-3</sup> from *HST* measurements of the $`[\text{S}\text{ii}]`$ doublet, while the overall density $`1\mathrm{}`$ from the nucleus varies between $`10^3`$ to $`4\times 10^3`$ cm<sup>-3</sup> (Capetti et al. CAM97 (1997)).
The ionization parameter<sup>2</sup><sup>2</sup>2$`UQ_{\mathrm{ion}}/4\pi r^2nc,`$ where $`Q_{\mathrm{ion}}`$ is the ionizing photon luminosity, $`n`$ is the hydrogen density at the illuminated face of the cloud, $`r`$ is the distance of the face of the cloud from the continuum source and $`c`$ is the speed of light. The local ionization parameter in the cloud falls with increasing depth due to absorption and geometrical dilution of the radiation field. also appears to vary across the NLR. The $`[\mathrm{O}\text{iii]}\lambda 5007/(\mathrm{H}\alpha +[\mathrm{N}\text{ii]}\lambda 6584)`$ ratio (Capetti et al. CAM97 (1997)) traces a high excitation core ($`\mathrm{log}U2.5`$) in the inner $`1\mathrm{}\times 2\mathrm{}`$ north of the nucleus, followed by a lowered ionization halo ($`\mathrm{log}U3.3`$) out to $`4\mathrm{}`$ from the nucleus, and then intermediate ionization filaments ($`\mathrm{log}U2.8`$) which extend out to the EELR.
Estimates of dust reddening in the NLR of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ range from $`E_{\text{b-v}}=0.07`$ for the continuum (Kriss et al 1992b ), to 0.20 (Marconi et al MWMO96 (1996)), 0.40 (Shields & Oke SO75 (1975); Neugebauer et al Neugebauer80 (1980); Ward Ward87 (1987)), and $`E_{\text{b-v}}=0.52`$ (Koski Koski78 (1978)) for the NLR. The analysis of Kraemer et al. (KRC98 (1998)) indicates that there may be some dust mixed with gas in varying amounts. The observed line ratios in the NLR suggest that the O/N abundance is less than solar. Netzer (Netzer97 (1997)) and Netzer & Turner (NT97 (1997)) interpret this as an indication of under-abundant oxygen (see also Sternberg, Genzel & Tacconi SGT94 (1994)), and propose that the abundances of He:C:N:O:Ne:Mg:Si:S:Ar:Fe relative to hydrogen are $`(100:3.7:1.1:2.7:1.1:0.37:0.35:0.16:0.037:0.4)\times 10^4`$, respectively. Kraemer et al. (KRC98 (1998)) interpret the line ratios as showing an overabundance of nitrogen, as well as hinting at higher than solar iron and neon, and propose abundance ratios of $`(100:3.4:3.6:6.8:2.2:0.33:0.31:0.15:0.037:0.8)\times 10^4`$.
Finally, Kraemer et al. (KRC98 (1998)) model the NLR emission with a multi-component gas model, and suggest that the low ionization emission lines are emitted by gas that is partially screened by a dense, optically thin component. They use their models to estimate that the filling factor is $`F10^4`$.
### 2.3 The ionizing continuum
Unlike the situation in Seyfert 1 galaxies, where it is possible to observe the intrinsic SED outside the obscured range, the intrinsic SED of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ cannot be directly observed even in the optical or X-ray bands. Only a small fraction of the AGN light is scattered into the line of sight, and can be observed against the host galaxy in polarized light. Pier et al (Pier94 (1994)) compiled various continuum measurements in the optical, UV and X-ray, and carefully took account of aperture differences, star-light contamination, reflection by dust and bremsstrahlung emission from the scattering plasma. They conclude that the resulting reflected SED is broadly similar to that observed in Seyfert 1 galaxies. Pier et al (Pier94 (1994)) also list various estimates of the fraction of light reflected by the scatterer. These values range from $`f_{\text{refl}}=10^3`$ (Bland-Hawthorn & Voit BH93 (1993)) to $`0.05`$ (Bland-Hawthorn, Sokoloski & Cecil BSC91 (1991)). Pier et al (Pier94 (1994)) argue that the most reliable estimate is $`f_{\text{refl}}0.01`$ to within a factor of a few. Because the reflectors are more than a hundred light years away from the continuum source (Miller et al. MGM91 (1990)), short-term continuum variability is unlikely to affect the reconstruction of the reflected SED.
We adopt the Pier et al. nuclear continuum SED in the UV and X-ray as the template SED (Fig. Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum), and further extend it from 10 keV to 100 keV with a slope of $`F_\nu \nu ^1`$. We enumerate on the unobserved UV to soft X-ray range to find the best fitting SED.
## 3 The observed line flux compilation
In addition to the ISO-SWS mid-IR lines fluxes presented in Lutz et al. (Lutz00 (2000)), we compiled a list of UV to IR lines from the literature. The compilation initially included a list of $`120`$ measured line fluxes, which were obtained over the last $`30`$ years using various instruments with different apertures, spectral resolutions and reduction techniques. The observed lines were taken from the following references, listed by spectral band with the aperture used (where given): UV lines from Kriss et al (1992b ) ($`18\mathrm{}`$); optical lines from Osterbrock & Parker (OP65 (1965)) (trailing long slit), Anderson (Anderson70 (1970)) ($`8\mathrm{}\times 8\mathrm{}`$), Wampler (Wampler71 (1971)) ($`10\mathrm{}`$), Koski (Koski78 (1978)) ($`2.7\mathrm{}\times 3.4\mathrm{}`$) and Neugebauer et al (Neugebauer80 (1980)) ($`10\mathrm{}\times 20\mathrm{}`$); optical to near-IR lines from Shields & Oke (SO75 (1975)) ($`10\mathrm{}\times 10\mathrm{}`$), near-IR lines from Oliva & Moorwood (OM90 (1990)) ($`6\mathrm{}\times 6\mathrm{}`$), Marconi et al (MWMO96 (1996)) (long slit of width $`4.4\mathrm{}`$), Osterbrock & Fulbright (OF96 (1996)) ($`3\mathrm{}\times 18\mathrm{}`$) and Osterbrock, Tran & Veilleux (OTV92 (1992)) (long slit of width $`1.2\mathrm{}`$); IR lines from Thompson (Thompson96 (1996)) ($`2\mathrm{}\times 10\mathrm{}`$) and the ISO-SWS ($`14\mathrm{}\times 20\mathrm{}`$ to $`20\mathrm{}\times 33\mathrm{}`$).
As is discussed in detail by Alexander et al. (Alexander99 (1999)), the result of such a compilation is generally not self-consistent, and only a small subset of the of the $`120`$ lines can be used. First, we excluded low and medium excitation lines ($`E_{\mathrm{ion}}100`$eV) observed through apertures whose smaller dimension is less than $`3\mathrm{}`$, so as to avoid significant loss of light due to incomplete coverage of the line emitting region. The observed emission from the high excitation lines is centrally concentrated in the inner $`<4\mathrm{}`$ (Marconi et al MWMO96 (1996)), which are covered even by the smallest apertures used. Second, we excluded lines with $`E_{\mathrm{ion}}<30`$eV to avoid using lines that may be photoionized primarily by young hot stars or other non-AGN, lower-energy excitation processes. This also reduces the loss of light bias. Third, we excluded narrow lines whose measured flux is uncertain, either because the flux is very low (flux less than $`5\times 10^{13}`$erg s<sup>-1</sup> cm<sup>-2</sup>, $`2\%`$ of the strongest line), or because it has a significant broad component, such as the $`\mathrm{N}\text{v}\lambda 1240`$ and $`\mathrm{C}\text{iv}\lambda 1549`$ lines. Fourth, we excluded the $`[\mathrm{Fe}\text{x]}\lambda 6734`$ and $`[\mathrm{Fe}\text{xi]}\lambda 7892`$ lines, whose collision strengths are highly uncertain and therefore cannot be modeled reliably. The final, much reduced line list includes 22 lines (Table 1), which we use for obtaining the best-fit SED. Whenever more than one measurement of the line exists, we quote the average flux and use the rms scatter as an error estimate. Important IR lines that were not included in the final list were nevertheless compared to the best fit model predictions to verify that there are no gross inconsistencies.
## 4 The photoionization models
The method of constructing photoionization models for the NLR and, in particular, the “$`\mathrm{log}^2S`$ fit procedure” for obtaining the best fitting SED and gas parameters is described in detail in Alexander et al. (Alexander99 (1999)), and is summarized here briefly. We parameterize the SED as a piece-wise broken power-law (Fig. Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum), and enumerate on the different possibilities of connecting the power-law segments. We test a large number of simplified NLR gas models, which consist of optically thick (radiation bounded) clouds whose ionized surfaces partially cover a spherical shell around the continuum source. Each cloud extends in the radial direction as far as it takes to effectively absorb all the ionizing UV photons (specifically, until the hydrogen ionization fraction falls below 2%). The gas clouds are parameterized by their chemical composition, the hydrogen density $`n`$, the ionization parameter $`U`$ at the irradiated face of the cloud, and the filling factor $`F`$. In addition, an asymmetry parameter $`A`$ expresses the ratio between the ionizing flux directed towards the NLR and that directed towards the observer. This describes situations where the continuum source is not isotropic, or where only a fraction $`f_{\mathrm{refl}}`$ of the continuum is reflected towards the observer ($`A=1/f_{\mathrm{refl}}`$).
For each NLR gas model, the fit procedure uses the observed line fluxes to derive the best-fit SED for that gas model, and from it to derive in a self-consistent way the corresponding covering factor $`C,`$ the inner NLR angular radius $`\theta `$, the width of the ionized region $`\mathrm{\Delta }\theta _{\mathrm{ion}}`$ (defined here as the radial extent of the Balmer lines emitting gas) and the reddening coefficient $`E_{\text{b-v}}`$. These parameters are constrained by the observations and so can be used to limit the range of acceptable NLR gas models. The $`\mathrm{log}^2S`$ fit procedure assigns a score $`S`$ to the best-fit model, which means that the model line fluxes fit the observed ones up to a factor $`S`$, on average. The worst-fitting line and the factor by which it deviates from the observed value are also recorded. Monte-Carlo simulations are used to calculate confidence limits on the best-fit SED. In addition, we calculate the correlation between the model-to-data line ratios and the lines wavelengths, ionization potentials, critical densities and “deplitivity” (the tendency of an element to be depleted into dust). These residual correlations test whether the remaining inconsistencies in the best-fit model are related to inaccurate modeling of the reddening, the spectral hardness / softness of the continuum, the gas density or its dust content. The correct model should not display any such correlations. The final, global best-fit SED is the one with the best $`S`$-score among all the NLR gas models that are consistent with the observed NLR geometry and reddening, whose worst-fitting line is not too far from the observed value, and which display no significant residual correlations.
We assume in all models that the gas clouds have constant density and that $`f_{\mathrm{refl}}=0.005`$. We investigate two classes of models. The first consists of models with a single type of cloud. We enumerate on different values for the ionization parameter ($`\mathrm{log}U=1,2`$ and $`3`$), gas density ($`n=2000`$ and $`10^4\mathrm{cm}^3`$) and filling factor ($`\mathrm{log}F=2,3`$ and $`4`$ at the ionized surface). We test two different sets of non-solar abundances, the low oxygen set and high nitrogen set (§2.2). We test two different possibilities for the radial run of the filling factor. The first is that $`F`$ is constant, which corresponds to either a static distribution of clouds, a rotating distribution of clouds, a linear constant velocity outflow, a strongly decelerating outflow at a constant opening angle, or an outflow where clouds are continuously added to the flow (e.g. from the molecular torus). The second is $`Fr^2`$, which corresponds to a constant velocity and constant opening angle outflow where the clouds are formed at the base of the flow. The second class of models have two types of gas clouds, and are constructed by combining all the possible pairs of one component models. We assume that the clouds do not obscure one another.
The photoionization calculations were carried out using the numerical photoionization code ion99, the 1999 version of the code ion described by Netzer (Netzer96 (1996)).
## 5 Results
### 5.1 Single component models
We find that single component models generally fail to fit the observed line ratios and the observed constraints on the geometry of the NLR. Low $`U`$ models ($`\mathrm{log}U=3`$) cannot reproduce the high excitation lines regardless of the values of the other model parameters. For the $`Fr^2`$ models, the low $`U`$ models significantly over-estimate the size of the NLR, while the high-$`U`$ models require unphysical covering factors exceeding unity. Single component models with a high $`U`$ and low constant filling factor do somewhat better, although such models are problematic because a constant filling factor is inconsistent with simple scenarios of NLR outflow. The best fit single component model has low oxygen abundance with $`\mathrm{log}U=1`$, $`\mathrm{log}n=3.3`$ and $`\mathrm{log}F=3`$. This model can reproduce the observed lines up to a factor of 2 and predicts $`\theta =0.8\mathrm{}`$, $`\mathrm{\Delta }\theta _{\mathrm{ion}}=4.8\mathrm{}`$, $`C=0.35`$ and $`E_{\text{b-v}}=0.18.`$ These values are roughly consistent with the observed constraints. However, this model under-predicts the observed $`\mathrm{O}\text{vi}\lambda 1035`$ flux by a factor of $`5`$ and shows a residual negative correlation with $`\lambda _0`$ at the 5% confidence level. Like all the models investigated here, the best fit SED displays a deep trough at 4 Ryd ($`\mathrm{log}f=27.4,29.0,27.4,28.2`$ at 2, 4, 8 and 16 Ryd, respectively). A fit of similar quality is obtained with the high nitrogen abundance set.
### 5.2 Two component models
A better fit to the observations is provided by the best-fit two component model (Table 2). This model fits the 22 observed line fluxes to within a factor of 1.9 on average (experience shows this is as well as one can expect for AGN photoionization models). The worst fitting line, $`[\mathrm{Ar}\text{vi]}\lambda 4.5\mu \mathrm{m}`$, is under-predicted by a factor of 4. The model-to-data ratios of individual lines are displayed in Fig. Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum. The low excitation IR lines ($`E_{\mathrm{ion}}<30\mathrm{eV}`$) such as $`[\mathrm{Ne}\text{ii]}\lambda 12.8\mu \mathrm{m}`$, $`[\mathrm{S}\text{iii]}\lambda 18.7\mu \mathrm{m}`$ and $`[\mathrm{S}\text{iii]}\lambda 33.5\mu \mathrm{m}`$, which were not used in the fit, are nevertheless all consistent with the observations to within a factor of 2. However, the model-to-data line ratios for these lines are not much smaller than 1, as would be expected if there is a significant contamination from gas photoionized by starbursts. The high excitation IR lines such as $`[\mathrm{S}\text{ix]}\lambda 1.25\mu \mathrm{m}`$, $`[\mathrm{Si}\text{x]}\lambda 1.4\mu \mathrm{m}`$ and $`[\mathrm{Si}\text{ix]}\lambda 2.6\mu \mathrm{m}`$, which were not used in the fit, are also consistent with the observations to within a factor of 3. There are no statistically significant residual correlations between these ratios and $`\lambda _0`$, $`E_{\mathrm{ion}}`$, the deplitivity or $`n_c`$. The best fit values of the extinction $`(E_{\text{b-v}}0.2`$) is in agreement with other estimates (§2.2). The covering factors are somewhat larger than is indicated by the observed opening angle of the emission cone (§2.1). The dense compact component requires a covering factor of $`C=0.45`$, which corresponds to a bi-cone with an opening angle of $`115\mathrm{°}`$. The less dense, more extended component requires a covering factor $`C=0.26`$, which corresponds to a single cone with an opening angle of $`130\mathrm{°}`$ or a bi-cone with opening angle of $`90\mathrm{°}`$. Figure Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum shows the best fit SED for this gas model together with the 99.9% confidence limits on it (the lower confidence limits at 2, 4 and 16 Ryd were not calculated as they extend beyond the SED grid). This SED shows the generic trough that is seen in the best-fit SED of all the models we investigated. This model has a less than solar oxygen abundance (Netzer Netzer97 (1997); Netzer & Turner NT97 (1997)). We find only a few two-component models with high nitrogen abundance that come close to a reasonable fit to the observations. Of these, almost all display a negative correlation with deplitivity at the 5% to 10% significance, which may indicate that the line emitting gas is more depleted than is assumed by the high nitrogen abundance set.
We conclude that best-fit two component model provides a better, but not an overwhelmingly better fit to the observed line ratios than the single component model. Its marked advantage over the single component models lies in its consistency with the observed NLR geometry and kinematics. This is further discussed in the next section.
## 6 Discussion
As was discussed in detail by Alexander et al. (Alexander99 (1999)), there are various degeneracies between the parameters that describe the gas model ($`n,`$ $`A`$, $`F`$, $`U`$, $`C`$, $`\theta `$, and $`\mathrm{\Delta }\theta _{\mathrm{ion}}`$). These degeneracies allow the fit procedure to converge to a robust best-fit SED even when the assumed values of $`n`$, $`A`$, $`F`$ or $`U`$ significantly differ from their true values, since this can be compensated to a large degree by a suitable modification of the gas geometry ($`C`$, $`\theta `$, and $`\mathrm{\Delta }\theta _{\mathrm{ion}}`$). This property of the fit procedure is especially important in the analysis of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$, where observations indicate that the actual properties of the NLR gas are much more complex than can be modeled by our family of simplified gas models. For this reason, we place more weight on the fact that all the best-fit SED models, whether one or two-component, display a deep trough at 4 Ryd than on the determination of the exact values of the gas parameters.
The trough in the SED is required for reproducing the relative line fluxes of the high and low ionization species. To check the robustness of this result, we attempted to re-fit the observed line fluxes with all of our one and two-component models using an approximate single power-law SED ($`\mathrm{log}f=25.8,26.6,27.4,28.2`$ at 2, 4, 8 and 16 Ryd, respectively), which was held fixed in the fit procedure. In all cases the low excitation lines were over-estimated with respect to the high excitation lines, regardless of the values of $`U`$, $`n`$ or $`F`$. For example, when the power-law SED is applied to the best-fit two-component gas model (Table 2), the low excitation lines ($`E_{\mathrm{ion}}50\mathrm{eV}`$) are over-estimated by the model by up to a factor of 13, whereas the high excitation lines ($`E_{\mathrm{ion}}100\mathrm{eV}`$) are under-estimated by up to a factor of 23. The overall mismatch of this SED with the observed line fluxes is reflected in both the poor fit score of $`S=4.1`$ and in the very strong residual anti-correlation between the line ratios and $`E_{\mathrm{ion}}`$, whose random probability is $`10^4`$. We caution against generalizing this result to mean that every AGN that exhibits a hard emission line spectrum has an absorbed ionizing SED. The emission line spectrum reflects the gas parameters, such as $`U`$, $`n`$ and $`F`$, no less than it does the ionizing SED. It is necessary to have some knowledge of the likely range of values for these parameters in order to interpret the hardness of the line spectrum. Alexander et al. (Alexander99 (1999)) provide a counter-example where subtle cancellations between the SED and the gas properties lead to a situation where an AGN (NGC4151) with a hard absorbed SED has a *softer* emission line spectrum than another AGN (Circinus) with a soft unabsorbed Big Blue Bump.
Although we do not claim to fix the gas parameters with certainty, the best fit model (Table 2) is reassuringly consistent with the observations, which broadly indicate that the integrated NLR emission originates in two components. Component A of the model can be interpreted as a system of dense ($`n=10^4`$ $`\mathrm{cm}^3`$ ), centrally concentrated ($`0.3\mathrm{}<\theta <1.4\mathrm{}`$) outflowing gas clouds ($`Fr^2`$) with a relatively high filling factor ($`F0.01`$) and high ionization parameter $`(\mathrm{log}U=1`$). Component B of the model can be interpreted as a more extended distribution ($`1.9\mathrm{}<\theta <4.5\mathrm{}`$) of lower density gas ($`n=2\times 10^3`$ $`\mathrm{cm}^3`$) with no net outflow ($`F=\mathrm{const}.`$), with a lower filling factor ($`F=0.001)`$ and a lower ionization parameter ($`\mathrm{log}U=2`$). Component A contributes 58% of the total line flux in 22 lines listed in Table 1, with the contribution to individual lines ranging from 45% of the low excitation $`[\mathrm{S}\text{iv]}\lambda 10.5\mu \mathrm{m}`$ ($`E_{\mathrm{ion}}=34.8\mathrm{eV}`$) to more than 99.9% of the very high excitation line $`[\mathrm{Si}\text{ix]}\lambda 3.9\mu \mathrm{m}`$ ($`E_{\mathrm{ion}}=303.2\mathrm{eV}`$). Its large covering factor indicates that there is probably a significant contribution of flux from the inner $`1\mathrm{}`$ of the diffuse SW emission cone as well as from the bright NE cone. Component B contributes the remaining 42% of the total line flux, mainly in the lower excitation lines. Its covering factor is small enough for it to be concentrated mostly in the NE bright emission cone, as is observed.
The best-fit procedure indicates that models with the low oxygen abundance set fit the observed line fluxes somewhat better than models with the high nitrogen abundance set. In particular, the $`\mathrm{O}\text{iii]}\lambda 1663`$, whose unusual relative weakness was an important argument for assuming non-solar abundances (Netzer Netzer97 (1997); Kraemer et al. KRC98 (1998)), is well reproduced by the best-fit low oxygen model with a model-to-data line flux ratio of 1.4 (Fig. Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum), even though it was not included in our fit since it didn’t pass the minimal flux criterion.
The trough in the best fit SED (Fig. Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum) can be interpreted as an absorption trough due to an absorber between the continuum source and the NLR. The energy resolution of the SED template, which is limited by the computational cost of enumerating over all the SED combinations, is too low to allow detailed modeling of the absorber. Figure Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum shows an example of how a quasi-thermal big blue bump that is absorbed by neutral hydrogen would appear in our low resolution SED reconstruction. We find that the trough is consistent, for example, with an absorber that shadows the entire NLR ($`C_{\mathrm{abs}}=1)`$ and has a column density of $`N_{H^0}=6\times 10^{19}\mathrm{cm}^2`$ in neutral hydrogen, or with an absorber that allows a small leakage of unfiltered radiation ($`C_{\mathrm{abs}}=0.999`$) and a column density of $`N_{H^0}=10^{20}\mathrm{cm}^2`$. A similar trough, consistent with an absorber of $`C_{\mathrm{abs}}>0.99`$ and $`N_{H^0}=5\times 10^{19}\mathrm{cm}^2`$, was discovered in the reconstructed SED of Seyfert 1 galaxy NGC 4151 (Alexander et al. Alexander99 (1999)). That absorber was also detected in the *HUT* absorption line spectra of the UV continuum of NGC 4151 (Kriss et al. 1992a , Kriss95 (1995)) . The *HUT* spectra of $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ (Kriss et al 1992b ) do not have a high enough S/N to allow the detection of absorption lines in the scattered UV continuum of this AGN or against the stellar background (G. Kriss, private comm.). We predict that future sensitive absorption line studies should reveal the presence of such an absorber.
The bias in our results due to the fact that we neglected the line emission from the absorbing gas is likely to be small if the absorber is similar to the dense, high velocity UV absorber that was detected in NGC 4151. Such an absorber will not emit forbidden lines, and its permitted lines will be broader than typical NLR lines. Only 3 of the 22 lines we used in our fit are permitted lines, and we did not use lines that are contaminated by broad components. A highly ionized and optically thin absorber will produce strong $`\mathrm{O}\text{vi}\lambda 1035`$ line emission in excess of the typical NLR emission. It is therefore interesting that unlike the two forbidden $`[\mathrm{O}\text{iii]}\lambda 5007`$ and $`[\mathrm{O}\text{iv]}\lambda 25.9\mu \mathrm{m}`$ lines and the semi-forbidden $`\mathrm{O}\text{iii]}\lambda 1663`$ line, which are well reproduced by the best fit model, the observed $`\mathrm{O}\text{vi}\lambda 1035`$ line is 3.2 stronger than predicted (Fig. Infrared spectroscopy of NGC 1068: Probing the obscured ionizing AGN continuum).
We have, up to now, applied our SED reconstruction method to *ISO-SWS* observations of IR coronal lines of three AGN: the Seyfert 2 Circinus galaxy (Moorwood et al. Moorwood96 (1996); Alexander et al. Alexander99 (1999)), the Seyfert 1 galaxy NGC 4151 (Alexander et al. Alexander99 (1999)), and the Seyfert 2 galaxy $`\mathrm{NGC}\mathrm{\hspace{0.17em}1068}`$ (this work). In one of these (Circinus), we detect a Big Blue Bump that peaks at $`50\mathrm{eV}`$. In the other two we detect deep troughs, which are consistent (but not exclusively so) with a Big Blue Bump that is absorbed by neutral gas interior to the NLR. Our findings thus far are consistent with the picture that luminous Seyfert galaxies are powered by thin accretion disks that produce a quasi-thermal Big Blue Bump, and that in a large fraction of them the NLR sees a partially absorbed ionizing continuum, as suggested by Kraemer et al. (KRC98 (1998)).
This work was supported by DARA under grants 50-QI-8610-8 and 50-QI-9492-3, and by the German-Israeli Foundation under grant I-0551-186.07/97. |
warning/0002/astro-ph0002139.html | ar5iv | text | # A NEARBY SUPERNOVAE SEARCH: EROS2
## 1 Type Ia supernovae and cosmology
Supernovae are classified in different subtypes, according to their spectral features. Type Ia supernovae (SNIa) are believed to be explosions of carbon-oxygen white dwarfs. SNIa progenitors are likely to be binary systems, composed of a red giant and an old C/O white dwarf. The latter accretes matter from its companion until it reaches the Chandrasekar mass ($`1.4M_{}`$), and then becomes unstable. This process leads to the total thermonuclear explosive burning of the white dwarf. Thus, the total energy released should be nearly constant from one SNIa to another. These objects may therefore be used as standard candles. Indeed, photometric and spectroscopic studies have shown that SNIa compose an homogeneous sample and their peak magnitudes present a small scatter ($`20\%`$) in all colors. Furthermore, these objects are very luminous — they have been detected up to redshifts $`z1.2`$ (Aldering et al. $`^\mathrm{?}`$). Thus they constitute powerful cosmological distance indicators.
It has been shown that the absolute maximum luminosities of SNIa correlate with other observables, like the post maximum decline rate, the color at maximum, the SN spectral features, or the galaxy type. When corrections based on such correlations are applied, the relative dispersion of the peak luminosities of SNIa can be reduced to 10% (Hamuy et al. $`^\mathrm{?}`$).
Measurements of the cosmological parameters $`H_0`$, $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$ have been made by analysing the apparent peak magnitude versus redshift relation (Perlmutter et al. $`^\mathrm{?}`$, Schmidt et al. $`^\mathrm{?}`$). These analyses rely heavily on the standardization procedure outlined above. For example, the evidence for a non zero $`\mathrm{\Lambda }`$ arises from a 20% flux decrease with respect to a ($`\mathrm{\Omega }=0.2`$) universe, which is comparable to the intrinsic luminosity spread. However, our SNIa knowledge is based on few objects, namely the 17 SNIa ($`z<0.1`$) discovered before maximum during the Calan-Tololo search. This is why a number of nearby SN searches have been launched in order to increase the set of well sampled SNIa and study further the standardization corrections mentionned above.
Supernovae rates as a function of redshift are a useful tool for studying the star formation history, or constraining the galactic chemical evolution scenarios. While probing the stellar evolution, they also bring valuable information on the SNIa progenitor system, and allow us to get a better insight into the physics processes involved in these events. SNIa rates have been measured at low redshift (see e.g. Cappellaro et al. $`^\mathrm{?}`$ ) with SNe discovered using photographic plates, and at high redshift ($`z0.4`$), with automatic subtraction of CCD images by Pain et al. $`^\mathrm{?}`$. EROS2 has obtained the first determination of the SNIa rate at $`z0.15`$ (Hardin et al. $`^\mathrm{?}`$).
## 2 The EROS2 nearby supernovae search
The EROS2 experiment is mainly devoted to the search for microlensing events towards the Magellanic clouds, and towards the Galactic bulge and disk. For this purpose, the collaboration operates a 1 meter telescope, installed at the European Southern Observatory of La Silla (Chile). This instrument was specially refurbished and automated in view of a microlensing survey. It is equipped with a dichroic beam splitter and two cameras to take images simultaneously in two wide pass-bands. Each camera comprises a mosaic of 8 $`2k\times 2k`$ thick CCD’s, covering a field of view of $`0.7^o(\alpha )\times 1.4^o(\delta )`$ with a pixel size of 0.6 arcseconds.
Since this setup is particularly well suited for discovering supernovae at $`z0.050.2`$, the EROS2 collaboration launched in 1997 a systematic nearby SN search aimed at the measurement of the nearby SN explosion rates and a detailed study of the correlations between the SNIa light curve shapes and their peak absolute luminosities.
### 2.1 The search strategy
Our SN search technique consists in comparing an image of a given field with a reference image of the same field taken two or three weeks before. For this purpose, we subtract the reference frame from the search frame, after a geometric and a photometric alignment, and a matching of the seeing. We then perform an object detection on the subtracted frame. Genuine candidates are selected among these objects by applying cuts tuned with a Monte-Carlo simulation, in order to reject variable stars, asteroids and subtraction artifacts. Finally, a visual scan allows us to eliminate the last spurious candidates.
### 2.2 The first stage : 1997-1998
During the first two years, 7 search campaigns have been conducted. We monitored fields from both celestial hemispheres. In order to avoid dust absorption they were chosen far from the Galactic plane. During these first searches, 35 SNe have been discovered. Spectra could be obtained for 10 of them with the ESO 3.6m and the ARC 3.5m telescopes. 7 of the SN were of type Ia, 1 of type Ic and 2 of type II. Using this first sample, a preliminary SNIa rate at $`z0.15`$ has been obtained (see section 3).
### 2.3 A worldwide SNIa search campaign
In the spring of 1999, we participated in a worldwide search campaign <sup>b</sup><sup>b</sup>bInvolving the following 9 groups : The Nearby Galaxies SN Search Team ($``$) (Strolger, Smith et al.), EROS2 ($``$) (Spiro et al.), KAIT ($``$) (Filippenko et al.), The Mount Stromlo Abell Cluster Supernovae Search ($``$) (Schmidt, Germany & Reiss), NEAT ($``$) (Helin, Pravdo & Rabinovitz), QUEST ($``$) (Schaefer et al.), SpaceWatch ($``$) (McMillan & Larsen), The Tenagra Observatories ($``$) (Schwartz), and The Wise Observatories Supernovae Search ($``$) (Gal-Yam et al.)., led by the Supernovae Cosmology Project and coordinated by Greg Aldering (SCP). The search involved 9 groups listed below.
EROS2 discovered a subset of 16 SN among the 41 supernovae found in this campaign. Among them, 19 (7 from EROS2) turned out to be of type Ia, discovered near maximum. An overview of the follow-up data for each SN can be found in table 1. Photometric and spectroscopic data from both these SN together with discoveries announced in the same period in IAU circulars are currently being analysed.
## 3 A first measurement of the SNIa rate at $`z0.15`$
Our preliminary determination of the SNIa rate at $`z0.15`$ relies on a sample of type Ia supernovae discovered during 2 search campaigns led in October and November 1997. 120 square degrees have been covered, 8 supernovae discovered. Among them, 4 were of type Ia.
Supernovae rates $``$ in the rest frame are usually expressed in SNu, i.e. in supernovae per unit time and per unit blue luminosity ($`SNe/10^{10}L__B/100yr`$). The number of supernovae of a given type discovered during a search is related to the rate $``$ of this type of SNe through
$$𝒩\times \underset{gal}{}L_{gal}\times T_{gal}$$
(1)
where $`L_{_{gal}}`$ is the absolute luminosity of the galaxy $`gal`$, and $`T_{gal}`$ is the control time during which a SNIa could have been detected. If $`\epsilon (t,z,\mathrm{})`$ is the search efficiency, i.e. the probability to detect a SNIa with redshift $`z`$ whose maximum occured at a time $`t`$ before the observation, $`T_{gal}`$ can be written as $`T_{gal}=_{\mathrm{}}^+\mathrm{}\epsilon (t,z,\mathrm{})𝑑t`$.
The sum $`𝒮=_{gal}L_{gal}\times T_{gal}`$ is computed by Monte-Carlo integration. Firstly the galaxies in the search fields are detected using the program sextractor (Bertin et al. $`^\mathrm{?}`$). Their apparent magnitudes are derived in the $`R_c`$ band from the EROS2 magnitudes. Since the redshift of each galaxy is not known, a value of $`z`$ is generated in the Monte-Carlo procedure, using a $`p(z|R_{c_{gal}})`$ pdf derived from the Schechter law with parameter values measured by the LCRS (Lin et al. $`^\mathrm{?}`$). The absolute luminosities of each galaxy can then be calculated. The detection efficiency is fully simulated. The SN rate we thus obtain is
$$=0.44_{0.210.07}^{+0.35+0.13}h^2\mathrm{SNu}.$$
(2)
By multiplying this value by the luminous density of the universe $`\rho _L=1.4\pm \mathrm{0.1\; 10}^8hL_{}Mpc^3`$ (Lin et al. $`^\mathrm{?}`$) we obtain the rate expressed in $`h^3\mathrm{Mpc}^3\mathrm{year}^1`$
$$=0.62_{0.290.11}^{+0.49+0.19}10^4h^3\mathrm{Mpc}^3\mathrm{year}^1.$$
(3)
## Conclusion
Since 1997, the EROS2 collaboration has conducted several campaigns of supernovae searches. Our discovery rate is about 1 SN every two hours of observations, which makes us competitive with respect to other teams carrying on searches at the same $`z`$. In a first stage, 35 SNe were discovered, the light curves of 7 SNIa were studied, and a first SNIa explosion rate at $`z0.15`$ was derived. In Spring 99, EROS2 participated in a worldwide search led by the Supernovae Cosmology Project, and discovered 8 of the 19 SNIa near maximum found by the consortium. Photometric and spectroscopic follow-up data are currently been analyzed, and results are expected to come out soon.
## References |
warning/0002/hep-th0002037.html | ar5iv | text | # Untitled Document
hep-th/0002037 RUNHETC-2000-05
Stability and BPS branes
Michael R. Douglas<sup>&</sup> , Bartomeu Fiol and Christian Römelsberger
Department of Physics and Astronomy
Rutgers University
Piscataway, NJ 08855–0849
<sup>&</sup>I.H.E.S., Le Bois-Marie, Bures-sur-Yvette, 91440 France
mrd, fiol, roemel@physics.rutgers.edu
We define the concept of $`\mathrm{\Pi }`$-stability, a generalization of $`\mu `$-stability of vector bundles, and argue that it characterizes $`𝒩=1`$ supersymmetric brane configurations and BPS states in very general string theory compactifications with $`𝒩=2`$ supersymmetry in four dimensions.
February 2000
1. Introduction
Compactifications of the heterotic string or type II superstring theory with D-branes include a choice of gauge field configuration. Characterizing the possibilities and analyzing their physics is a central problem in this subject.
The most intensively studied case is a gauge field on a Calabi-Yau manifold preserving $`𝒩=1`$ supersymmetry in four dimensions. As is well-known, solutions of the Yang-Mills equations preserving this supersymmetry correspond by the work of Donaldson, Uhlenbeck and Yau to holomorphic vector bundles satisfying the condition of $`\mu `$-stability, a quasi-topological condition depending on the Kähler class of the Calabi-Yau.
In general, quantities which depend on the Kähler class are modified in string theory by world-sheet instanton corrections. In the case of heterotic strings, space-time instanton corrections can also enter. The true picture of this moduli space can in some cases be obtained by duality; mirror symmetry in type II compactification and type II-heterotic duality in heterotic compactification. These corrections can drastically alter the large volume and classical picture, and this strongly suggests that the $`\mu `$-stability condition must be significantly modified as well.
Based on recent work on BPS D-brane configurations in Calabi-Yau compactification, we propose a generalization of the $`\mu `$-stability condition which takes these corrections into account. The dependence on the Kähler class is replaced by dependence on the periods $`\mathrm{\Pi }`$ of the Calabi-Yau and thus we call it “$`\mathrm{\Pi }`$-stability.” The usual definitions of Calabi-Yau periods and $`𝒩=2`$ central charges must be generalized slightly (in a way already suggested by mathematicians) to make the definition, as we will also explain.
$`\mathrm{\Pi }`$-stability is a precisely defined condition and can be studied using the same mathematical techniques as $`\mu `$-stability; we will consider the noncompact Calabi-Yau $`𝒪_{P^2}(3)`$ in detail in an upcoming work . In the present work we will state the ideas and assumptions which lead us to this proposal and check that it is compatible with the known physics of BPS branes and marginal stability in solvable examples; the large volume limit, the orbifold limit and the large complex structure limit, in which it is related to a condition governing stability of special Lagrangian manifolds formulated by Joyce .
2. The proposal
Let $`_c`$ and $`_k`$ be the complex structure and complexified Kähler moduli spaces of a Calabi-Yau $`M`$ in string theory, with complex dimension $`b_{2,1}`$ and $`b_{1,1}`$. $`_k`$ is best defined (when possible) as the complex structure moduli space of a mirror manifold $`W`$. Let $`z^i`$ be local coordinates on $`_k`$ and $`\mathrm{\Pi }_aC^{2b_{1,1}+2}`$ be a vector of periods as defined in special geometry (e.g. when a mirror exists, the periods $`\mathrm{\Omega }`$ of the holomorphic three-form of the mirror). These are defined up to overall normalization; we choose a particular (but arbitrary) normalization at each point in $`_k`$.
Let $`E`$ be a holomorphic cycle carrying a vector bundle, or some generalization of this idea appropriate to string theory. The proposal rests on the idea that these can be defined knowing only the complex structure of $`M`$. We require a category of these and an idea of homomorphism. $`E^{}`$ is a subbundle of $`E`$ if $`E^{}E`$ and there exists a holomorphic embedding of $`E^{}`$ in $`E`$ or in other words an injective holomorphic map from sections of $`E^{}`$ to $`E`$. More generally, $`E^{}`$ is a subobject of $`E`$ if there is a monomorphism (an injective homomorphism) in $`\mathrm{Hom}(E^{},E)`$.
Let $`Q(E)`$ be the charge of $`E`$ or appropriate generalization, defined so that the central charge of a brane associated with $`E`$ is $`Z=Q(E)\mathrm{\Pi }`$. Clearly the precise definition of $`Q(E)`$ depends on our basis for $`\mathrm{\Pi }`$. In the A picture (special Lagrangians on $`W`$), $`Z=_\mathrm{\Sigma }\mathrm{\Omega }`$ for a brane wrapped on the cycle $`\mathrm{\Sigma }`$ and we are just choosing a basis for $`H^3(W)`$. In the B picture, a definition with strong motivations from D-brane physics and the mathematics of K-theory takes $`Q(E)=\mathrm{ch}(E)\sqrt{\widehat{A}}`$ where $`\mathrm{ch}(E)`$ is the Chern character and $`\sqrt{\widehat{A}}`$ a topological invariant of the CY. By taking a different basis for the periods, one could also work in conventions where $`Q(E)=\mathrm{ch}(E)`$, which tend to be more convenient for comparisons with the mathematical literature on vector bundles. In any case, we want to emphasize that the concepts entering our definitions are independent of convention.
We next define the “grade” $`\phi (E)`$ of a BPS brane at a point in moduli space with periods $`\mathrm{\Pi }`$ to be
$$\begin{array}{cc}\hfill \phi (E)& =\frac{1}{\pi }\mathrm{arg}Z(E)\hfill \\ & =\frac{1}{\pi }\mathrm{Im}\mathrm{log}Z(E).\hfill \end{array}$$
One might also write $`\phi (E;\mathrm{\Pi })`$ to make the dependence on the periods explicit. The terminology is intended to be an analog both of the slope of a vector bundle and of the grading in a derived category. Note that the grade will ultimately be defined to take values in $`R`$, not $`[0,2)`$. Given the grading at some point in moduli space, we will define them elsewhere by analytic continuation of $`Z(E)`$. Clearly the well-definedness of this will require some discussion, which we will make below.
We now define $`E`$ to be $`\mathrm{\Pi }`$-stable at a point in moduli space with periods $`\mathrm{\Pi }`$ if, for every subobject $`E^{}E`$, we have
$$\phi (E^{})<\phi (E).$$
We then conjecture that the BPS branes in the theory (for bulk moduli $`\mathrm{\Pi }`$) are the $`\mathrm{\Pi }`$-stable objects with unbroken gauge symmetry $`U(1)`$.
When (2.1) degenerates to equality, it is clear that a decay of $`E`$ to products including $`E^{}`$ would be physically allowed. The additional physical information in the statement of $`\mathrm{\Pi }`$-stability is the statements that certain decays $`EE^{}+\mathrm{}`$ are not possible despite the degeneration of (2.1) (namely, those for which $`E^{}`$ is not a subobject), that the bound state $`E`$ will exist on a specific side of this line and not on the other, that the other cases $`\phi (E)\phi (E^{})Z`$ do not lead to decays, and that all objects not destabilized by particular subobjects actually exist. We proceed to check these points in various limits of the theory.
3. The large volume limit
In the large volume limit of $`_k`$ $`\mathrm{\Pi }`$-stability (2.1) reduces to $`\mu `$-stability. In this limit, the periods can be associated with the $`2k`$-cycles of $`M`$, and in terms of the triple intersection form
$$c_{ijk}=_M\omega _i\omega _j\omega _k$$
we have (up to lower order corrections)<sup>*</sup><sup>*</sup> The conventions are the ones in which $`Q_{2p}=\mathrm{ch}_{3p}(E)`$.
$$\begin{array}{cc}& \mathrm{\Pi }_6=\frac{1}{6}c_{ijk}t^it^jt^k\hfill \\ & \mathrm{\Pi }_4^i=\frac{1}{2}c_{ijk}t^jt^k\hfill \\ & \mathrm{\Pi }_2^i=t^iB^i+iV^i\hfill \\ & \mathrm{\Pi }_0=1\hfill \end{array}$$
with $`B=\omega _i\mathrm{Re}t^i`$ and $`J=\omega _i\mathrm{Im}t^i`$. The leading terms ($`|B|V`$) in (2.1) then take the form
$$\begin{array}{cc}\hfill \phi (E)& =\frac{1}{\pi }\mathrm{Im}\mathrm{log}\mathrm{\Pi }_6+\mathrm{Im}\frac{1}{2\pi \mathrm{\Pi }_6\mathrm{rank}E}JJ\mathrm{c}_1(E)\hfill \\ & \frac{3}{2}+\frac{3}{\pi V}\frac{1}{\mathrm{rank}E}JJ\mathrm{c}_1(E)+\mathrm{}\hfill \end{array}.$$
The leading nonconstant term for $`V1`$ and $`V|B|`$ is proportional to the slope, so in this limit $`\mathrm{\Pi }`$-stability reduces to $`\mu `$-stability .
Some features of the world-volume physics of $`\mu `$-stability have been discussed by Sharpe . For $`b^{1,1}>1`$, one can have “walls” in Kähler moduli space on which the $`\mu `$-stability of bundles changes. These have been much studied for complex surfaces for application to Donaldson theory. Physically, a wall-crossing leading to decay of a bound state corresponds to an enhancement of gauge symmetry on the wall (where the bundle is semistable) followed by D-term supersymmetry breaking, in the simplest case governed by the potential
$$V=(|\varphi |^2\zeta )^2.$$
The same relation to the D-terms will hold for $`\mathrm{\Pi }`$-stability.
One can get a first indication of why we will need to extend the grading beyond the region $`[0,2)`$ by dropping the condition $`|B|V`$. Since string theory only depends on the gauge-invariant combination $`BF`$, the situation with arbitrary values of $`B`$ can be related to that for $`|B|V`$ by considering branes with non-zero $`c_1=F`$, i.e. tensoring our bundles with appropriate line bundles. Once we start to consider $`c_1`$ of order $`V`$ (in string units), the grading will leave the region $`[0,2)`$: for example, using line bundles we can get $`\frac{1}{\pi }\mathrm{Im}\mathrm{log}(F+iV)^3`$ which takes values in $`[0,3)`$.
Although one might wonder if the stability of a brane with such large values of the flux changes from the usual large volume idea, we have no evidence for this. Actually, if we do not extend the grading to a nonperiodic variable, the definition (2.1) of stability is not sensible.
According to our definition of subobject, a single six-brane carrying a line bundle with $`JJc_1=n`$ (let us call it $`O(n)`$) will have infinitely many subobjects, namely the branes $`O(m)`$ with $`m<n`$. Although these are not considered subobjects for the usual definition of $`\mu `$-stability (they are never relevant for this definition anyways), they can become relevant away from the large volume limit, and it is not natural to drop them. If we do not extend the grading, these lead to many nonsensical predictions; in particular decays of branes into heavier constituents.
Requiring that in the $`c_1\mathrm{}`$ limit, the grading goes over to that for the D$`0`$-brane suggests setting the gradings for pure (trivial bundle) $`2p`$-branes to be $`\phi (\mathrm{\Sigma }_{2p})=3p/2`$, in the conventions of this section. (Any choice for the origin $`\phi =0`$ is of course a convention.)
We will give further arguments below for why extending the grading is sensible and correct in string theory.
4. Considerations from world-volume effective theory
As we go away from the large volume limit, our primary tools will be the constraints of $`𝒩=1`$ supersymmetry on the world-volume effective theory, and the “decoupling” statement of , that superpotential and D terms only depend respectively on the complex and Kähler moduli for B branes and the mirror for A branes. This statement will be further justified elsewhere.
Clearly the configurations of a single BPS brane on a CY can be described by a $`d=4`$, $`𝒩=1`$ effective world-volume theory. We will also consider non-BPS combinations of branes and non-BPS branes which can be described by non-supersymmetric vacua of such a theory, obtained by combining the theories of the various BPS constituents and adding degrees of freedom corresponding to open strings stretched between these branes.
It is not strictly true that all BPS bound states are described by supersymmetric vacua of the resulting world-volume theory. This condition is too restrictive as it ignores the possibility that a different $`𝒩=1`$ supersymmetry is unbroken, given by a combination of the linearly realized $`𝒩=1`$ supersymmetry and an inhomogeneous $`𝒩=1`$ symmetry present in all theories containing a decoupled $`U(1)`$ (such as D-brane bound states). This latter is simply the shift of the decoupled gaugino $`\delta \chi =ϵ^{}`$.
Such supersymmetry preserving vacua are characterized by having a non-zero potential which comes entirely from an overall constant shift of the $`U(1)`$ D terms (i.e. $`D_i=D_j`$ for every $`U(1)`$ factor). Other supersymmetry breaking vacua, in particular those in which the F terms are non-zero, correspond to non-BPS bound states. The states found by Sen in K3 compactification are an example which can be understood this way.<sup>*</sup><sup>*</sup> This point was developed in a discussion with A. Sen.
A way to distinguish the non-BPS situation from the BPS vacuum preserving a different $`𝒩=1`$ supersymmetry is to note that D terms do not give masses to fermions, while F terms do. Thus the possibility of two BPS branes combining to a BPS bound state is signalled by a massless Ramond open string fermion living in a chiral multiplet (it is massless when gauge symmetry is unbroken; i.e. at the unstable maximum of the D term potential). We will use these characterizations of BPS ground states below.
We finally remark on the relation between “non-BPS branes” and supersymmetry breaking in physical string compactifications. First, to get supersymmetry breaking one should break all of the $`𝒩=1`$ supersymmetries, so indeed one is looking for “non-BPS” branes with breaking by F and D terms together. One must furthermore ensure breaking for all values of the bulk moduli, which has not been accomplished in examples considered so far.
5. D-branes and Homological Algebra
To study D-branes and their stability at arbitrary points in moduli space, we need to describe them as objects which allow a study of their moduli space and their subobjects (potential decay products). A natural framework for this seems to be homological algebra \[15,,14\], and in particular the language of extension groups, $`\mathrm{Ext}^p(V,W)`$. We will try now to give some physical intuition for the role of $`\mathrm{Ext}^p(V,W)`$.
For bundles, $`\mathrm{Ext}^p(V,W)`$ is the same as $`H^p(M,V^{}W)`$. Recall that one physical appearance of the complex cohomology groups is that they count fermion zero modes. As is well-known, holomorphic $`p`$-forms on a Calabi-Yau manifold are directly related to spinors with chirality determined by the parity of $`p`$. An element of $`H^p(M,V^{}W)`$ will thus correspond to a zero mode of the Dirac operator coupled to the bundle $`V^{}W`$.
In the general discussion of BPS branes as quantum objects and in $`(0,2)`$ sigma models of heterotic string compactification , sheaves can sometimes be used instead of bundles, and the generalization to $`\mathrm{Ext}`$ has been used in this context.
Another context in which the $`\mathrm{Ext}`$ groups are useful turns out to be quiver gauge theories. Indeed these provide very elementary examples, which will be discussed at length in . Let us give the briefest introduction here, to make some necessary points.
We recall that a quiver is a directed graph with vertices $`vV`$ and arrows $`aA`$ from vertices $`ta`$ to $`ha`$; an associated gauge theory is labelled by a dimension (we also use the terms weight or charge) vector $`n`$. It has gauge group $`_vU(n_v)`$, matter content $`R_{i,\overline{j}}^a`$ transforming in $`(n_{ta},\overline{n}_{ha})`$ and a superpotential.
A moduli space of solutions to the superpotential constraints $`W^{}=0`$ is a complex variety, and the set of these for various $`n_i`$ provides another example of the type of category of holomorphic objects we have in mind. Following the quiver literature, we will refer to points in these moduli spaces as “representations” of the quiver.
A homomorphism between two representations $`R`$ and $`S`$ is defined as a set of linear maps $`\varphi _v:R_vS_v`$ for each vertex $`v`$ satisfying $`S^a\varphi _{ta}=\varphi _{ha}R^a`$ for each arrow $`a`$. We say that $`R`$ is a subrepresentation of $`S`$ if there is an injective homomorphism from $`R`$ to $`S`$.
For $`p=0`$, we have $`\mathrm{Ext}^0(V,X)=\mathrm{Hom}(V,X)`$. The definition we gave of subobject in section 2 makes sense in any abelian category: we say that $`V`$ is a subobject of $`X`$ is there is an injective homomorphism from $`V`$ to $`X`$. We will see in section 7 that subobjects play the same role in quiver gauge theory that they did in the discussion of bundles.
We note that the relation of subobject is not necessarily determined by the charges of the two objects; two objects of the same charge might differ in the charges of subobjects they admit. A common situation is that all “generic” objects of a given charge admit the same subobjects, while certain degenerate objects (for example semisimple ones) admit more. If so, we can (a bit loosely) talk about one weight vector $`n(E^{})`$ being a subobject of another weight vector $`n(E)`$. This means that a generic representation $`E`$ of weight $`n(E)`$ will have a subrepresentation $`E^{}`$ of weight $`n(E^{})`$.
Moving to $`p=1`$ and $`\mathrm{Ext}`$ (one sometimes leaves off the superscript $`1`$), it is well known that elements of $`H^1(M,\mathrm{End}V)`$ correspond to deformations of the complex structure of $`V`$. A related statement which is relevant for bound state problems is that $`H^1(M,V^{}W)`$ corresponds to possible deformations of the direct sum bundle $`VW`$. Generalizing these observations to arbitrary points in moduli space, one criterion we might apply to find out if $`V`$ and $`W`$ can form a bound state is to ask if $`\mathrm{Ext}(V,W)\mathrm{Ext}(W,V)0`$.
In general, an extension $`X`$ is associated with an exact sequence
$$0W\stackrel{\varphi }{}XV0$$
which does not split: $`XWV`$. This is also the definition of extension for quiver representations, and can be used to define $`\mathrm{Ext}^1(V,W)`$ in this context.
The second arrow in this sequence, representing $`\varphi \mathrm{Hom}(W,X)`$, gives us a way to think of a homomorphism as associated with a particular way to form a bound state. Another picture is that these are “potential gauge symmetries,” which can become unbroken when the bound state becomes marginally stable. This will be signalled by the appearance of a $`\mathrm{Hom}(X,W)`$ which (if $`W`$ and $`X`$ are not isomorphic) implies that $`X`$ is no longer simple.
A primary tool for determining the dimensions of these groups is the Grothendieck-Riemann-Roch theorem, which can be thought of as a special case of the index theorem for the Dirac operator applicable for holomorphic bundles. This generalizes to sheaves and quivers, and allows evaluating
$$\chi (V,W)=\underset{i}{}(1)^idim\mathrm{Ext}^i(V,W)$$
in terms of the Chern classes of $`V`$, $`W`$ and $`M`$. Although this does not determine any of the individual dimensions directly, given appropriate further assumptions one might use it to make statements such as $`dim\mathrm{Hom}=0`$ implies $`\chi 0`$.
For bundles on a Calabi-Yau, $`\chi (E^{},E)`$ is mirror to the intersection number $`I(E^{},E)`$ between three-branes, and is antisymmetric. For bundles on a divisor of a CY, the two are related as
$$I(E^{},E)=\chi (E^{},E)\chi (E,E^{}).$$
Whereas the intersection number counts all fermionic massless strings between $`E^{}`$ and $`E`$, in this case $`\chi `$ distinguishes some of these and gives more information.
As an elementary example, let us consider the quiver with two nodes and three arrows between them (and no superpotential). Its representations correspond to arbitrary configurations of three chiral multiplets which transform as $`(\overline{n}_1,n_2)`$ under the gauge group $`U(n_1)\times U(n_2)`$. The analog of (5.1) in this case is \[13,,10\]:
$$\chi (E^{},E)=dim\mathrm{Hom}(E^{},E)dim\mathrm{Ext}(E^{},E)=n_1^{}n_1+n_2^{}n_23n_1^{}n_2.$$
This theory also turns out to describe the moduli spaces of certain bundles on $`P^2`$ and (5.1) is the corresponding intersection form on the local mirror to $`C^3/Z^3`$.
The simplest example of a subobject is $`E^{}=(\mathrm{0\; 1})`$ which is a subobject of any quiver representation $`(n_1n_2)`$ with $`n_2>0`$ (the relation defining the homomorphism degenerates). As a more subtle example, we give the pair $`E^{}=(\mathrm{1\; 0})`$ and $`E=(\mathrm{3\; 1})`$ for which $`dim\mathrm{Hom}(E^{},E)=0`$, $`dim\mathrm{Hom}(E,E^{})=3`$, and the Ext groups are zero. This allows us to illustrate several points: $`i)`$ since $`\mathrm{Hom}(E^{},E)=0`$, $`E^{}`$ can not be a subobject of $`E`$; $`ii)`$ although $`\mathrm{Hom}(E,E^{})0`$, $`E`$ is clearly not a subobject of $`E^{}`$; $`iii)`$ the vanishing of the Ext groups corresponds to the fact that $`(\mathrm{4\; 1})`$ is not a bound state (all such configurations have unbroken gauge symmetry); $`iv)`$ the intersection number $`I(E^{},E)=3`$ by itself is not enough information to see this.
We gave two examples of categories which appear in string theory, but an important question is whether some category of holomorphic objects can describe all B branes (in some background CY) over all of Kähler moduli space. As we discuss in the conclusions, the derived category of coherent sheaves is a natural candidate.
6. Marginal stability and special Lagrangian geometry
For a given brane (bundle) $`E`$ and subbrane $`E^{}`$, the inequality in condition (2.1) will degenerate to equality on walls of real codimension one in $`_k`$. These are the familiar “lines of marginal stability” in $`𝒩=2`$, $`d=4`$ supersymmetric theories which physically are the only lines on which the condition for $`E`$ to exist as a BPS brane could change. We now argue that the bound state $`E`$ will exist on the side of the line predicted by (2.1) and not on the other.
The mirror interpretation of these processes involves joining and splitting of special Lagrangian submanifolds on the mirror manifold $`W`$, as was studied by Joyce . As one varies the complex structure of $`W`$, it is possible for a pair of special Lagrangian manifolds of homology class $`[\mathrm{\Sigma }_1]`$ and $`[\mathrm{\Sigma }_2]`$ to intercommute producing a single manifold of class $`[\mathrm{\Sigma }_1]+[\mathrm{\Sigma }_2]`$. Conversely, a single brane can become unstable to split into a pair.
Joyce found a condition (, section 7) which applies to the local neighborhood of an intersection and to branes with small differences in the grade (in our terminology), and predicts on which side of the marginal stability line the bound state will be stable. This condition agrees precisely with (2.1) if we require that the intersection number $`EE^{}>0`$. If it is negative, the role of the two branes is exchanged.
In string theory, this process can also be understood as a stretched open string between the two three-branes becoming tachyonic. Approximating the neighborhood of the intersection as flat space, there are six complex fields describing light strings stretched between D$`3`$-branes; their squared masses are linear in the angles of rotation. These masses can be described as a combination of superpotential and D-term masses in an effective $`𝒩=1`$ field theory, and crossing the line of marginal stability changes the sign of the D-term mass, as was pointed out in .
As we argued in section 4, only the solutions with vacuum energy coming entirely from D terms can correspond to BPS bound states, and these can be identified by the presence of massless fermions. This means that the geometric condition $`EE^{}>0`$ for A branes to intersect is not actually the condition which governs BPS decay. The correct condition instead, as suggested in section 5, is that $`\mathrm{Hom}(E^{},E)0`$ and $`\mathrm{Hom}(E,E^{})=0`$.
This is a stronger condition than non-zero intersection number and the idea that a non-zero intersection number implies the existence of such a decay is contradicted in numerous examples on the B side. Indeed in orbifold theories one can realize the counterexample given in section 5. A further complication with predicting decays on the A side is that the superpotential can also have world-sheet instanton corrections , which could lift the degrees of freedom responsible for the decay.
Deriving holomorphic properties from the special Lagrangian picture looks difficult at present. In practical applications of mirror symmetry, one is usually better off doing computations on the side which does not receive stringy corrections. For the holomorphic structure and specifically the computation of $`dim\mathrm{Hom}(E^{},E)`$ this means the B side.
Given the existence of the homomorphism (so the decay can happen), we are basically asserting that the special Lagrangian picture correctly predicts the direction of the decay. This seems almost beyond doubt in the large volume limit on $`W`$ or equivalently the large complex structure limit in $`_c`$. Given our claim that $`_c`$ and $`_k`$ are effectively decoupled for the question of stability, this result is very strong.
7. The orbifold limit
A very different region of moduli space which can be treated exactly is the orbifold $`C^3/\mathrm{\Gamma }`$ with $`\mathrm{\Gamma }SU(3)`$ and the non-compact Calabi-Yaus obtained by substringy resolution of the singularity. As explained in \[8,,3,,4\], very general D-branes on these spaces are described by quiver gauge theories, with RR charges mapped into the ranks of the gauge groups. We now ask whether a BPS state exists with a particular charge vector. In this formalism it will be a bound state of fractional branes, and we need to know whether the associated gauge theory admits an $`𝒩=1`$ supersymmetric vacuum (in the sense of section 4) which breaks the gauge symmetry to $`U(1)`$.
The question of marginal stability in this context becomes the following. Let us imagine we have some solutions to the superpotential constraints (quiver representations): for what values of the Kähler moduli do they correspond to BPS states?
The dependence of the gauge theory potential on the Kähler moduli of the background CY is through Fayet-Iliopoulos terms; in other words the moment maps for the $`U(1)`$’s. Denote these as $`\theta _i`$; they will satisfy $`_i\theta _i=0`$. For a cyclic orbifold $`C^3/Z_n`$ they satisfy no other relations; the real dimension of $`_k`$ equals the number of remaining FI terms.
The question of whether such a moduli space actually contains a supersymmetric vacuum can be answered using the work of King on stability of quiver representations . This will be true if and only if the representation $`R`$ is a direct sum of $`\theta `$-stable representations. A $`\theta `$-stable representation $`R`$ is one for which
$$\underset{v}{}\theta _vn_v(R)=0$$
and for every subrepresentation $`R^{}`$ we have
$$\underset{v}{}\theta _vn_v(R^{})>0.$$
The condition (7.1) simply follows by taking traces of the D-flatness conditions. The condition (7.1) is proven by the techniques of geometric invariant theory. In general one finds a solution of the D-flatness conditions by minimizing the potential in a complexified gauge orbit. A subrepresentation violating (7.1) exists if and only if a one-parameter subgroup of a certain central extension (depending on $`\theta `$) of the complexified gauge group with a limit point exists; in this case the minimum of the potential is at the limit point and off of the original complexified gauge orbit.
A very simple example illustrating (7.1) is the quiver of section 5. The fact that $`(\mathrm{0\; 1})`$ is a subobject of any $`(n_1n_2)`$ with $`n_2>0`$ implies the (obvious) condition that non-trivial supersymmetric vacua exist only if $`\theta _2>0`$. More complicated quivers are required to illustrate the possibility of more complicated decays .
As we discussed earlier, the most general vacuum corresponding to a BPS state can have a constant non-zero potential arising from D terms. Given the physical FI parameters $`\zeta _i`$, such a vacuum can be obtained by using a solution to the D-flatness conditions for a different set of ’FI terms’ $`\theta _i`$. Explicitly, we have
$$\begin{array}{cc}\hfill V& =\underset{i}{}\mathrm{tr}(D_i\theta _i+\theta _i\zeta _i)^2\hfill \\ & =\underset{i}{}n_i(\theta _i\zeta _i)^2.\hfill \end{array}$$
The supersymmetric minimum will use the choice of $`\theta `$ satisfying (7.1) which minimizes the total energy.
Very near the orbifold point, we can use a quadratic approximation to the potential and kinetic term. The $`\theta `$ minimizing the potential is then
$$\theta =\zeta \frac{\zeta n}{en}e$$
where $`e`$ is the vector with all components $`1`$. The condition for $`\theta `$-stability is then
$$\frac{\zeta n^{}}{en^{}}>\frac{\zeta n}{en}.$$
To compare this result with the prediction of $`\mathrm{\Pi }`$-stability, we need an expression for the FI terms in terms of the periods $`\mathrm{\Pi }`$. The orbifold points have an enhanced discrete symmetry and this allows us to write $`\mathrm{\Pi }_k\frac{1}{n}z_me^{2\pi imk/n}`$ at linear order. On the other hand, the FI terms are cyclically permuted under the symmetry; there exists a basis (the twist fields at the orbifold point) in which<sup>*</sup><sup>*</sup> We have only checked these signs carefully for $`C^3/Z_3`$, where they follow if we define the line from the orbifold point to the conifold point to be real $`z>0`$. $`\zeta _m=\mathrm{Im}\mathrm{\Pi }_m`$. Substituting this into (2.1) reproduces (7.1).
Thus $`\mathrm{\Pi }`$-stability appears to be correct in this limit as well.
8. Why the grading?
One might consider a simpler definition of stability which only depends on $`Z/Z^{}`$ as a conventional complex variable, but this turns out not to be possible. First of all, there cannot be a decay $`EE^{}+\mathrm{}`$ when $`Z/Z^{}`$ is on the negative real axis, by conservation of energy. (The decay $`E\overline{E}^{}+\mathrm{}`$ might be possible but is covered independently by checking whether $`\overline{E}^{}`$ is a subobject of $`E`$.) Thus we have no decay when $`\phi (E^{})\phi (E)=1`$.
This argument does not rule out a decay when $`\phi (E^{})\phi (E)=2`$. If this were possible, we would need to keep track of which sheet of the complex plane $`Z/Z^{}`$ sits on anyways, to get the direction of decay correct. So there cannot be a condition for stability which depends only on central charges; it must see the grading.
One should next ask whether the definition of grading in terms of analytically continued central charges is sensible. Indeed, there is no obvious prescription for adding graded central charges, so one cannot analytically continue the periods $`\mathrm{\Pi }`$ and then define $`Z=Q\mathrm{\Pi }`$. One can certainly analytically continue all the central charges separately, but one might then expect compatibility conditions between them.
A better conceptual basis for this definition uses the idea of “graded Lagrangian submanifold.” This idea originated in the work of Fukaya and Kontsevich \[12,,23\] and has appeared in other discussions of mirror symmetry , but it has not played a direct physical role until now.
The basic point is that the symplectic group, $`Sp(2n)`$, has a non-zero first homotopy group $`\pi _1Z`$. A general variation of a Lagrangian submanifold can be described by an element of $`Sp(2n)`$ at each point. Given a metric, we can consider orthonormal frames, and reduce the action to $`U(n)`$, its maximal compact subgroup.
For variations of special Lagrangian manifolds, one must have the same $`U(1)`$ element at each point on the surface, since this acts on the holomorphic three-form. Thus any closed loop, such as one induced from a closed loop in complex moduli space (or Teichmüller space), can be associated with a winding number in $`\pi _1(U(n))`$. This allows us to extend the gradings of the sL-manifolds throughout the Teichmüller space.
We still have the problem that there is no clear prescription for adding or subtracting graded central charges. The only obvious way to avoid this problem is to only allow decays to constituents which have the same grade. We have also found in concrete examples that postulating that decays happen for more than one value of $`\phi (E^{})\phi (E)`$ generally leads to inconsistencies such as decay to constituents which are more massive than the original object. These two points lead us to rule out such decays in our proposal (2.1).
We should say that, in our opinion, this is the weakest point in the arguments for our proposal. Although our proposal is clearly the simplest of this type, at present it also seems conceivable that the correct condition in string theory is more complicated, with decays at more than one even integral value of $`\phi (E^{})\phi (E)`$. Although somewhat problematic, a reason not to rule out this possibility is the idea that $`\phi (E)`$ in string theory could be a periodic variable. Although this is not immediately incompatible with (2.1), many further consistency conditions would need to be satisfied. In any case, we expect that further study of the explicit $`C^3/Z_3`$ example will pick out (at most) one viable proposal.
We conclude with the remark that the explicit analyses in the large volume and orbifold limits determine the gradings for the objects under consideration, providing a starting point for the definition by analytic continuation of the central charges. Nontrivial gradings will then arise when one moves large distances in moduli space.
9. Conclusions
D-branes in type II string compactifications with $`𝒩=2`$ supersymmetry in the bulk are specified by a choice of embedding and a choice of gauge bundle. We proposed a general criterion, $`\mathrm{\Pi }`$-stability, for determining the configurations which preserve $`𝒩=1`$ supersymmetry. The condition depends on the moduli of the Calabi-Yau and combines elements of the A and B brane mirror pictures.
In a given example the criterion has two elements, which correspond respectively to the problems of solving “F-flatness” (superpotential) and “D-flatness” conditions in the $`𝒩=1`$ effective theory.
The problem of finding the set of F-flat configurations is independent of the D terms and given our decoupling assumption is independent of the Kähler moduli (although it could happen that in different regimes of Kähler moduli space, very different configurations survive the D-flatness conditions). In the large volume limit these are holomorphic sheaves; more generally one expects a similar category of holomorphic objects which admit a definition of homomorphism. As another example of the category of holomorphic objects, we discussed the orbifold limit of non-compact CY’s, in which the category is that of quiver representations. We will discuss the example of $`C^3/Z_3`$ in detail in .
The correct category must admit an action by the monodromy transformations of the CY, as discussed in \[17,,27\]. In principle it can contain constituents which are not simultaneously BPS. The question of what states are simultaneously BPS depends on Kähler moduli, and thus (given the decoupling) it cannot enter in defining the space of holomorphic objects. It is not yet clear whether such a general description can be made with a finite number of degrees of freedom.
Following the seminal proposal of , the “derived category” based on the category of sheaves (which in the orbifold example is the same as the derived category based on the quivers) is a natural candidate to explore.<sup>*</sup><sup>*</sup> (Note added in v3): As pointed out to us by E. Sharpe and by R. Thomas, the derived category does not have a clear notion of subobject, so it is not at all obvious that it can be used as the category of our proposal. In any case we regard it as a useful clue to the correct category, as several of its features do have analogs in string theory \[30,,31,,11\].
Given these objects, $`\mathrm{\Pi }`$-stability is a precise definition of stability of holomorphic objects which can be analyzed just knowing the periods at the point in moduli space of interest, and the inclusion relations between objects. Computing periods is a well-studied problem and obtaining the gradings appears to be easy as well. The inclusion relations are not necessarily easy to get, but they are clearly necessary for analyzing any definition of stability motivated by geometric invariant theory, and seem to carry direct physical information about the products of a decay process. These considerations lead us to believe that this is the simplest general form of the stability condition one could propose. We checked its validity in several limits; conversely, a failure of the condition would appear to contradict one or more elements of the currently accepted picture of branes on CY. (As we noted in section 8, there is a similar but more complicated variant proposal which we have not ruled out at this point.) It is also worth mentioning that there are many other consistency checks one could make with string theory; for example that objects are only destabilized by lighter constituents.
Our proposal was originally motivated by the physics of D-branes in weakly coupled type II theory, but the concepts required for its statement as well as the principles justifying it are quite general, and one might expect it to apply to BPS states in quite general $`𝒩=2`$ theories or at least those which can be embedded in or are dual to D-brane theories in type II strings. For example, it will be interesting to see if marginal stability in $`𝒩=2`$ supersymmetric gauge theory (as studied for example in \[1,,22\]) can be described by using a suitable category of objects.
In and subsequent work we will discuss the concrete BPS branes which arise in particular Calabi-Yaus, and analyze their $`\mathrm{\Pi }`$-stability. The direct physical applications of such work might include a better understanding of dualities of $`𝒩=2`$ and $`𝒩=1`$ theories, computations of black hole entropy, and quantitative approaches to the study of supersymmetry breaking.
We would like to thank O. Aharony, N. Berkovits, D.-E. Diaconescu, S. Kachru, A. Klemm, C. Lazaroiu, J. Maldacena, M. Mariño, G. Moore, H. Ooguri, A. Sen, E. Sharpe, R. Thomas and E. Witten for helpful discussions and comments.
This research was supported in part by DOE grant DE-FG02-96ER40959.
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warning/0002/hep-ph0002127.html | ar5iv | text | # 1 Introduction
## 1 Introduction
For better understanding and precise measuring the weak-action properties of heavy quarks, governed by the QCD forces, we need as wide as possible collection of snapshots with hadrons, containing the heavy quarks. Then we can provide the study of heavy quarks dynamics by testing the various conditions, determining the forming of bound states as well as the entering of strong interactions into the weak processes. So, a new lab for such investigations is a doubly heavy long-lived quarkonium $`B_c`$ recently observed by the CDF Collaboration for the first time.
This meson is similar to the charmonium and bottomonium in the spectroscopy, since it is composed by two nonrelativistic heavy quarks, so that the NRQCD approach is well justified to the system. The modern predictions for the mass spectra of $`\overline{b}c`$ levels were obtained in refs. in the framework of potential models and lattice simulations. The arrangement of excitations is close to what was observed in the charmonium and bottomonium. However, the feature of $`B_c`$-mesons is an absence of annihilation into light quarks, gluons and leptons due to QCD and QED, that implies the higher excitations decay into the low lying levels and ground state due to the emission of photons and pion pairs. The measured value of $`B_c`$ mass yet has a large uncertainty
$$M_{B_c}=6.40\pm 0.39\pm 0.13\mathrm{GeV},$$
in agreement with the theoretical expectations.
The production mechanism for the $`B_c`$-meson was studied in refs. . The most simple picture takes place for the production in the $`e^+e^{}`$-annihilation, where the universal perturbative fragmentation functions can be analytically calculated for the S-, P- and D-wave levels in the framework of factorization for the hard production of quarks and their soft binding into the hadron, which can be reliably described in the potential models. In hadron collisions, the fragmentation regime takes place at the transverse momenta $`p_Tm_{B_c}`$, and at $`p_Tm_{B_c}`$ the subleading terms in $`1/p_T`$ or higher twists have to be taken into account. This can be calculated in the framework of factorization approach by a careful evaluation of complete set of diagrams in the given $`\alpha _s`$-order, $`O(\alpha _s^4)`$. The non-fragmentational contributions dominate at $`p_Tm_{B_c}`$ .
The measured $`B_c`$ lifetime
$$\tau [B_c]=0.46_{0.16}^{+0.18}\pm 0.03\mathrm{ps},$$
agrees with the estimates obtained in the framework of both the OPE combined with the NRQCD evaluation of hadronic matrix elements and potential quark models, where one has to sum up the dominating exclusive modes to calculate the total $`B_c`$ width
$$\tau _{\mathrm{OPE},\mathrm{PM}}[B_c]=0.55\pm 0.15\mathrm{ps}.$$
The accurate measurement of $`B_c`$ lifetime could allow one to distinguish various parameter dependencies such as the optimal heavy quark masses, which basically determine the theoretical uncertainties in OPE.
At present, the calculations of $`B_c`$ decays in the framework of QCD sum rules were performed in . The authors of got the results, where the form factors are about 3 times less than the values expected in the potential quark models, and the semileptonic and hadronic widths of $`B_c`$ are one order of magnitude less than those in OPE. The reason for such the disagreement was pointed out in and studied in : in the QCD sum rules for the heavy quarkonia the Coulomb-like corrections are significant, since they correspond to summing up the ladder diagrams, where $`\alpha _s/v`$ is not a small parameter, as the heavy quarks move nonrelativistically, $`v1`$. The Coulomb rescaling of quark-quarkonium vertex enhances the estimates of form factors in the QCD sum rules for the $`B_c^+\psi (\eta _c)l^+\nu `$ decays, where the initial and recoil mesons are both the heavy quarkonia. In the framework of NRQCD at the recoil momentum close to zero one derives the spin symmetry relations for the form factors of semileptonic $`B_c`$ decays . In the strict limit of $`v_1=v_2`$, where $`v_{1,2}`$ denote the four-velocities of initial and recoil mesons, respectively, the authors of found a single relation between the form factors <sup>4</sup><sup>4</sup>4In refs. the relations were studied in the framework of potential models.. In the soft limit $`v_1v_21`$ at $`v_1v_2`$ was considered, and the generalized spin symmetry relations were obtained for the $`B_c\psi (\eta _c)`$ transitions: four equations, including that of . Moreover, the gluon condensate term was calculated in both QCD and NRQCD, so that it enforced a convergency of the method.
In the present paper we calculate the $`B_c`$ decays due to the $`cs`$ weak transition in the framework of QCD sum rules, taking into account the Coulomb-like $`\alpha _s/v`$-corrections for the heavy quarkonium in the initial state. In the semileptonic decays the hadronic final state is saturated by the pseudoscalar $`B_s`$ and vector $`B_s^{}`$ mesons, so that we need the values of their leptonic constants entering the sum rules and determining the normalization of form factors. For this purpose, we reanalyze the two-point sum rules for the $`B`$ mesons to take into account the product of quark and gluon condensates in addition to the previous consideration of terms with the quark and mixed condensates. We demonstrate the significant role of the product term for the convergency of method and reevaluate the constants $`f_B`$ as well as $`f_{B_s}`$. Taking into account the dependence on the threshold energy $`E_c`$ of hadronic continuum in the $`\overline{b}s`$ system in both the value of $`f_{B_s}`$ extracted from the two-point sum rules and the form factors in the three-point sum rules, we observe the stability of form factors versus $`E_c`$, which indicates the convergency of sum rules.
The spin symmetries of leading terms in the lagrangians of HQET for the singly heavy hadrons (here $`B_s^{()}`$) and NRQCD for the doubly heavy mesons (here $`B_c`$) result in the relations between the form factors of semileptonic $`B_cB_s^{()}`$ decays. We derive two generalized relations in the soft limit $`v_1v_21`$: one equation in addition to what was found previously in ref.. The relations are in a good agreement with the sum rules calculations up to the accuracy better than 10%, that shows a low contribution of next-to-leading $`1/m_Q`$-terms.
We perform the numerical estimates of semileptonic $`B_c`$ widths and use the factorization approach to evaluate the hadronic modes. Summing up the dominating exclusive modes, we calculate the lifetime of $`B_c`$, which agree with the experimental data and the predictions of OPE and quark models. We discuss the preferable prescription for the normalization point of nonleptonic weak lagrangian for the charmed quark and present our optimal estimate of total $`B_c`$ width. We stress that in the QCD sum rules to the given order in $`\alpha _s`$, the uncertainty in the values of heavy quark masses is much less than in OPE. This fact leads to a more definite prediction on the $`B_c`$ lifetime.
The paper is organized as follows. Section 2 is devoted to the general formulation of three-point sum rules for the $`B_c`$ decays with account of Coulomb-like corrections. The analysis of two-point sum rules for the leptonic constant of singly heavy meson with the introduction of term allowing for the product of quark and gluon condensates is presented in Section 3, where we also show the convergency of three-point sum rules with respect to a dependence on the threshold energy of continuum in the heavy-light system. We estimate the form factors of semileptonic $`B_cB_s^{()}`$ decays. The relations between the form factors of semileptonic decays as follows from the spin symmetry of HQET and NRQCD are derived in Section 4 in the soft limit of zero recoil. Section 5 contains the description how the nonleptonic decays modes are calculated and the $`B_c`$ lifetime is evaluated. We discuss the optimal estimation of lifetimes for the heavy hadrons in Section 6. In Conclusion we summarize the results.
## 2 Three-point sum rules for the $`B_c`$ meson.
In this paper the approach of three-point QCD sum rules is used to study the form factors of semileptonic and nonleptonic decay rates for the c $``$ s transition in decays of $`B_c`$ meson. From the two-point sum rules we extract the values for the leptonic constants of mesons in the initial and final states. In our consideration we use the following notations:
$$0|\overline{q}_1i\gamma _5q_2|P(p)=\frac{f_PM_P^2}{m_1+m_2},$$
(1)
and
$$0|\overline{q}_1i\gamma _\mu q_2|V(p,ϵ)=iϵ_\mu M_Vf_V,$$
(2)
where P and V represent the scalar and vector mesons, and $`m_1`$, $`m_2`$ are the quark masses.
The hadronic matrix elements for the semileptonic $`cs`$transition in the $`B_c`$ decays can be written down as follows:
$`B_s(p_2)|V_\mu |B_c(p_1)`$ $`=`$ $`f_+(p_1+p_2)_\mu +f_{}q_\mu ,`$ (3)
$`{\displaystyle \frac{1}{i}}B_s^{}(p_2)|V_\mu |B_c(p_1)`$ $`=`$ $`iF_Vϵ_{\mu \nu \alpha \beta }ϵ^\nu (p_1+p_2)^\alpha q^\beta ,`$ (4)
$`{\displaystyle \frac{1}{i}}B_s^{}(p_2)|A_\mu |B_c(p_1)`$ $`=`$ $`F_0^Aϵ_\mu ^{}+F_+^A(ϵ^{}p_1)(p_1+p_2)_\mu +F_{}^A(ϵ^{}p_1)q_\mu ,`$ (5)
where $`q_\mu =(p_1p_2)_\mu `$ and $`ϵ^\mu =ϵ^\mu (p_2)`$ is the polarization vector of $`B_s^{}`$ meson. $`V_\mu `$ and $`A_\mu `$ are the flavour changing vector and axial electroweak currents. The form factors $`f_\pm ,F_V,F_0^A`$ and $`F_\pm ^A`$ are functions of $`q^2`$ only. It should be noted that since the leptonic current $`l_\mu =\overline{l}\gamma _\mu (1+\gamma _5)\nu _l`$ is transversal in the limit of massless leptons, the probabilities of semileptonic decays are independent of $`f_{}`$ and $`F_{}^A`$ (the $`\tau ^+\nu _\tau `$ mode is forbidden by the energy conservation). Following the standard procedure for the evaluation of form factors in the framework of QCD sum rules , we consider the three-point functions
$`\mathrm{\Pi }_\mu (p_1,p_2,q^2)`$ $`=`$ $`i^2{\displaystyle }dxdye^{i(p_2xp_1y)}`$ (6)
$`0|T\{\overline{q}_1(x)\gamma _5q_2(x),V_\mu (0),\overline{b}(y)\gamma _5c(y)\}|0,`$
$`\mathrm{\Pi }_{\mu \nu }^{V,A}(p_1,p_2,q^2)`$ $`=`$ $`i^2{\displaystyle }dxdye^{i(p_2xp_1y)}`$ (7)
$`0|T\{\overline{q}_1(x)\gamma _\mu q_2(x),J_\nu ^{V,A}(0),\overline{b}(y)\gamma _5c(y)\}|0,`$
where $`\overline{q}_1(x)\gamma _5q_2(x)`$ and $`\overline{q}_1(x)\gamma _\nu q_2(x)`$ denote interpolating currents for $`B_s`$ and $`B_s^{}`$, correspondingly. $`J_\mu ^{V,A}`$ are the currents $`V_\mu `$ and $`A_\mu `$ of relevance to the various cases.
The Lorentz structures in the correlators can be written down as
$`\mathrm{\Pi }_\mu `$ $`=`$ $`\mathrm{\Pi }_+(p_1+p_2)_\mu +\mathrm{\Pi }_{}q_\mu ,`$ (8)
$`\mathrm{\Pi }_{\mu \nu }^V`$ $`=`$ $`i\mathrm{\Pi }_Vϵ_{\mu \nu \alpha \beta }p_2^\alpha p_1^\beta ,`$ (9)
$`\mathrm{\Pi }_{\mu \nu }^A`$ $`=`$ $`\mathrm{\Pi }_0^Ag_{\mu \nu }+\mathrm{\Pi }_1^Ap_{2,\mu }p_{1,\nu }+\mathrm{\Pi }_2^Ap_{1,\mu }p_{1,\nu }+\mathrm{\Pi }_3^Ap_{2,\mu }p_{2,\nu }+\mathrm{\Pi }_4^Ap_{1,\mu }p_{2,\nu }.`$ (10)
The form factors $`f_\pm `$, $`F_V`$, $`F_0^A`$ and $`F_\pm ^A`$ are determined from the amplitudes $`\mathrm{\Pi }_\pm `$, $`\mathrm{\Pi }_V`$, $`\mathrm{\Pi }_0^A`$ and $`\mathrm{\Pi }_\pm ^A=\frac{1}{2}(\mathrm{\Pi }_1^A\pm \mathrm{\Pi }_2^A)`$, respectively. In (8)-(10) the scalar amplitudes $`\mathrm{\Pi }_i`$ are the functions of kinematical invariants, i.e. $`\mathrm{\Pi }_i=\mathrm{\Pi }_i(p_1^2,p_2^2,q^2)`$.
The leading QCD term is a triangle quark-loop diagram in Fig. 1, for which we can write down the double dispersion representation at $`q^20`$
$$\mathrm{\Pi }_i^{pert}(p_1^2,p_2^2,q^2)=\frac{1}{(2\pi )^2}\frac{\rho _i^{pert}(s_1,s_2,Q^2)}{(s_1p_1^2)(s_2p_2^2)}𝑑s_1𝑑s_2+\text{subtractions},$$
(11)
where $`Q^2=q^20`$. The integration region in (11) is determined by the condition
$$1<\frac{2s_1s_2+(s_1+s_2q^2)(m_b^2m_c^2s_1)}{\lambda ^{1/2}(s_1,s_2,q^2)\lambda ^{1/2}(m_c^2,s_1,m_b^2)}<1,$$
(12)
where $`\lambda (x_1,x_2,x_3)=(x_1+x_2x_3)^24x_1x_2`$. The expressions for spectral densities $`\rho _i^{pert}(s_1,s_2,Q^2)`$ are given in Appendix A.
Now let us proceed with the physical part of three-point sum rules. The connection to hadrons in the framework of QCD sum rules is obtained by matching the resulting QCD expressions of current correlators with the spectral representation, derived from a double dispersion relation at $`q^20`$.
$$\mathrm{\Pi }_i(p_1^2,p_2^2,q^2)=\frac{1}{(2\pi )^2}\frac{\rho _i^{phys}(s_1,s_2,Q^2)}{(s_1p_1^2)(s_2p_2^2)}𝑑s_1𝑑s_2+\text{subtractions}.$$
(13)
Assuming that the dispersion relation (13) is well convergent, the physical spectral functions are generally saturated by the ground hadronic states and a continuum starting at some effective thresholds $`s_1^{th}`$ and $`s_2^{th}`$
$`\rho _i^{phys}(s_1,s_2,Q^2)`$ $`=`$ $`\rho _i^{res}(s_1,s_2,Q^2)+`$
$`\theta (s_1s_1^{th})\theta (s_2s_2^{th})\rho _i^{cont}(s_1,s_2,Q^2),`$
where the resonance term is expressed through the product of leptonic constant and form factor for the transition under consideration, so that
$`\rho _i^{res}(s_1,s_2,Q^2)`$ $`=`$ $`0|\overline{b}\gamma _5(\gamma _\mu )s|B_s(B_s^{})B_s(B_s^{})|F_i(Q^2)|B_cB_c|\overline{b}\gamma _5c|0)`$ (15)
$`(2\pi )^2\delta (s_1M_1^2)\delta (s_2M_2^2)+\text{higher state contributions},`$
where $`M_{1,2}`$ denote the masses of hadrons in the initial and final states. The continuum of higher states is modelled by the perturbative absorptive part of $`\mathrm{\Pi }_i`$, i.e. by $`\rho _i`$. Then, the expressions for the form factors $`F_i`$ can be derived by equating the representations for the three-point functions $`\mathrm{\Pi }_i`$ in (11) and (13), which means the formulation of sum rules.
For the heavy quarkonium $`\overline{b}c`$, where the relative velocity of quark movement is small, an essential role is taken by the Coulomb-like $`\alpha _s/v`$-corrections. They are caused by the ladder diagram, shown in Fig. 2. It is well known that an account for this corrections in two-point sum rules numerically leads to a double-triple multiplication of Born value of spectral density . In our case it leads to the finite renormalization for $`\rho _i`$ , so that
$$\rho _i^c=𝐂\rho _i,$$
(16)
$$𝐂=\frac{|\mathrm{\Psi }_{\overline{b}c}^C(0)|}{|\mathrm{\Psi }_{\overline{b}c}^{free}(0)|}=\sqrt{\frac{4\pi \alpha _s}{3v}(1\mathrm{exp}\{\frac{4\pi \alpha _s}{3v}\})^1},$$
(17)
where $`v`$ is the relative velocity of quarks in the $`\overline{b}c`$-system,
$$v=\sqrt{1\frac{4m_bm_c}{p_1^2(m_bm_c)^2}}.$$
(18)
To the moment, the procedure of calculations is completely described.
## 3 Numerical results on the form factors and the semileptonic decay widths
We evaluate the form factors in the scheme of spectral density moments. This scheme is not strongly sensitive to the value of the $`\overline{b}c`$-system threshold energy. In our calculations $`E_c^{\overline{b}c}=1.2\text{GeV}`$. The two-point sum rules for the $`B_c`$ meson with account for the Coulomb-like corrections give $`\alpha _s^c(\overline{b}c)`$=0.45, which corresponds to $`f_{B_c}`$=400 MeV . The quark masses are fixed by the calculations of leptonic constants $`f_\mathrm{\Psi }`$ and $`f_\mathrm{{\rm Y}}`$ in the same order over $`\alpha _s`$. The requirement of stability in the sum rules including the contributions of higher excitations, results in quite an accurate determination of masses $`m_c=1.40\pm 0.03`$GeV and $`m_b=4.60\pm 0.02`$GeV, which are in a good agreement with the recent estimates in , where the quark masses free off a renormalon ambiguity were introduced. The values of leptonic constants $`f_\mathrm{\Psi }`$, $`f_\mathrm{{\rm Y}}`$ linearly depend on the Coulomb-exchange $`\alpha _s`$. We find $`\alpha _s^c(\overline{c}c)0.60`$, $`\alpha _s^c(\overline{b}b)0.37`$, which obey the remormalization group evolution with the appropriate scale prescription, depending on the quarkonium contents. In this way, we can extract the above values of $`\alpha _s^c(\overline{b}c)`$ and $`f_{B_c}`$. Note, that the heavy ($`Q_1\overline{Q}_2`$)-quarkonia constants obey the scaling relation
$$\frac{f_n^2}{M_n}\left(\frac{M_n(m_1+m_2)}{4m_1m_2}\right)^2=\frac{c}{n},$$
(19)
where $`n`$ denotes the radial excitation number of nS-level, and $`c`$ is independent of heavy quark flavors.
The leptonic constant for the $`B_s`$ meson is extracted from the two-point sum rules. The Borel improved sum rules for the $`B`$ meson leptonic constant have the following form:
$$f_B^2M_Be^{\overline{\mathrm{\Lambda }}(\mu )\tau }=K^2\frac{3}{\pi ^2}C(\mu )\underset{0}{\overset{\omega _0(\mu )}{}}𝑑\omega \omega ^2e^{\omega \tau }+\overline{q}q(1\frac{m_0^2\tau ^2}{16}+\frac{\pi ^2\tau ^4}{288}\frac{\alpha _s}{\pi }G^2),$$
(20)
where we use $`\overline{q}q=(0.23\text{GeV})^3`$, $`m_0^2=0.8\text{GeV}^2`$, $`\frac{\alpha _s}{\pi }G^2=1.7710^2\text{GeV}^4`$ as the central values, and $`M_B`$=5.28 GeV. The K-factor is due to $`\alpha _s`$-corrections. We expect it is large, but we suppose the appearance of the same factor in evaluating the $`\alpha _s`$-corrections to the heavy-light vertex in the triangle diagram. For this factor we have the following expression :
$$K^2=\left\{\underset{0}{\overset{E_c\tau }{}}z^2e^z𝑑z\right\}^1\underset{0}{\overset{E_c\tau }{}}z^2e^z\left\{1+\frac{2\alpha _s(\stackrel{~}{\mathrm{\Lambda }})}{\pi }\left(\frac{13}{6}+\frac{2\pi ^2}{9}\text{ln}z\right)\right\}𝑑z,$$
(21)
where we suppose that the scale $`\stackrel{~}{\mathrm{\Lambda }}`$ is equal to 1.25 GeV. The dependence of K-factor on the Borel parameter $`\tau `$ and the threshold energy $`E_c`$ is shown in Fig. 3. The K-factor is not sensitive to $`E_c`$ changing in the range $`1.0÷1.5`$GeV.
We see that NLO corrections to the leptonic constant are about $`40\%`$. Using the Padé approximation, we find that higher orders corrections can be about $`30\%`$. So, we hold the K factor in conservative limits $`1.4÷1.7`$. It is quite reasonable to suppose its cancellation in evaluating the semileptonic form factors due to the renormalization of heavy-light vertex in the triangle diagram.
The contribution of quark condensate term is not sensitive to the variation of $`\overline{q}q`$ in the limits from -$`(0.23\text{GeV})^3`$ to -$`(0.27\text{GeV})^3`$ (this variation corresponds to the renormalization group evolution and the insertion of $`\alpha _s`$-corrections to this term, so that $`K`$ can be putted as the overall factor).
In the limit of semi-local duality $`\tau 0`$ we get the relation: $`\overline{\mathrm{\Lambda }}(\mu )=\frac{3}{4}\omega _0(\mu )`$ (the contribution of the quark condensate term to this equation is about $`15\%`$ ). We introduce the renormalization invariant quantities
$$\omega _{0,dual}^{ren}=C^{1/3}(\mu )\omega _0(\mu ),\overline{\mathrm{\Lambda }}_{dual}^{ren}=\frac{3}{4}\omega _{0,dual}^{ren}.$$
For $`\overline{\mathrm{\Lambda }}_{dual}^{ren}`$ we have $`\overline{\mathrm{\Lambda }}_{dual}^{ren}=M_Bm_b=0.63\text{GeV}`$, and we obtain that in the semi-local duality the threshold energy $`\omega _{0,dual}^{ren}=0.84\text{GeV}`$. Neglecting the quark condensate term in the leptonic constant we have
$$f_B^2M_B=K^2\frac{3}{\pi ^2}(\omega _{0,dual}^{ren})^3.$$
(22)
Since in the three-point sum rules we use the scheme of moments and search for a stable region, in the general Borel scheme for $`f_B`$ we have to consider the stability at $`\tau 0`$ with the extended region of resonance contribution. We expect, that the sum rules in (20) with the redefined $`\omega ^{ren}`$ and $`\overline{\mathrm{\Lambda }}^{ren}`$, as mentioned, have a stability point at $`\tau \frac{1}{\overline{\mathrm{\Lambda }}}`$
$$f_B^2M_Be^{\overline{\mathrm{\Lambda }}\tau }=K^2\frac{3}{\pi ^2}\underset{0}{\overset{E_c}{}}𝑑\omega \omega ^2e^{\omega \tau }+\overline{q}q(1\frac{m_0^2\tau ^2}{16}+\frac{\pi ^2\tau ^4}{288}\frac{\alpha _s}{\pi }G^2),$$
(23)
where $`E_c`$ is already not equal to $`\omega _{0,dual}^{ren}`$. Demanding a low deviation of $`\overline{\mathrm{\Lambda }}`$ from $`\overline{\mathrm{\Lambda }}^{ren}=0.63`$ GeV, we find that sum rules in Eq.(23) can lead to the results, which are in a good agreement with the semi-local duality if $`E_c=1.1÷1.3`$GeV (see Fig. 4). Then the optimal value of Borel parameter $`\tau =\tau _m6.5\text{GeV}^1`$. We write down
$$f_B^2M_Be^{\overline{\mathrm{\Lambda }}\tau _m}=K^2\frac{3}{\pi ^2}RE_c^3+\overline{q}q(1\frac{m_0^2\tau _m^2}{16}+\frac{\pi ^2\tau _m^4}{288}\frac{\alpha _s}{\pi }G^2),$$
(24)
where $`R`$ denotes the average value of $`e^{\omega \tau _m}`$. So, we find the $`E_c^{3/2}`$-dependence of $`f_B\sqrt{M_B}`$, whereas the contribution of condensate is numerically suppressed, as expected from the semi-local duality. The results of general Borel scheme calculations of $`f_B\sqrt{M_B}`$ ignoring the overall $`K`$-factor are presented in Fig. 4. We observe two stability regions. The stability region at $`\tau =2÷4`$ corresponds to that of considered in . The results for the leptonic constant $`f_B`$ obtained from this region is about 1.5 greater than the value obtained from the stability region at $`\tau =6÷7`$. The second region appears only when we introduce the term with the product of quark and gluon condensates. The similar situation has been observed in the NRQCD sum rules for doubly heavy baryons . The product of quark and gluon condensates was not taken into account in , and therefore, the intermediate stability point was observed only. Fixing the optimal values of $`f_B^2M_B`$ in Eq.(24) from Fig. 4, we can invert the sum rules to study the dependence of $`\overline{\mathrm{\Lambda }}`$ on $`\tau `$, as shown in Fig. 5, where the optimal values of $`\overline{\mathrm{\Lambda }}`$ agree with the semi-local duality and the estimate $`\overline{\mathrm{\Lambda }}^{ren}=M_Bm_b`$. Note that the intermediate stability point $`\tau 4`$ exhibit a low variation of $`\overline{\mathrm{\Lambda }}`$ close to 0.4 GeV, which was obtained in , and usually given by the potential models (see, for instance, ).
The calculated physical quantity should be independent of parameters in the sum rule scheme. However, we truncate both the operator product expansion and the perturbative series for the Wilson coefficients by fixed orders in the opreator dimension and $`\alpha _s`$, respectively. This fact leads to that the results depend on the scheme parameters, say, the Borel variable. Moreover, the physical part of sum rules is modelled by the contributions of resonances and a continuum term starting at a threshold, that introduces the dependence on the threshold value and suggests that the stability of results can be improved by the terms of excited resonaces in addition to the ground state.
Let us, at first, consider the results on the leptonic constant in the limit of semi-local duality, which requires the stability at $`\tau 0`$ and corresponds to the duality region containing the ground state, only. The sum rules show that the condensate contributions are given by the polinomials over $`\tau `$, and the leading correction at $`\tau 0`$ is the term with the quark condensate. Then, we expect that the region of stability will extend at $`\tau >0`$ if we will add the higher condensates with the appropriate ratios of their values. The central values of condensates as mentioned above correspond to the results shown in Fig. 4. The stability of semi-local sum rules can be improved by a variation of $`\overline{\mathrm{\Lambda }}`$ and $`\omega _0`$, which is not important for the current discussion. In order to clarify this statement we present the results of semi-local sum rules for the leptonic constant of $`B`$ meson in Fig. 6, where we put $`\overline{\mathrm{\Lambda }}=0.54`$ GeV. Note, that the value of leptonic constant is the same as we have found in the pertubative limit of semi-local duality. Then, we state that the value of leptonic constant obtained at $`\tau 0`$ agrees with the value corresponding to the stability at $`\tau 7`$ GeV<sup>-1</sup>, i.e. in the second point of local extremum in Fig. 4. In order to confirm, we present also the result for the other ratio of condensate values (the lower value of mixed condensate and the upper value of gluon condensate in the regions mentioned below) in Fig. 7. We see that the semi-local duality is completely broken at unappropriate choice of condensate values.
Second, we study the general Borel sum rules in the scheme, where we do not require the stability at $`\tau 0`$, which means that we extend the duality region to incorporate possible excitation contributions. Then, we expect that the estimate of leptonic constant should result in the same value obtained in the semi-local duality. We see that this can be reached at the same values of condensates involved as well as in the same region of $`\tau `$, which is again given by the second extremum (see Figs. 4 and 6). Certainly, the extension of duality region leads to a failure of stability at low $`\tau `$, which is caused by additional terms given by the excitations beyond the control of the method.
At higher values of $`\tau `$ the stability would be reached by introduction of higher condensates contributing as the polinomials of higher powers. Of course, we can achieve the better stability in the general Borel scheme by adjusting the condensate values, but this will destroy the semi-local duality, that indicates the divergency of the method, while the convergency demands an appropriate choice of condensate values as we put.
The criterion on the covergency of both the semi-local duality and general Borel sum rules was ignored in , where the greater value of leptonic constant was obtained (see Fig. 7).
The same notes can be done in the discussion on $`\overline{\mathrm{\Lambda }}`$. For the sake of comparison, we present the results in Figs. 5 and 8, corresponding to the inverted sum rules for the leptonic constant given in Figs. 4 and 6.
The physical meaning of gluon condensate as it stands in the operator product expansion taken between the observed states is independent of the scheme of calculations. In the sum rules the gluon condensate contributes in the region, where the excited states are suppressed, while the ground state corresponding to the current under consideration dominates. So, the condensate essentially determines the binding energy of quarks in the hadron containing the heavy quark. The definition of sum rules scheme includes the parameter determining the threshold energy of continuum contribution in addition to the resonance term. So, the difference between the schemes of semi-local duality and usual Borel representation is due to the variation of duality region. So, the semi-local duality means the duality for the region contaning the ground state only, while the usual Borel (or moment) scheme explore the extended region containing several hadronic states: the ground state and its excitations. However, the gluon condensate essentially contributes at the scheme parameters, when the excitations are suppressed. Thus, its value is the same for both schemes used, i.e. for the semi-local duality and in the Borel scheme. The value of gluon condensate has been varied in the range $`\frac{\alpha _s}{\pi }G^2=(1.5÷2)10^2\text{GeV}^4`$. As for the mixed condesate we have used the range $`m_0^2=(0.72÷0.88)\text{GeV}^2`$. The variation of these parameters does not change the qualitative picture for the lepton constant as it has been discussed, while the numerical uncertainty of its value is less than 7%.
Numerically, multiplying the result taken from Fig. 4, by the K-factor we find the value $`f_B=140÷170`$MeV, which is in a good agreement with the recent lattice results and the estimates in the QCD SR by other authors . So, we can conclude that the $`1/m_b`$-corrections are not valuable for $`f_B`$. The uncertainty of estimates is basically connected with the higher orders in $`\alpha _s`$. For the vector $`B^{}`$ meson constant $`f_B^{}`$ we put $`\frac{f_B^{}}{f_B}=1.11`$(see ).
For the leptonic constant of $`B_s`$ meson we explore the following relation $`\frac{f_{B_s}}{f_B}=1.16`$, which expresses the flavor SU(3)-symmetry violation for B mesons .
Remember, in sum rules the heavy quark masses are fixed by the two-point sum rules for bottomonia and charmonia with the precision of 20 MeV. In our consideration the quark masses are equal to $`m_b`$=4.6 GeV, $`m_c`$=1.4 GeV, and we use $`m_s`$=0.15 GeV, which agrees with the various estimates . The uncertainties in the values of form factors are basically determined by the variation of $`b`$-quark mass, while changing the other quark masses in the ranges $`m_s=0.14÷0.16`$ GeV and $`m_c=1.35÷1.45`$ GeV, results in the uncertainty less than 2%. In Figs. 9, 10 and 11 we present the results in the scheme of spectral density moments.
We have investigated the dependence of form factors on the $`\overline{b}s`$ threshold energy of continuum in the range $`E_c=1.1÷1.3`$ GeV. The characteristic forms of this dependence are shown in Figs. 12 and 13. We see that the optimal choice for the $`\overline{b}s`$ system threshold energy is 1.2 GeV. In Table 1 we present the results of sum rules for the form factors in comparison with estimates in the framework of potential models . We see a good agreement of estimates in the QCD sum rules with the values in the quark model. For the sake of completeness the quark model expressions for the form factors are given in Appendix B.
In the form factors were derived using the similar SR technique but without the Coulomb-like corrections in the $`\overline{b}c`$ system, which enhance the form factors about three times, as we have found.
Let us discuss the uncertainties in the sum rules and other approaches. So, the potential models with similar choices of parameters result in the form factor values, which are slightly model-dependent. The corresponding accuracy is about 10%. Then, we expect that the potential models give good reference points for the appropriate numerical values.
The accuracy of sum rules under consideration is basically determined by the variation of heavy quark masses. Indeed, the significant $`\alpha _s`$ correction to the leptonic constant of $`B_s`$ meson should cancel the same factor for the renormalization of quark-meson vertex in the triangle diagram. The dependence on the choice of threshold energy in the $`\overline{b}s`$-channel can be optimized and, hence, minimized. The variation of threshold energy in the $`\overline{b}c`$-channel give the error less than 1%. The effective coulomb constant is fixed from the two-point sum rules for the heavy quarkonium, and its variation is less than 2%, which gives the same uncertainty for the form factors. The heavy quark masses are determined by the two-point sum rules for the heavy quarkonia, too. However, their variations result in the most essential uncertainty. Summing up all of mentioned variations we estimate $`\delta f/f5`$%.
The enhancment of form factors in the decays of heavy quarkonium considered in the framework of sum rules was also found in the decays $`B_c\psi l\nu `$ . It is important to stress that the current consideration removes the contradiction between the estimates in OPE, potential models and the values in the SR. We have calculated the form factors in the SR consistent with the values in the OPE and potential models.
We use the three-point sum rules to determine the dependence on $`q^2`$ of form factors in vicinity of $`q^2=0`$, where the method works even in the approximation by a bare quark loop . Naively, we expect a simple pole form
$$f(q^2)=\frac{f(0)}{1\frac{q^2}{M_{pole}^2}},$$
(25)
so that the first derivative $`f^{}(0)=f(0)\frac{1}{M_{pole}^2}`$ can be evaluated in the framework of sum rules to estimate the “pole” mass $`M_{pole}`$, which may deviate from the value of physical mass of corresponding bound state. Of course, the bare quark loop approximation cannot be justified at $`q^2m_{pole}^2`$, while the three-point sum rules in the form of double dispersion relation cannot be explored in the region of resonances at $`q^2>0`$ because of so-called “non–Landau” singularities indicating the presence of strongly bound states in the $`q^2`$-channel. That is why we calculate the form factors as well as their Borel transforms at $`q^2<0`$.
Numerically, we have found $`M_{pole}`$=1.3$`÷`$1.4 GeV for the form factors with the $`B_s^{}`$ decay modes, and $`M_{pole}`$=1.8$`÷`$1.9 GeV for the decay form factors with the $`B_s`$ modes.
The semileptonic widths are presented in Table 2. We have supposed the quark mixing matrix element $`|V_{cs}|`$=0.975 . The mesons masses are equal to $`M_{B_c}`$=6.25 GeV, $`M_{B_s}`$=5.37 GeV, $`M_{B_s^{}}`$=5.41 GeV .
These results agree with the values obtained in the framework of covariant quark model : $`\mathrm{\Gamma }`$($`B_se^+\nu _e`$)=$`4.710^{14}`$GeV, $`\mathrm{\Gamma }`$($`B_s^{}e^+\nu _e`$)=$`7.410^{14}`$GeV, as we could expect looking at Table 1.
## 4 The symmetry relations
At the recoil momentum close to zero, the heavy quarks in both the initial and final states have small relative velocities inside the hadrons, so that the dynamics of heavy quarks is essentially nonrelativistic. This allows us to use the combined NRQCD/HQET approximation in the study of mesonic form factors. The expansion in the small relative velocities to the leading order leads to various relations between the different form factors. Solving these relations results in the introduction of an universal form factor (an analogue of the Isgur-Wise function) at $`q^2q_{max}^2`$.
We consider the soft limit
$`v_1^\mu v_2^\mu ,`$
$`w=v_1v_21,`$ (26)
where $`v_{1,2}^\mu =p_{1,2}^\mu /\sqrt{p_{1,2}^2}`$ are the four-velocities of heavy mesons in the initial and final states. The study of region (26) is reasonable enough, because in the rest frame of $`B_c`$meson ($`p_1^\mu =(\sqrt{p_1^2},\stackrel{}{0})`$), the four-velocities differ only by a small value $`|\stackrel{}{p}_2|`$ $`(p_2^\mu =(E_2,\stackrel{}{p}_2)`$, whereas their scalar product $`w`$ deviates from unity only due to a term, proportional to the square of $`|\stackrel{}{p}_2|`$: $`w=\sqrt{1+\frac{|\stackrel{}{p}_2|^2}{p_2^2}}1+\frac{1}{2}\frac{|\stackrel{}{p}_2|^2}{p_2^2}`$. Thus, in the linear approximation at $`|\stackrel{}{p}_2|0`$, relations (26) are valid and take place. Here we would like to note, that (26) generalizes the investigation of , where the case of $`v_1=v_2`$ was considered. This condition severely restricts the relations of spin symmetry for the form factors and, as a consequence, it provides a single connection between the form factors. In the soft limit of zero recoil we find
$$\stackrel{~}{v}_3^\mu =\frac{1}{2}(v_1^\mu +v_2^\mu )$$
(27)
for the four-velocity of spectator $`b`$-quark, and
$$\stackrel{~}{v}_1^\mu =v_1^\mu +\frac{m_3}{2m_1}(v_1v_2)^\mu $$
(28)
for the decaying $`c`$-quark. The matrix element of $`J_\mu =\overline{Q}_1\mathrm{\Gamma }_\mu q_2`$ with the spin structure $`\mathrm{\Gamma }_\mu =\{\gamma _\mu ,\gamma _5\gamma _\mu \}`$ has the form
$`H_{Q_1\overline{Q}_3}|J_\mu |H_{q_2\overline{Q}_3}`$ $`=`$ $`tr[\mathrm{\Gamma }_\mu (1+\stackrel{~}{v}_1^\mu \gamma _\mu )\mathrm{\Gamma }_1(1+\stackrel{~}{v}_3^\nu \gamma _\nu )`$ (29)
$`\mathrm{\Gamma }_2\rho _{light}]h,`$
where $`\mathrm{\Gamma }_1`$ determines the spin state in the heavy meson $`Q_1\overline{Q}_3`$ (in our case it is pseudoscalar, so that $`\mathrm{\Gamma }_1=\gamma _5`$), $`\mathrm{\Gamma }_2`$ determines the spin wave function of quarkonium in the final state $`\mathrm{\Gamma }_2=\{\gamma _5,ϵ^\mu \gamma _\mu \}`$ for the pseudoscalar and vector states, respectively ($`H=P,V`$). The ‘propagator of the light quark’<sup>5</sup><sup>5</sup>5More strictly, we determine the spin structure of matrix element and the term given by light degrees of freedom. is taken in a general form
$$\rho _{light}=1+B(\text{ / }v_2\text{ / }v_1)+C(\text{ / }v_2+\text{ / }v_1)+D\text{ / }v_2\text{ / }v_1,$$
(30)
where $`B,C,D`$ are the functions of $`w`$. The quantity $`h`$ in (29) at $`w1`$ is an universal factor independent of the spin state of meson. So, for the form factors, discussed in our paper, we have
$`P_{Q_1\overline{Q}_3}|\overline{Q}_1\gamma ^\mu Q_3|P_{q_2\overline{Q}_3}`$ $`=`$ $`(c_1^Pv_1^\mu +c_2^Pv_2^\mu )h,`$ (31)
$`P_{Q_1\overline{Q}_3}|\overline{Q}_1\gamma ^\mu Q_3|V_{q_2\overline{Q}_3}`$ $`=`$ $`ic_Vϵ^{\mu \nu \alpha \beta }ϵ_\nu v_{1\alpha }v_{2\beta }h,`$ (32)
$`P_{Q_1\overline{Q}_3}|\overline{Q}_1\gamma _5\gamma ^\mu Q_3|V_{q_2\overline{Q}_3}`$ $`=`$ $`(c_ϵϵ^\mu +c_1v_1^\mu (ϵv_1)+c_2v_2^\mu (ϵv_1))h,`$ (33)
where
$`c_ϵ`$ $`=`$ $`2,`$
$`c_V`$ $`=`$ $`1\stackrel{~}{B}{\displaystyle \frac{m_3}{2m_1}},`$ (34)
$`c_1^P`$ $`=`$ $`1\stackrel{~}{B}+{\displaystyle \frac{m_3}{2m_1}},`$
$`c_2^P`$ $`=`$ $`1+\stackrel{~}{B}{\displaystyle \frac{m_3}{2m_1}},`$
and $`\stackrel{~}{B}=\frac{B2D}{1+C}`$. The rest coefficients $`c_{1,2}`$ depend on the C and D parameters. We have the symmetry relations for the following form factors<sup>6</sup><sup>6</sup>6To remove an error in the analogous second relation for the $`B_cJ/\mathrm{\Psi }(\eta _c)`$ transition should have the missed factor 2 in front of $`c_ϵ`$.:
$`f_+(c_1^P_2c_2^P_1)f_{}(c_1^P_2+c_2^P_1)`$ $`=`$ $`0,`$
$`F_0^Ac_V2c_ϵF_V_1_2`$ $`=`$ $`0,`$ (35)
$`F_0^Ac_1^P+c_ϵ_1(f_++f_{})`$ $`=`$ $`0,`$
where $`_1=m_1+m_3`$, $`_2=m_2+m_3`$. Equating the second relation in (35), for example, we obtain
$$\stackrel{~}{B}=\frac{2m_1+m_3}{2m_1}+\frac{4m_3(m_1+m_3)F_V}{F_0^A}10.0,$$
(36)
where all form factors are taken at $`q_{max}^2`$. Substituting $`\stackrel{~}{B}`$ in first and third relations, we get $`f_+2.0`$ and $`f_{}8.3`$. These values have to be compared with the corresponding form factors obtained in the QCD sum rules: $`f_+(q_{max}^2)=1.8`$ and $`f_{}(q_{max}^2)=8.1`$, where we suppose the pole like behaviour of form factors (see Eq.(25)). Thus, we find that in the QCD sum rules, relations (35) are valid with the accuracy better than $`10\%`$ at $`q^2=q_{max}^2`$. The deviation could increase at $`q^2<q_{max}^2`$ because of variations in the pole masses governing the evolution of form factors. However, in $`B_c^+B_s^{()}l^+\nu `$ decays the phase space is restricted, so that the changes of form factors are about 50%, while their ratios develop more slowly.
## 5 Nonleptonic decays and the lifetime
The hadronic decay widths can be obtained on the basis of assumption on the factorization for the weak transition between the quarkonia and the final two-body hadronic states. For the dominant nonleptonic decay modes $`B_c^+B_s^{()}\pi ^+(\rho ^+)`$ the effective Hamiltonian can be written down as
$$H_{eff}=\frac{G_F}{2\sqrt{2}}V_{cs}V_{ud}^{}\{C_+(\mu )O_++C_{}(\mu )O_{}\},$$
(37)
where
$$O_\pm =(\overline{u}_i\gamma _\nu (1\gamma _5)d_i)(\overline{s}_j\gamma ^\nu (1\gamma _5)c_j)\pm (\overline{u}_i\gamma _\nu (1\gamma _5)d_j)(\overline{s}_i\gamma ^\nu (1\gamma _5)c_j),$$
(38)
where $`i,j`$ run over the colors. The factors $`C_\pm (\mu )`$ account for the strong corrections to the corresponding four-fermion operators caused by hard gluons. The review on the evaluation of $`C_\pm (\mu )`$ can be found in . In the present paper, dealing with the QCD sum rules in the leading order over $`\alpha _s`$, we explore the $`C_\pm (\mu )`$-evolution to the leading log accuracy. The $`B_c^+B_s\pi ^+`$ amplitude, for example, takes the form
$$A(B_c^+B_s\pi ^+)=\frac{G_F}{\sqrt{2}}V_{cs}V_{ud}a_1(\mu )\pi ^+|\overline{u}\gamma _\nu (1\gamma _5)d|0B_s|\overline{s}\gamma ^\nu (1\gamma _5)c|B_c,$$
(39)
where $`a_1(\mu )=\frac{1}{2N_c}(C_+(\mu )(N_c+1)+C_{}(\mu )(N_c1))`$ at $`N_c=3`$ being the number of colors. In our calculations we put the following light meson parameters: $`m_{\pi ^+}`$=0.14 GeV, $`m_{\rho ^+}`$=0.77 GeV, $`f_{\pi ^+}`$=0.132 GeV, $`f_{\rho ^+}`$=0.208 GeV. The results are collected in Table 3.
It is worth noting that the sum of widths for transitions $`B_c^+B_s(B_s^{})\pi ^+(\rho ^+)`$ is $`10\%`$ larger than the width for the transition $`B_c^+B_s(B_s^{})+lighthadrons`$, which is calculated using the simple formula
$$\mathrm{\Gamma }[B_c^+B_s(B_s^{})+lighthadrons]=N_ca_1^2(\mu )\mathrm{\Gamma }[B_c^+B_s(B_s^{})e^+\nu _e],$$
where we neglect the contributions given by the modes with the factor $`a_2`$ instead of $`a_1`$. In addition, the deviation between these estimates can be caused by the corresponding ‘bag’ factor appearing in the formulation of factorization approach and vacuum saturation in the connection of leptonic form factors to the hadronic ones. The modern lattice estimates show that the ‘bag’ parameters are about 7% less than 1 .
In the parton approximation we could expect
$$\mathrm{\Gamma }[B_c^+B_s^{()}+lighthadrons]=(2C_+^2(\mu )+C_{}^2(\mu ))\mathrm{\Gamma }[B_c^+B_s^{()}e^+\nu _e],$$
which results in the estimate very close to the value obtained as the sum of exclusive modes at $`\mu >0.9`$GeV. The deviation between these two estimates slightly increase at $`\frac{m_c}{2}<\mu <0.9`$GeV. Concerning the comparison of hadronic width summing up the exclusive decay modes with the estimate based on the quark-hadron duality, we insist that the deviation between these two estimates is less that 10% and, hence, it cannot be treated as an essential argument against the validity of our calculations.
We estimate the lifetime using the fact that the dominant modes of the $`B_c`$meson decays are the $`cs,bc`$ transitions with the $`B_s^{()}`$ and J/$`\psi `$, $`\eta _c`$ final states respectively, and the electroweak annihilation <sup>7</sup><sup>7</sup>7The $`\overline{b}\overline{c}c\overline{s}`$ transition is negligibly small in the $`B_c`$ decays because of destructive Pauli interference for the charmed quark in the initial state and the product of decay ..
We stress that in the $`B_cB_s`$ decays caused by the weak decays of charmed quark the possible hadronic final states are the charged mesons $`\pi `$, $`\rho `$ and $`K`$ or the multi-particle states like $`\pi \pi `$, $`\pi \pi \pi `$ or $`K\pi `$. First, the states with the kaon are suppressed by the Cabibbo angle, and we neglect their contributions in the total nonleptonic width of $`B_c`$ (the corresponding error of estimates is about 4%). The method for the calculation of multi-particle branching fractions was offered by Bjorken in his pioneering paper devoted to the decays of hadrons containing heavy quarks . He supposed a simple relation for the yields of pions, as given by the Poisson distribution with the average value of pions determined by the energy release. The Bjorken’s model of multi-particle yields in the decays does not distinguish the resonant and continuum final states as well as $`B_s`$ and $`B_s^{}`$. So, we check that the ratio of $`R_{2\pi }=\mathrm{\Gamma }(B_c^+B_s^{()}\pi ^+)/\mathrm{\Gamma }(B_c^+B_s^{()}\rho ^+)0.82`$ calculated in the framework of sum rules is close to the estimate $`R_{2\pi }=\mathrm{\Gamma }(B_c^+B_s^{()}\pi ^+)/\mathrm{\Gamma }(B_c^+B_s^{()}\pi ^+\pi ^0)0.85`$ given by Bjorken. Then, we see that the non-resonant multi-particle states are suppressed in comparison with the resonance yields. The same fact can be found, once we consider the $`K\pi `$, $`K\pi \pi `$ and $`K^{}\pi `$, $`K\rho `$ branching fractions in the decays of $`D`$ mesons as measured experimentally. This consideration confirms us that we take into account all of significant nonleptonic decay modes of $`B_c`$. In order to estimate the contribution of non-resonant $`3\pi `$ modes of $`B_c`$ decays into $`B_s^{()}`$ we use the Bjorken’s technique, i.e. the Poisson distribution with the average value corrected to agree with the non-resonant $`3\pi `$-modes in the decays of $`D`$ mesons. We have chosen the following branching ratios: $`\mathrm{BR}(D^+K^{}\pi ^+\pi ^+)=9.0\pm 0.6`$% and $`\mathrm{BR}(D^+K^{}\pi ^+\pi ^+\pi ^0|_{\mathrm{non}\mathrm{resonant}})=1.2\pm 0.6`$%. So, we have found $`\overline{n}\frac{1}{8}`$, which means that $`\mathrm{BR}(B_c^+B_s^{()}(3\pi )^+)0.2`$%, while $`\mathrm{BR}(B_c^+B_s^{()}(2\pi )^+|_{\mathrm{non}\mathrm{resonant}})3`$%. We see that the neglected modes contribute to the total width of $`B_c`$ as a small fraction in the limits of uncertainty envolved.
The width of beauty decay in the sum rules was derived using the similar methods in : $`\mathrm{\Gamma }(B_c^+\overline{c}c+X)=(28\pm 5)10^{14}`$GeV. The width of the electroweak annihilation is taken from as $`1210^{14}`$GeV.
In Fig. 14 we present the $`B_c`$meson lifetime calculated in the QCD SR under consideration. We also show the results of the lifetime evaluation in the framework of Operator Product Expansion in NRQCD .
In contrast to OPE, where the basic uncertainty is given by the variation of heavy quark masses, these parameters are fixed by the two-point sum rules for bottomonia and charmonia, so that the accuracy of SR calculations for the total width of $`B_c`$ is determined by the choice of scale $`\mu `$ for the hadronic weak lagrangian in decays of charmed quark. We show this dependence in Fig. 14, where $`\frac{m_c}{2}<\mu <m_c`$ and the dark shaded region corresponds to the scales preferred by data on the charmed meson lifetimes. The discussion on the optimal choice of scale in hadronic decays is addressed in the next section.
## 6 Discussion on the lifetimes of heavy hadrons
At present the ordinary prescription for the normalization point of lagrangian generating the nonleptonic decays of heavy quark $`Q`$, is $`\mu m_Q`$. The motivation is the following: the characteristic scale is determined by the energy release given by the heavy quark mass. Therefore, we can argue the operator product expansion in the inverse powers of $`m_Q`$, wherein we can factorize the Wilson coefficients taken in the perturbative QCD and the matrix elements of operators over the hadronic states with $`\mu `$ usually posed to $`m_Q`$. This prescription is in a qualitative agreement with the current data on the measured lifetimes of charmed and beauty hadrons and their branching ratios for the semileptonic decay modes, say.
Let us consider this issue in more details. To the moment, the analysis of decays in the QCD sum rules is restricted by the leading order (LO) in $`\alpha _s`$ (except the Coulomb-like corrections in the heavy quarkonia). The corresponding parameters, the heavy quark masses, are also fixed to the same order (they have to be reevaluated in next-to-leading order). Therefore, for the sake of consistency, we use the LO expressions, which are given by the partonic approximation improved by $`1/m_Q`$-corrections in HQET or NRQCD. In this way, we can write down the following formulae:
1. The semileptonic branching fraction of $`D^0`$ meson is given by
$$\mathrm{BR}_{sl}[D^0]=\frac{1}{2+2C_+^2(\mu )+C_{}^2(\mu )}.$$
(40)
2. The difference between the total widths of charmed mesons is determined by the Pauli interference in decays of $`D^+`$, so that
$$\mathrm{\Gamma }[D^+]\mathrm{\Gamma }[D^0]=\mathrm{cos}^4\theta _cG_\mathrm{F}^2\frac{m_c^3f_D^2}{8\pi }\left[(C_+^2(\mu )C_{}^2(\mu ))B+\frac{1}{3}(C_+^2(\mu )+C_{}^2(\mu ))\stackrel{~}{B}\right],$$
(41)
where $`\theta _c`$ is the Cabibbo angle, $`f_D`$ is the leptonic constant of $`D`$ meson, and the ‘bag’ constants are defined by
$`D|(\overline{c}\mathrm{\Gamma }_\mu q)(\overline{q}\mathrm{\Gamma }_\mu c)|D`$ $`=`$ $`f_D^2M_D^2B,`$ (42)
$`D|(\overline{c}\mathrm{\Gamma }_\mu c)(\overline{q}\mathrm{\Gamma }_\mu q)|D`$ $`=`$ $`{\displaystyle \frac{1}{3}}f_D^2M_D^2\stackrel{~}{B},`$ (43)
with $`\mathrm{\Gamma }_\mu =\gamma _\mu (1\gamma _5)`$.
The similar expression can be derived for the beauty mesons
$$\mathrm{\Gamma }[B^+]\mathrm{\Gamma }[B^0]=|V_{bc}|^2G_\mathrm{F}^2\frac{m_b^3f_B^2}{8\pi }\left[(C_+^2(\mu )C_{}^2(\mu ))B+\frac{1}{3}(C_+^2(\mu )+C_{}^2(\mu ))\stackrel{~}{B}\right].$$
(44)
3. The semileptonic branching fraction of $`B^0`$ meson is given by
$$\mathrm{BR}_{sl}[B^0]=\frac{1}{2+0.22+[2C_+^2(\mu )+C_{}^2(\mu )](1+k)},$$
(45)
where the fraction of $`0.22`$ is due to the $`\tau \nu _\tau `$-contribution, while the value of $`k`$ denotes the fraction of $`bc\overline{c}s`$ transition in the nonleptonic decays. In the same way, the average yield of charm in the decays of beauty mesons is equal to
$$n_c=\frac{2+0.22+[2C_+^2(\mu )+C_{}^2(\mu )](1+2k)}{2+0.22+[2C_+^2(\mu )+C_{}^2(\mu )](1+k)}.$$
(46)
As for numerical applications, one usually puts
$`\mu _D=m_c`$ in decays of $`D`$ mesons,
$`\mu _B=m_b`$ in decays of $`B`$ mesons,
$`f_Df_B200`$ MeV, and
$`B=\stackrel{~}{B}=1`$ naively motivated by nonrelativistic potential models.
At $`k=0.4`$ , this set (marked as SETMQ column) results in the estimates shown in Table 4 in comparison with the experimental data .
First, we note that the semileptonic width of $`D^0`$ is well described to the given order and at chosen value of $`m_c`$, while its branching ratio is in a valuable contradiction with the data indicating a more higher enforcement of nonleptonic modes. Second, the qualitative agreement of predictions with the measured differences of $`\mathrm{\Gamma }[D^+]\mathrm{\Gamma }[D^0]`$ and $`\mathrm{\Gamma }[B^+]\mathrm{\Gamma }[B^0]`$ is mainly based on the assumption of $`\stackrel{~}{B}1`$. Recent consideration of charmed baryon lifetimes by M.Voloshin clearly drawn a conclusion that the naive picture of color structure as given by the potential models (i.e. the purely antisymmetric color-composition of valence flavors) is significantly broken. A similar statement was obtained in the description of $`D^{}`$ meson production at HERA, where the authors of found that the $`O(\alpha _{em}\alpha _s^3)`$-calculations for the differential cross sections of $`c\overline{q}`$-pair composing the meson, are able to reproduce the measured spectra, if we introduce the valuable contribution by the color-octet state in addition to the singlet one. So, the four-quark singlet-operator results in
$$O_{(1)}=D^{}|(\overline{c}\gamma _\mu q)(\overline{q}\gamma _\mu c)|D^{}\left(g_{\mu \nu }+\frac{p_\mu p_\nu }{p^2}\right)\frac{1}{12M},$$
(47)
which is reduced to $`O_{(1)}=|\mathrm{\Psi }_{(1)}(0)|^2`$ in the framework of nonrelativistic potential model, where
$$|D^{}=\sqrt{2M}\frac{d^3q}{(2\pi )^3}\mathrm{\Psi }_{(1)}(q)\frac{\delta _{ij}}{\sqrt{3}}\frac{1}{\sqrt{2}}\overline{c}_i\text{ }ϵ\text{ }/q_j|0,$$
(48)
with $`ϵ_\alpha `$ denoting the polarization vector and $`\mathrm{\Psi }_{(1)}(q)`$ being the wave function.
The term of color-octet was parameterized by
$$O_{(8)}=D^{}|(\overline{c}\lambda ^a\gamma _\mu q)(\overline{q}\lambda ^a\gamma _\mu c)|D^{}\left(g_{\mu \nu }+\frac{p_\mu p_\nu }{p^2}\right)\frac{1}{64M},$$
(49)
which is, in a similar way, can be represented as $`O_{(8)}=|\mathrm{\Psi }_{(8)}(0)|^2`$, if we introduce the additional Fock state
$$\sqrt{2M}\frac{d^3q}{(2\pi )^3}\mathrm{\Psi }_{(8)}(q)\frac{\lambda _{ij}^a}{\sqrt{2}}n_a\frac{1}{\sqrt{2}}\overline{c}_i\text{ }ϵ\text{ }/q_j|0,$$
(50)
where the $`\overline{c}`$ and $`q`$ fields represent the rapid valence quarks, and $`n_a`$ is a random color-vector determined by soft degrees of freedom inside the meson (i.e. by the quark-gluon sea). We have $`n_an_b=\delta _{ab}`$ in the production, while $`n_an_b=\frac{1}{8}\delta _{ab}`$ in decays.
The data on the $`D^{}`$ production in DIS give $`O_{(8)}/O_{(1)}1.3`$. We can analogously expect that in the $`D`$ meson the ratio of four-quark matrix elements is close to that in $`D^{}`$. So, $`O_{(8)}[D]/O_{(1)}[D]1`$, where $`O_{(1,8)}[D]`$ can be obtained from the above expressions for $`D^{}`$ by the substitution of $`\gamma _5\gamma _\mu `$ for $`\gamma _\mu `$ and removing the transverse projector. Then we straightforwardly find
$$\stackrel{~}{B}=1+\frac{O_{(8)}[D]}{O_{(1)}[D]}2.$$
(51)
Putting $`\stackrel{~}{B}=2`$ into the SETMQ we get the values given in the SETMQ<sub>2</sub> column in Table 4. So, the choice of $`\mu =m_Q`$ is in a deep contradiction with the observed differences of total widths for the heavy mesons at the most reasonable value of $`\stackrel{~}{B}=2`$.
In this position we argue the following: There are other physical scales in the problem, which are characteristic for two hadronic systems in the decay process. The first system is the decaying hadron. The second is the transition current $`cs`$ or $`\overline{b}\overline{c}`$, where the form factor behaviour versus the transferred momentum is determined by the $`c\overline{s}`$ and $`\overline{b}c`$ states. Those hadronic systems have the following scales, being the average squares of heavy quark momentum $`p^2`$, which are phenomenologically equal to
$$\mu _{c\overline{u}}^2=\mu _{b\overline{d}}^2=2T\overline{\mathrm{\Lambda }},$$
(52)
where according to the potential models $`\overline{\mathrm{\Lambda }}0.4`$ GeV is the binding energy of heavy quark (i.e. the constituent mass of light quark), $`T0.45`$ GeV is the average kinetic energy in the system. $`T`$ determines the 2S-1S splitting and it is approximately independent of quark flavors. Analogously, we put
$$\mu _{c\overline{s}}^2=\mu _{b\overline{s}}^2=2T\overline{\mathrm{\Lambda }}_s,$$
(53)
where $`\overline{\mathrm{\Lambda }}_s=\overline{\mathrm{\Lambda }}+(m_{D_s}m_D)=\overline{\mathrm{\Lambda }}+(m_{B_s}m_B)0.5`$ GeV. For the $`bc`$ current we put
$$\mu _{c\overline{b}}^2=2Tm_{bc},$$
(54)
with $`m_{bc}m_bm_c/(m_b+m_c)m_Bm_D/(m_B+m_D)1.3`$ GeV.
Let us suppose that the decay scale is given by the following combinations:
$`\mu _D^2`$ $`=`$ $`\mu _{c\overline{u}}\mu _{c\overline{s}},`$
$`\mu _{D_s}^2`$ $`=`$ $`\mu _{c\overline{s}}^2,`$ (55)
$`\mu _B^2`$ $`=`$ $`\mu _{b\overline{u}}\mu _{c\overline{b}},`$
$`\mu _{B_s}^2`$ $`=`$ $`\mu _{b\overline{s}}\mu _{c\overline{b}},`$
At $`\stackrel{~}{B}=2`$, this set of parameters with $`f_Bf_D175`$ MeV, and $`k=0.18`$ is represented by the SETH column in Table 4. We see a good agreement with the data. The result of $`\mu _B`$ and $`k`$ variations is also shown in Fig. 15.
We stress that the theoretical OPE prediction for the contribution of $`bc\overline{c}s`$ mode, $`k=0.4`$, is significantly overestimated, to our opinion. Indeed, looking at the semileptonic decay with $`\tau ^+\nu _\tau `$ in the final state we see that at the same values of form factors as in modes with light leptons the heavy lepton mode is suppressed because of the restricted phase space, so that the reduction factor equals 0.22. The same effect has to take place in the $`b`$-quark decays with two charmed quarks in the final state. So, since the sum of $`D`$ and $`K`$ masses is greater than the $`\tau `$ mass we could expect even the greater suppression than for the heavy lepton mode. Nevertheless, OPE operates with the quark masses at so moderate release of energy, that certainly leads to the overestimation of phase space in the decay of $`bc\overline{c}s`$. We use more realistic values for $`k`$ close to $`0.2`$.
Next, the contribution of four-quark operators are suppressed in decays of $`\mathrm{\Lambda }_b`$ baryon. So, one expects that the deviation of its total width from that of $`B`$ mesons is about 2-3% . To the moment, the experimental result is far away from this expectation (see Table 4). We can reproduce the data on the $`\mathrm{\Lambda }_b`$ lifetime, if we introduce the light diquark scale
$$\mu _{ud}^2=2\frac{T}{2}\overline{\mathrm{\Lambda }},$$
(56)
where $`T/2`$ corresponds to the half tension of color string inside the diquark. Then we pose
$$\mu _{\mathrm{\Lambda }_b}^2=\mu _{ud}\mu _{c\overline{b}},$$
(57)
which result in a good agreement with the experimental data.
To finish this discussion we draw the conclusion: a probable way to reach the agreement between the theoretical predictions and available experimental data on the lifetimes and inclusive decay widths of heavy hadrons is to suggest the different normalization points in the effective nonleptonic lagrangian for the heavy quark weak decays as dependent on the hadron. This assumption provide us with quite acceptable results to the leading order in $`\alpha _s`$. The variation of normalization point shows the sensitivity of calculations to the higher orders in the QCD coupling constant, which indicates, first, the necessity to proceed with the higher orders, and, second, the appropriate choice of scale can allow us to decrease the scale-dependent higher orders terms.
Thus, we suppose that the preferable choice of scale in the $`cs`$ decays of $`B_c`$ is equal to
$$\mu _{B_c}^2=\mu _{c\overline{b}}\mu _{c\overline{s}}(0.85\mathrm{GeV})^2,$$
(58)
and at $`a_1(\mu _{B_c})=1.20`$ in the charmed quark decays we predict
$$\tau [B_c]=0.48\pm 0.05\mathrm{ps}.$$
(59)
## 7 Conclusion
We have investigated the semileptonic decays of $`B_c`$ meson due to the weak decays of charmed quark in the framework of three-point sum rules in QCD. We have found the important role played by the Coulomb-like $`\alpha _s/v`$-corrections. As in the case of two-point sum rules, the form factors are about three times enhanced due to the Coulomb renormalization of quark-meson vertex for the heavy quarkonium $`B_c`$. We have studied the dependence of form factors on the threshold energy, which determines the continuum region of $`\overline{b}s`$ system. The obtained dependence has the stability region, serving as the test of convergency for the sum rule method. The HQET two-point sum rules for the leptonic constant $`f_{B_s}`$ and $`f_{B_s^{}}`$ have been reanalyzed to introduce the term caused by the product of quark and gluon condensates. This contribution essentially improves the stability of SR results for the leptonic constants of B mesons, yielding: $`f_B=140÷170`$ MeV.
We have studied the soft limit for the form factors in combined HQET/NRQCD technique at the recoil momentum close to zero, which allows us to derive the generalized relations due to the spin symmetry of effective lagrangian. The relations are in a good agreement with the full QCD results, which means that the corrections to the form factors in both relative velocity of heavy quarks inside the $`\overline{b}c`$ quarkonium and the inverse heavy quark masses are small within the accuracy of the method.
Next, we have studied the nonleptonic decays, using the assumption on the factorization of the weak transition. The results on the widths and branching fractions for various decay modes of $`B_c`$ are collected in Tables.
Finally, we have estimated the $`B_c`$ meson lifetime, and showed the dependence on the scale for the hadronic weak lagrangian in decays of charmed quark
$$\tau [B_c]=0.48\pm 0.05\mathrm{ps}.$$
Our estimates are in a good agreement with the theoretical predictions for the lifetime in both the potential models and OPE as well as with the experimental data.
This work was in part supported by the Russian Foundation for Basic Research, grants 99-02-16558 and 00-15-96645.
## 8 Appendix A
For the perturbative spectral densities $`\rho _i(s_1,s_2,Q^2)`$ we have the following expressions :
$`\rho _+(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{2k^{3/2}}}\{{\displaystyle \frac{k}{2}}(\mathrm{\Delta }_1+\mathrm{\Delta }_2)k[m_3(m_3m_1)+m_3(m_3m_2)]`$
$`[2(s_2\mathrm{\Delta }_1+s_1\mathrm{\Delta }_2)u(\mathrm{\Delta }_1+\mathrm{\Delta }_2)]`$
$`[m_3^2{\displaystyle \frac{u}{2}}+m_1m_2m_2m_3m_1m_3]\},`$
$`\rho _{}(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{2k^{3/2}}}\{{\displaystyle \frac{k}{2}}(\mathrm{\Delta }_1\mathrm{\Delta }_2)k[m_3(m_3m_1)m_3(m_3m_2)]+`$
$`[2(s_2\mathrm{\Delta }_1s_1\mathrm{\Delta }_2)+u(\mathrm{\Delta }_1\mathrm{\Delta }_2)]`$
$`[m_3^2{\displaystyle \frac{u}{2}}+m_1m_2m_2m_3m_1m_3]\},`$
$`\rho _V(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{k^{3/2}}}\{(2s_1\mathrm{\Delta }_2u\mathrm{\Delta }_1)(m_3m_2)`$ (62)
$`+(2s_2\mathrm{\Delta }_1u\mathrm{\Delta }_2)(m_3m_1)+m_3k\},`$
$`\rho _0^A(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{k^{1/2}}}\{(m_1m_3)[m_3^2+{\displaystyle \frac{1}{k}}(s_1\mathrm{\Delta }_2^2+s_2\mathrm{\Delta }_1^2u\mathrm{\Delta }_1\mathrm{\Delta }_2)]`$
$`m_2(m_3^2{\displaystyle \frac{\mathrm{\Delta }_1}{2}})m_1(m_3^2{\displaystyle \frac{\mathrm{\Delta }_2}{2}})`$
$`+m_3[m_3^2{\displaystyle \frac{1}{2}}(\mathrm{\Delta }_1+\mathrm{\Delta }_2u)+m_1m_2]\},`$
$`\rho _+^A(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{2k^{3/2}}}\{m_1[2s_2\mathrm{\Delta }_1u\mathrm{\Delta }_2+4\mathrm{\Delta }_1\mathrm{\Delta }_2+2\mathrm{\Delta }_2^2]`$
$`m_1m_3^2[4s_22u]+m_2[2s_1\mathrm{\Delta }_2u\mathrm{\Delta }_1]m_3[2(3s_2\mathrm{\Delta }_1+s_1\mathrm{\Delta }_2)`$
$`u(3\mathrm{\Delta }_2+\mathrm{\Delta }_1)+k+4\mathrm{\Delta }_1\mathrm{\Delta }_2+2\mathrm{\Delta }_2^2+m_3^2(4s_22u)]`$
$`+{\displaystyle \frac{6}{k}}(m_1m_3)[4s_1s_2\mathrm{\Delta }_1\mathrm{\Delta }_2u(2s_2\mathrm{\Delta }_1\mathrm{\Delta }_2+s_1\mathrm{\Delta }_2^2+s_2\mathrm{\Delta }_1^2)`$
$`+2s_2(s_1\mathrm{\Delta }_2^2+s_2\mathrm{\Delta }_1^2)]\},`$
$`\rho _{}^A(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{2k^{5/2}}}\{kum_3(2m_1m_32m_3^2+u)+12(m_1m_3)s_2^2\mathrm{\Delta }_1^2+`$
$`k\mathrm{\Delta }_2[(m_1+m_3)u2s_1(m_2m_3)]+2\mathrm{\Delta }_2^2(k+3us_1)(m_1m_3)`$
$`+\mathrm{\Delta }_1[ku(m_2m_3)+2\mathrm{\Delta }_2(k3u^2)(m_1m_3)]+`$
$`2s_2(m_1m_3)[2km_3^2k\mathrm{\Delta }_1+3u\mathrm{\Delta }_1^26u\mathrm{\Delta }_1\mathrm{\Delta }_2]`$
$`2s_1s_2(km_33\mathrm{\Delta }_2^2(m_1m_3))],`$
$`\rho _+^{}_{}{}^{}A(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{2k^{5/2}}}\{2(m_1m_3)[(k3us_2)\mathrm{\Delta }_1^2+6s_1^2\mathrm{\Delta }_2^2]+`$
$`ku(m_1m_3)(2m_3^2+\mathrm{\Delta }_2)+ku^2m_3+\mathrm{\Delta }_1[ku(2m_1m_23m_3)`$
$`2(m_1m_3)(ks_2k\mathrm{\Delta }_2+3u^2\mathrm{\Delta }_2)]`$
$`2s_1[(m_1m_3)(2km_3^26u\mathrm{\Delta }_1\mathrm{\Delta }_23u\mathrm{\Delta }_2^2)+`$
$`2s_2(km_3+3m_1\mathrm{\Delta }_1^23m_3\mathrm{\Delta }_1^2)+k\mathrm{\Delta }_2(2m_1m_23m_3)]\},`$
$`\rho _{}^{}_{}{}^{}A(s_1,s_2,Q^2)`$ $`=`$ $`{\displaystyle \frac{3}{2k^{5/2}}}\{2(m_1m_3)[(k+3us_2)\mathrm{\Delta }_1^2+6s_1^2\mathrm{\Delta }_2^2]+`$
$`ku(m_1m_3)(2m_3^2+\mathrm{\Delta }_2)+ku^2m_3+\mathrm{\Delta }_1[ku(2m_1m_2+m_3)`$
$`2(m_1m_3)(ks_2k\mathrm{\Delta }_2+3u^2\mathrm{\Delta }_2)]+`$
$`2s_1[(m_1m_3)(2km_3^26u\mathrm{\Delta }_1\mathrm{\Delta }_2+3u\mathrm{\Delta }_2^2)`$
$`2s_2(km_33m_1\mathrm{\Delta }_1^2+3m_3\mathrm{\Delta }_1^2)+k\mathrm{\Delta }_2(2m_1+m_2m_3)]\}.`$
Here $`k=(s_1+s_2+Q^2)^24s_1s_2`$, $`u=s_1+s_2+Q^2`$, $`\mathrm{\Delta }_1=s_1m_1^2+m_3^2`$ and $`\mathrm{\Delta }_2=s_2m_2^2+m_3^2`$. $`m_1,m_2`$ and $`m_3`$ are the masses of quark flavours relevant to the various decays, see prescriptions in Fig. 1.
## 9 Appendix B
Here we list the expression for the form factors of semileptonic decays $`B_cB_s^{()}`$ taken from the potential model .
$`f_+`$ $`=`$ $`{\displaystyle \frac{(\stackrel{~}{m}_c+\stackrel{~}{m}_s)}{2\stackrel{~}{m}_s}}\sqrt{{\displaystyle \frac{M_{B_s}}{M_{B_c}}}}\xi (w),`$ (68)
$`f_{}`$ $`=`$ $`{\displaystyle \frac{(\stackrel{~}{m}_c\stackrel{~}{m}_s+2\stackrel{~}{m}_b)}{2\stackrel{~}{m}_s}}\sqrt{{\displaystyle \frac{M_{B_s}}{M_{B_c}}}}\xi (w),`$ (69)
$`F_0^A`$ $`=`$ $`{\displaystyle \frac{M_{B_c}^2+M_{B_s^{}}^2q^2+2M_{B_c}(\stackrel{~}{m}_s\stackrel{~}{m}_b)}{2\stackrel{~}{m}_s}}\sqrt{{\displaystyle \frac{M_{B_s^{}}}{M_{B_c}}}}\xi (w),`$ (70)
$`F_+^A`$ $`=`$ $`{\displaystyle \frac{12\stackrel{~}{m}_b/M_{B_c}}{2\stackrel{~}{m}_s}}\sqrt{{\displaystyle \frac{M_{B_s^{}}}{M_{B_c}}}}\xi (w),`$ (71)
$`F_{}^A`$ $`=`$ $`{\displaystyle \frac{1+2\stackrel{~}{m}_b/M_{B_c}}{2\stackrel{~}{m}_s}}\sqrt{{\displaystyle \frac{M_{B_s^{}}}{M_{B_c}}}}\xi (w),`$ (72)
where
$$\xi (w)=\frac{2\omega \omega _x}{\omega ^2+\omega _x^2}\sqrt{\frac{2\omega \omega _x}{\omega ^2w^2+\omega _x^2}}\mathrm{exp}\left(\frac{\stackrel{~}{m}_b^2(w^21)}{\omega ^2w^2+\omega _x^2}\right),$$
(73)
$$\omega =2\pi \left(\frac{M_{B_c}\stackrel{~}{f}_{B_c}^2}{12}\right)^{1/3},\omega _x=2\pi \left(\frac{M_{B_s^{()}}\stackrel{~}{f}_{B_s^{()}}^2}{12}\right)^{1/3},$$
(74)
where $`w`$ is the product of $`B_c`$ and $`B_s^{()}`$ four-velocities. The quark masses and the leptonic constants have the values generally used in the calculations in the framework of potential models
$$\stackrel{~}{m}_b=4.8\text{GeV},\stackrel{~}{m}_c=1.5\text{GeV},\stackrel{~}{m}_s=0.55\text{GeV},\stackrel{~}{f}_{B_c}=0.47\text{GeV},\stackrel{~}{f}_{B_s^{()}}=0.17\text{GeV}.$$ |
warning/0002/hep-ph0002136.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Although, on the one hand, all experimental observations up to now are perfectly consistent with $`CPT`$ symmetry, and, on the other hand, the standard field theory implies that this symmetry should hold exactly, continued experimental, phenomenological and theoretical studies of this and related symmetries are warranted. In this connection, we like to recall, on the one hand, that $`CP`$ symmetry is violated only at such a tiny level as $`10^3`$ , while $`CPT`$ symmetry is tested at best up to the level one order smaller \[3-7\] and, on the other hand, that some of the premises of the $`CPT`$ theorem, e.g., locality, are being challenged by, say, the superstring model.
In a series of papers \[4-7\], we have demonstrated how one may identify or constrain possible violation of $`CP`$, $`T`$ and $`CPT`$ symmetries in the $`K^0`$-$`\overline{K^0}`$ system in a way as phenomenological and comprehensive as possible. For this purpose, we have first introduced parameters which represent violation of these symmetries in mixing parameters and decay amplitudes in a well-defined way and related them to the experimentally measured quantities. We have then carried out numerical analyses, with the aid of the Bell-Steinberger relation and with all the available data on $`2\pi `$, $`3\pi `$, $`\pi ^+\pi ^{}\gamma `$ and $`\pi \mathrm{}\nu _{\mathrm{}}`$ decays used as inputs, to derive constraints to these symmetry-violating parameters. It has been shown among other things that the new results on the asymptotic leptonic asymmetries obtained by CPLEAR Collaboration allow one for the first time to constrain to some extent possible $`CPT`$ violation in the $`\pi \mathrm{}\nu _{\mathrm{}}`$ decay modes.<sup>a</sup><sup>a</sup>aWe afterwards became aware that CPLEAR Collaboration themselves had also, by an analysis more or less similar to ours, reached the similar conclusion independently.
The present work is a continuation of the previous works, which is new particularly in the following points:
(1) The new results on $`\text{Re}(\epsilon ^{}/\epsilon )`$, etc., from the Fermilab KTeV and CERN NA48 experiments , along with CPLEAR’s new data \[13-16\] and the latest version of the data compiled by Particle Data Group(PDG) , are used as inputs.
(2) A particular attention is paid to clarify what can be said without recourse to the Bell-Steinberger relation and what can be said with the aid of this relation.
(3) A case study with either $`CPT`$ or $`T`$ symmetry assumed is also carried out.
(4) The relevant decay amplitudes are parametrized in a convenient form, with freedom associated with rephasing of both the initial and final states, as discussed explicitly and thoroughly in , taken into account.
The paper is organized as follows. The theoretical framework used to describe the $`K^0`$ \- $`\overline{K^0}`$ system , including the Bell-Steinberger relation, is recapitulated in Sec.2 and the experimentally measured quantities related to $`CP`$ violation in decay modes of interest to us are enumerated in Sec.3. We then parametrize the mixing parameters and decay amplitudes in a convenient and well-defined way and give conditions imposed by $`CP`$, $`T`$ and/or $`CPT`$ symmetries on these parameters in Sec.4. In Sec.5, experimentally measured quantities are expressed in terms of the parameters defined, treating them as first order small. In Sec.6, paying particular attention to the data provided by KTeV Collaboration and by NA48 Collaboration, a numerical analysis is performed, while, in Sec.7, with most of the available experimental data, including those reported by CPLEAR Collaboration, used as inputs, a more comprehensive numerical analysis is performed. Sec.8 is devoted to a case study, in which the case with $`CPT`$ symmetry assumed and the case with $`T`$ symmetry assumed are considered separately. The results of the analyses are summarized and some concluding remarks are given in Sec.9.
## 2 The $`K^0`$-$`\overline{K^0}`$ mixing and the Bell-Steinberger <br>relation
Let $`|K^0`$ and $`|\overline{K^0}`$ be eigenstates of the strong interaction with strangeness $`S=+1`$ and $`1`$, related to each other by $`(CP)`$, $`(CPT)`$ and $`T`$ operations as
$$\begin{array}{cc}(CP)|K^0=e^{i\alpha _K}|\overline{K^0},& (CPT)|K^0=e^{i\beta _K}|\overline{K^0},\\ (CP)|\overline{K^0}=e^{i\alpha _K}|K^0,& (CPT)|\overline{K^0}=e^{i\beta _K}|K^0,\\ T|K^0=e^{i(\beta _K\alpha _K)}|K^0,& T|\overline{K^0}=e^{i(\beta _K+\alpha _K)}|\overline{K^0}.\end{array}$$
(2.1)
Note here that, given the first two where $`\alpha _K`$ and $`\beta _K`$ are arbitrary real parameters, the rest follow from the assumptions $`(CP)T=T(CP)=(CPT)`$, $`(CP)^2=(CPT)^2=1`$, and anti-linearity of $`T`$ and $`(CPT)`$. When the weak interaction $`H_\text{w}`$ is switched on, the $`K^0`$ and $`\overline{K^0}`$ states decay into other states, generically denoted as $`|n`$, and get mixed. The time evolution of the arbitrary state
$$|\mathrm{\Psi }(t)=c_1(t)|K_1+c_2(t)|K_2,$$
with
$$|K_1|K^0,|K_2|\overline{K^0},$$
is described by a Schrödinger-like equation
$$i\frac{d}{dt}|\mathrm{\Psi }=\mathrm{\Lambda }|\mathrm{\Psi },$$
or
$$i\frac{d}{dt}\left(\begin{array}{c}c_1(t)\\ c_2(t)\end{array}\right)=\mathrm{\Lambda }\left(\begin{array}{c}c_1(t)\\ c_2(t)\end{array}\right).$$
(2.2)
The operator or $`2\times 2`$ matrix $`\mathrm{\Lambda }`$ may be written as
$$\mathrm{\Lambda }Mi\mathrm{\Gamma }/2,$$
(2.3)
with $`M`$ (mass matrix) and $`\mathrm{\Gamma }`$ (decay or width matrix) given, to the second order in $`H_\text{w}`$, by
$`M_{ij}`$ $``$ $`K_i|M|K_j`$ (2.4a)
$`=`$ $`m_K\delta _{ij}+K_i|H_\text{w}|K_j`$
$`+{\displaystyle \underset{n}{}}P{\displaystyle \frac{K_i|H_\text{w}|nn|H_\text{w}|K_j}{m_KE_n}},`$
$`\mathrm{\Gamma }_{ij}`$ $``$ $`K_i|\mathrm{\Gamma }|K_j`$
$`=`$ $`2\pi {\displaystyle \underset{n}{}}K_i|H_\text{w}|nn|H_\text{w}|K_j\delta (m_KE_n),`$
where the operator $`P`$ projects out the principal value. The two eigenstates of $`\mathrm{\Lambda }`$ and their respective eigenvalue may be written as
$`|K_S`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{|p_S|^2+|q_S|^2}}}\left(p_S|K^0+q_S|\overline{K^0}\right),`$ (2.5a)
$`|K_L`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{|p_L|^2+|q_L|^2}}}\left(p_L|K^0q_L|\overline{K^0}\right);`$ (2.5b)
$`\lambda _S=m_Si{\displaystyle \frac{\gamma _S}{2}},`$ (2.6a)
$`\lambda _L=m_Li{\displaystyle \frac{\gamma _L}{2}}.`$ (2.6b)
$`m_{S,L}=\text{Re}(\lambda _{S,L})`$ and $`\gamma _{S,L}=2\text{Im}(\lambda _{S,L})`$ are the mass and the total decay width of the $`K_{S,L}`$ state respectively. By definition, $`\gamma _S>\gamma _L`$ or $`\tau _S<\tau _L`$ ($`\tau _{S,L}1/\gamma _{S,L}`$), and the suffices $`S`$ and $`L`$ stand for ”short-lived” and ”long-lived” respectively. The eigenvalues $`\lambda _{S,L}`$ and the ratios of the mixing parameters $`q_{S,L}/p_{S,L}`$ are related to the elements of the mass-width matrix $`\mathrm{\Lambda }`$ as
$$\lambda _{S,L}=\pm E+(\mathrm{\Lambda }_{11}+\mathrm{\Lambda }_{22})/2,$$
(2.7)
$$q_{S,L}/p_{S,L}=\mathrm{\Lambda }_{21}/[E\pm (\mathrm{\Lambda }_{11}\mathrm{\Lambda }_{22})/2],$$
(2.8)
where
$$E[\mathrm{\Lambda }_{12}\mathrm{\Lambda }_{21}+(\mathrm{\Lambda }_{11}\mathrm{\Lambda }_{22})^2/4]^{1/2}.$$
(2.9)
¿From the eigenvalue equation of $`\mathrm{\Lambda }`$, one may readily derive the well-known Bell-Steinberger relation :
$$\left[\frac{\gamma _S+\gamma _L}{2}i(m_Sm_L)\right]K_S|K_L=K_S|\mathrm{\Gamma }|K_L,$$
(2.10)
where
$$K_S|\mathrm{\Gamma }|K_L=2\pi \underset{n}{}K_S|H_\text{w}|nn|H_\text{w}|K_L\delta (m_KE_n).$$
(2.11)
## 3 Decay modes
The $`K^0`$ and $`\overline{K^0}`$ (or $`K_S`$ and $`K_L`$) states have many decay modes, among which we are interested in $`2\pi `$, $`3\pi `$, $`\pi ^+\pi ^{}\gamma `$ and semileptonic modes.
### 3.1 $`2\pi `$ modes
The experimentally measured quantities related to $`CP`$ violation are $`\eta _+`$ and $`\eta _{00}`$ defined by
$$\eta _+|\eta _+|e^{i\varphi _+}\frac{\pi ^+\pi ^{},\text{outgoing}|H_\text{w}|K_L}{\pi ^+\pi ^{},\text{outgoing}|H_\text{w}|K_S},$$
(3.1a)
$$\eta _{00}|\eta _{00}|e^{i\varphi _{00}}\frac{\pi ^0\pi ^0,\text{outgoing}|H_\text{w}|K_L}{\pi ^0\pi ^0,\text{outgoing}|H_\text{w}|K_S}.$$
(3.1b)
Defining
$$\omega \frac{(2\pi )_2|H_\text{w}|K_S}{(2\pi )_0|H_\text{w}|K_S},$$
(3.2)
$$\eta _I|\eta _I|e^{i\varphi _I}\frac{(2\pi )_I|H_\text{w}|K_L}{(2\pi )_I|H_\text{w}|K_S},$$
(3.3)
where $`I`$=1 or 2 stands for the isospin of the $`2\pi `$ states, one gets
$`\eta _+`$ $`=`$ $`{\displaystyle \frac{\eta _0+\eta _2\omega ^{}}{1+\omega ^{}}},`$ (3.4a)
$`\eta _{00}`$ $`=`$ $`{\displaystyle \frac{\eta _02\eta _2\omega ^{}}{12\omega ^{}}},`$ (3.4b)
where
$$\omega ^{}\frac{1}{\sqrt{2}}\omega e^{i(\delta _2\delta _0)},$$
(3.5)
$`\delta _I`$ being the S-wave $`\pi \pi `$ scattering phase shift for the isospin $`I`$ state at an energy of the rest mass of $`K^0`$. $`\omega `$ is a measure of deviation from the $`\mathrm{\Delta }I=1/2`$ rule, and may be inferred, for example, from
$`r`$ $``$ $`{\displaystyle \frac{\gamma _S(\pi ^+\pi ^{})2\gamma _S(\pi ^0\pi ^0)}{\gamma _S(\pi ^+\pi ^{})+\gamma _S(\pi ^0\pi ^0)}}`$ (3.6)
$`=`$ $`{\displaystyle \frac{4\text{Re}(\omega ^{})2|\omega ^{}|^2}{1+2|\omega ^{}|^2}}.`$
Here and in the following, $`\gamma _{S,L}(n)`$ denotes the partial width for $`K_{S,L}`$ to decay into the final state $`|n`$.
### 3.2 $`3\pi `$ and $`\pi ^+\pi ^{}\gamma `$ modes
The experimentally measured quantities are
$$\eta _{+0}=\frac{\pi ^+\pi ^{}\pi ^0,\text{outgoing}|H_\text{w}|K_S}{\pi ^+\pi ^{}\pi ^0,\text{outgoing}|H_\text{w}|K_L},$$
(3.7a)
$$\eta _{000}=\frac{\pi ^0\pi ^0\pi ^0,\text{outgoing}|H_\text{w}|K_S}{\pi ^0\pi ^0\pi ^0,\text{outgoing}|H_\text{w}|K_L},$$
(3.7b)
$$\eta _{+\gamma }=\frac{\pi ^+\pi ^{}\gamma ,\text{outgoing}|H_\text{w}|K_L}{\pi ^+\pi ^{}\gamma ,\text{outgoing}|H_\text{w}|K_S}.$$
(3.8)
We shall treat the $`3\pi `$ $`(\pi ^+\pi ^{}\gamma )`$ states as purely $`CP`$-odd ($`CP`$-even).
### 3.3 Semileptonic modes
The well measured time-independent asymmetry parameter related to $`CP`$ violation in semi-leptonic decay modes is
$$d_L^{\mathrm{}}=\frac{\gamma _L(\pi ^{}\mathrm{}^+\nu _{\mathrm{}})\gamma _L(\pi ^+\mathrm{}^{}\overline{\nu }_{\mathrm{}})}{\gamma _L(\pi ^{}\mathrm{}^+\nu _{\mathrm{}})+\gamma _L(\pi ^+\mathrm{}^{}\overline{\nu }_{\mathrm{}})},$$
(3.9)
where $`\mathrm{}=e`$ or $`\mu `$. CPLEAR Collaboration \[9,14-16\] have furthermore for the first time measured two kinds of time-dependent asymmetry parameters
$$d_1^{\mathrm{}}(t)=\frac{|\mathrm{}^+|H_\text{w}|\overline{K^0}(t)|^2|\mathrm{}^{}|H_\text{w}|K^0(t)|^2}{|\mathrm{}^+|H_\text{w}|\overline{K^0}(t)|^2+|\mathrm{}^{}|H_\text{w}|K^0(t)|^2},$$
(3.10a)
$$d_2^{\mathrm{}}(t)=\frac{|\mathrm{}^{}|H_\text{w}|\overline{K^0}(t)|^2|\mathrm{}^+|H_\text{w}|K^0(t)|^2}{|\mathrm{}^{}|H_\text{w}|\overline{K^0}(t)|^2+|\mathrm{}^+|H_\text{w}|K^0(t)|^2},$$
(3.10b)
where $`|\mathrm{}^+=|\pi ^{}\mathrm{}^+\nu _{\mathrm{}}`$ and $`|\mathrm{}^{}=|\pi ^+\mathrm{}^{}\overline{\nu }_{\mathrm{}}`$.
## 4 Parametrization and conditions imposed by $`CP`$, $`T`$ and $`CPT`$ symmetries
We shall parametrize the ratios of the mixing parameters $`q_S/p_S`$ and $`q_L/p_L`$ as
$$\begin{array}{c}\frac{q_S}{p_S}=e^{i\alpha _K}\frac{1\epsilon _S}{1+\epsilon _S},\\ \\ \frac{q_L}{p_L}=e^{i\alpha _K}\frac{1\epsilon _L}{1+\epsilon _L},\end{array}$$
(4.1)
and $`\epsilon _{S,L}`$ further as
$$\epsilon _{S,L}=\epsilon \pm \delta .$$
(4.2)
¿From Eqs.(2.7), (2.8) and (2.9), treating $`\epsilon `$ and $`\delta `$ as small parameters, one may derive
$`\mathrm{\Delta }m`$ $``$ $`2\text{Re}(M_{12}e^{i\alpha _K}),`$ (4.3a)
$`\mathrm{\Delta }\gamma `$ $``$ $`2\text{Re}(\mathrm{\Gamma }_{12}e^{i\alpha _K}),`$ (4.3b)
$`\epsilon `$ $``$ $`(\mathrm{\Lambda }_{12}e^{i\alpha _K}\mathrm{\Lambda }_{21}e^{i\alpha _K})/2\mathrm{\Delta }\lambda ,`$ (4.4a)
$`\delta `$ $``$ $`(\mathrm{\Lambda }_{11}\mathrm{\Lambda }_{22})/2\mathrm{\Delta }\lambda ,`$ (4.4b)
from which it follows that
$$\epsilon _{}\text{Re}[\epsilon \mathrm{exp}(i\varphi _{SW})]\frac{2\text{Im}(M_{12}e^{i\alpha _K})}{\sqrt{(\gamma _S\gamma _L)^2+4(\mathrm{\Delta }m)^2}},$$
(4.5a)
$$\epsilon _{}\text{Im}[\epsilon \mathrm{exp}(i\varphi _{SW})]\frac{\text{Im}(\mathrm{\Gamma }_{12}e^{i\alpha _K})}{\sqrt{(\gamma _S\gamma _L)^2+4(\mathrm{\Delta }m)^2}},$$
(4.5b)
$$\delta _{}\text{Re}[\delta \mathrm{exp}(i\varphi _{SW})]\frac{(\mathrm{\Gamma }_{11}\mathrm{\Gamma }_{22})}{2\sqrt{(\gamma _S\gamma _L)^2+4(\mathrm{\Delta }m)^2}},$$
(4.6a)
$$\delta _{}\text{Im}[\delta \mathrm{exp}(i\varphi _{SW})]\frac{(M_{11}M_{22})}{\sqrt{(\gamma _S\gamma _L)^2+4(\mathrm{\Delta }m)^2}},$$
(4.6b)
where
$`\mathrm{\Delta }mm_Sm_L,\mathrm{\Delta }\gamma `$ $``$ $`\gamma _S\gamma _L,\mathrm{\Delta }\lambda \lambda _S\lambda _L,`$ (4.7a)
$`\varphi _{SW}`$ $``$ $`\mathrm{tan}^1\left({\displaystyle \frac{2\mathrm{\Delta }m}{\mathrm{\Delta }\gamma }}\right).`$ (4.7b)
$`\varphi _{SW}`$ is often called the superweak phase.
Paying particular attention to the $`2\pi `$ and semileptonic decay modes, we shall parametrize amplitudes for $`K^0`$ and $`\overline{K^0}`$ to decay into $`|(2\pi )_I`$ as
$`(2\pi )_I|H_\text{w}|K^0`$ $`=`$ $`F_I(1+\epsilon _I)e^{i\alpha _K/2},`$ (4.8a)
$`(2\pi )_I|H_\text{w}|\overline{K^0}`$ $`=`$ $`F_I(1\epsilon _I)e^{i\alpha _K/2},`$ (4.8b)
and amplitudes for $`K^0`$ and $`\overline{K^0}`$ to decay into $`|\mathrm{}^+`$ and $`|\mathrm{}^{}`$ as
$`\mathrm{}^+|H_\text{w}|K^0`$ $`=`$ $`F_{\mathrm{}}(1+\epsilon _{\mathrm{}})e^{i\alpha _K/2},`$ (4.9a)
$`\mathrm{}^{}|H_\text{w}|\overline{K^0}`$ $`=`$ $`F_{\mathrm{}}(1\epsilon _{\mathrm{}})e^{i\alpha _K/2},`$ (4.9b)
$`\mathrm{}^+|H_\text{w}|\overline{K^0}`$ $`=`$ $`x_\mathrm{}+F_{\mathrm{}}(1+\epsilon _{\mathrm{}})e^{i\alpha _K/2},`$ (4.9c)
$`\mathrm{}^{}|H_\text{w}|K^0`$ $`=`$ $`x_{\mathrm{}}^{}F_{\mathrm{}}(1\epsilon _{\mathrm{}})e^{i\alpha _K/2}.`$ (4.9d)
$`x_\mathrm{}+`$ and $`x_{\mathrm{}}`$, which measure violation of the $`\mathrm{\Delta }S=\mathrm{\Delta }Q`$ rule, will further be parametrized as
$$x_\mathrm{}+=x_{\mathrm{}}^{(+)}+x_{\mathrm{}}^{()},x_{\mathrm{}}=x_{\mathrm{}}^{(+)}x_{\mathrm{}}^{()}.$$
(4.10)
Our amplitude parameters $`F_I`$, $`\epsilon _I`$, $`F_{\mathrm{}}`$, $`\epsilon _{\mathrm{}}`$, $`x_{\mathrm{}}^{(+)}`$ and $`x_{\mathrm{}}^{()}`$, and our mixing parameters $`\epsilon `$ and $`\delta `$ as well, are all invariant with respect to rephasing of $`|K^0`$ and $`|\overline{K^0}`$,
$$|K^0|K^0^{}=|K^0e^{i\xi _K},|\overline{K^0}|\overline{K^0}^{}=|\overline{K^0}e^{i\xi _K},$$
(4.11)
in spite that $`\alpha _K`$ itself is not invariant with respect to this rephasing . $`F_I`$, $`\epsilon _I`$, $`F_{\mathrm{}}`$ and $`\epsilon _{\mathrm{}}`$ are however not invariant with respect to rephasing of the final states ,
$`|(2\pi )_I|(2\pi )_I^{}=|(2\pi )_Ie^{i\xi _I},`$ (4.12a)
$`|\mathrm{}^+|\mathrm{}^+^{}=|\mathrm{}^+e^{i\xi _\mathrm{}+},|\mathrm{}^{}|\mathrm{}^{}^{}=|\mathrm{}^{}e^{i\xi _{\mathrm{}}},`$ (4.12b)
nor are the relative $`CP`$ and $`CPT`$ phases $`\alpha _{\mathrm{}}`$, $`\beta _I`$ and $`\beta _{\mathrm{}}`$ defined in such a way as
$`CP|\mathrm{}^+=e^{i\alpha _{\mathrm{}}}|\mathrm{}^{},`$ (4.13a)
$`CPT|(2\pi )_I=e^{i\beta _I}|(2\pi )_I,CPT|\mathrm{}^+=e^{i\beta _{\mathrm{}}}|\mathrm{}^{}.`$ (4.13b)
One may convince himself that freedom associated with choice of $`\xi _I`$, $`\xi _\mathrm{}++\xi _{\mathrm{}}`$ and $`\xi _\mathrm{}+\xi _{\mathrm{}}`$ allows one, without loss of generality, to take<sup>b</sup><sup>b</sup>bNote that, although freedom associated with $`\xi _K`$ and $`\xi _\mathrm{}+\xi _{\mathrm{}}`$ allows one to take $`\alpha _K=0`$ and $`\alpha _{\mathrm{}}=0`$ (instead of $`\text{Im}(\epsilon _{\mathrm{}})=0`$) respectively, we prefer not to do so. Note also that our parametrization (4.9a,b) is similar to, but different from the one more widely adopted ,
$$\mathrm{}^+|H_\text{w}|K^0=F_{\mathrm{}}(1y_{\mathrm{}}),\mathrm{}^{}|H_\text{w}|\overline{K^0}=F_{\mathrm{}}^{}(1+y_{\mathrm{}}^{}),$$
and that, nevertheless, our $`\text{Re}(\epsilon _{\mathrm{}})`$ is exactly equivalent to $`\text{Re}(y_{\mathrm{}})`$ introduced through these equations and also to $`\text{Re}(y)`$ defined in \[14-16\].
$$\text{Im}(F_I)=0,\text{Im}(F_{\mathrm{}})=0,\text{Im}(\epsilon _{\mathrm{}})=0,$$
(4.14)
respectively, and that $`CP`$, $`T`$ and $`CPT`$ symmetries impose such conditions as
$$\begin{array}{ccc}CP\text{symmetry}& :& \epsilon =0,\delta =0,\epsilon _I=0,\text{Re}(\epsilon _{\mathrm{}})=0,\hfill \\ & & \text{Im}(x_{\mathrm{}}^{(+)})=0,\text{Re}(x_{\mathrm{}}^{()})=0;\hfill \\ T\text{symmetry}& :& \epsilon =0,\text{Im}(\epsilon _I)=0,\text{Im}(x_{\mathrm{}}^{(+)})=0,\hfill \\ & & \text{Im}(x_{\mathrm{}}^{()})=0;\hfill \\ CPT\text{symmetry}& :& \delta =0,\text{Re}(\epsilon _I)=0,\text{Re}(\epsilon _{\mathrm{}})=0,\hfill \\ & & \text{Re}(x_{\mathrm{}}^{()})=0,\text{Im}(x_{\mathrm{}}^{()})=0.\hfill \end{array}$$
(4.15)
Among these parameters, $`\epsilon `$ and $`\delta `$ will be referred to as indirect parameters and the rest as direct parameters.<sup>c</sup><sup>c</sup>cAs emphasized in , classification of the symmetry-violating parameters into ”direct” and ”indirect” ones makes sense only when they are defined in such a way that they are invariant under rephasing of $`|K^0`$ and $`|\overline{K^0}`$, Eq.(4.11).
## 5 Formulas relevant for numerical analysis
We shall adopt a phase convention which gives Eq.(4.14). Observed or expected smallness of violation of $`CP`$, $`T`$ and $`CPT`$ symmetries and of the $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }Q=\mathrm{\Delta }S`$ rules allows us to treat all our parameters, $`\epsilon `$, $`\delta `$, $`\epsilon _I`$ , $`\epsilon _{\mathrm{}}`$, $`x_{\mathrm{}}^{(+)}`$, $`x_{\mathrm{}}^{()}`$ as well as $`\omega ^{}`$ as small, <sup>d</sup><sup>d</sup>dAs a matter of fact, we have already assumed that $`CP`$, $`T`$ and $`CPT`$ violations are small in deriving Eqs.(4.3a) $``$ (4.6b) and in parametrizing the relevant amplitudes as in Eqs.(4.8) and (4.9). and, from Eqs.(3.2), (3.3), (3.4a,b), (3.6), (3.9) and (3.10a,b), one finds, to the leading order in these small parameters,
$$\omega \text{Re}(F_2)/\text{Re}(F_0),$$
(5.1)
$$\eta _I\epsilon \delta +\epsilon _I,$$
(5.2)
$`\eta _+`$ $``$ $`\eta _0+\epsilon ^{},`$ (5.3a)
$`\eta _{00}`$ $``$ $`\eta _02\epsilon ^{},`$ (5.3b)
$$r4\text{Re}(\omega ^{}),$$
(5.4)
$$d_L^{\mathrm{}}2\text{Re}(\epsilon \delta )+2\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()}),$$
(5.5)
$`d_1^{\mathrm{}}(t\tau _S)4\text{Re}(\epsilon )+2\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()}),`$ (5.6a)
$`d_2^{\mathrm{}}(t\tau _S)4\text{Re}(\delta )2\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()}),`$ (5.6b)
where
$$\epsilon ^{}(\eta _2\eta _0)\omega ^{}.$$
(5.7)
Note that $`d_L^{\mathrm{}}`$, $`d_1^{\mathrm{}}(t\tau _S)`$ and $`d_2^{\mathrm{}}(t\tau _S)`$ are not independent:
$$d_L^{\mathrm{}}[d_1^{\mathrm{}}(t\tau _S)d_2^{\mathrm{}}(t\tau _S)]/2.$$
(5.8)
¿From Eqs.(5.3a,b), it follows that
$$\eta _0(2/3)\eta _+(1+(1/2)|\eta _{00}/\eta _+|e^{i\mathrm{\Delta }\varphi }),$$
(5.9)
and, treating $`|\epsilon ^{}/\eta _0|`$ as a small quantity, which is justifiable empirically (see below), one further obtains
$$\eta _{00}/\eta _+13\epsilon ^{}/\eta _0,$$
(5.10)
or
$`\text{Re}(\epsilon ^{}/\eta _0)`$ $``$ $`(1/6)(1|\eta _{00}/\eta _+|^2),`$ (5.11a)
$`\text{Im}(\epsilon ^{}/\eta _0)`$ $``$ $`(1/3)\mathrm{\Delta }\varphi ,`$ (5.11b)
where
$$\mathrm{\Delta }\varphi \varphi _{00}\varphi _+.$$
(5.12)
On the other hand, from Eqs.(3.5), (5.1), (5.2) and (5.7), one may derive
$$\epsilon ^{}/\eta _0=i\text{Re}(\omega ^{})(\epsilon _2\epsilon _0)e^{i\mathrm{\Delta }\varphi ^{}}/[|\eta _0|\mathrm{cos}(\delta _2\delta _0)],$$
(5.13)
where
$$\mathrm{\Delta }\varphi ^{}\varphi _0\delta _2+\delta _0\pi /2.$$
(5.14)
Furthermore, noting that
$$K_S|K_L2[\text{Re}(\epsilon )i\text{Im}(\delta )],$$
one may use the Bell-Steinberger relation, Eq.(2.10), to express $`\text{Re}(\epsilon )`$ and $`\text{Im}(\delta )`$ in terms of measured quantities. By taking $`2\pi `$, $`3\pi `$, $`\pi ^+\pi ^{}\gamma `$ and $`\pi \mathrm{}\nu _{\mathrm{}}`$ intermediate states into account in Eq.(2.11) and making use of the fact $`\gamma _S\gamma _L`$, we derive
$`\text{Re}(\epsilon )`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\gamma _S^2+4(\mathrm{\Delta }m)^2}}}\times `$ (5.15)
$`[\gamma _S(\pi ^+\pi ^{})|\eta _+|\mathrm{cos}(\varphi _+\varphi _{SW})`$
$`+\gamma _S(\pi ^0\pi ^0)|\eta _{00}|\mathrm{cos}(\varphi _{00}\varphi _{SW})`$
$`+\gamma _S(\pi ^+\pi ^{}\gamma )|\eta _{+\gamma }|\mathrm{cos}(\varphi _{+\gamma }\varphi _{SW})`$
$`+\gamma _L(\pi ^+\pi ^{}\pi ^0)\{\text{Re}(\eta _{+0})\mathrm{cos}\varphi _{SW}\text{Im}(\eta _{+0})\mathrm{sin}\varphi _{SW}\}`$
$`+\gamma _L(\pi ^0\pi ^0\pi ^0)\{\text{Re}(\eta _{000})\mathrm{cos}\varphi _{SW}\text{Im}(\eta _{000})\mathrm{sin}\varphi _{SW}\}`$
$`+2{\displaystyle \underset{\mathrm{}}{}}\gamma _L(\pi \mathrm{}\nu _{\mathrm{}})\{\text{Re}(\epsilon _{\mathrm{}})\mathrm{cos}\varphi _{SW}\text{Im}(x_{\mathrm{}}^{(+)})\mathrm{sin}\varphi _{SW}\}],`$
$`\text{Im}(\delta )`$ $``$ $`{\displaystyle \frac{1}{\sqrt{\gamma _S^2+4(\mathrm{\Delta }m)^2}}}\times `$ (5.16)
$`[\gamma _S(\pi ^+\pi ^{})|\eta _+|\mathrm{sin}(\varphi _+\varphi _{SW})`$
$`\gamma _S(\pi ^0\pi ^0)|\eta _{00}|\mathrm{sin}(\varphi _{00}\varphi _{SW})`$
$`\gamma _S(\pi ^+\pi ^{}\gamma )|\eta _{+\gamma }|\mathrm{sin}(\varphi _{+\gamma }\varphi _{SW})`$
$`+\gamma _L(\pi ^+\pi ^{}\pi ^0)\{\text{Re}(\eta _{+0})\mathrm{sin}\varphi _{SW}+\text{Im}(\eta _{+0})\mathrm{cos}\varphi _{SW}\}`$
$`+\gamma _L(\pi ^0\pi ^0\pi ^0)\{\text{Re}(\eta _{000})\mathrm{sin}\varphi _{SW}+\text{Im}(\eta _{000})\mathrm{cos}\varphi _{SW}\}`$
$`+2{\displaystyle \underset{\mathrm{}}{}}\gamma _L(\pi \mathrm{}\nu _{\mathrm{}})\{\text{Re}(\epsilon _{\mathrm{}})\mathrm{sin}\varphi _{SW}+\text{Im}(x_{\mathrm{}}^{(+)})\mathrm{cos}\varphi _{SW}\}].`$
If, however, one retains the contribution of the $`2\pi `$ intermediate states alone, which is justfiable empirically, the Bell-Steinberger relation gives simply
$$\text{Re}(\epsilon )i\text{Im}(\delta )|\eta _0|e^{i\mathrm{\Delta }\varphi \mathrm{"}}\mathrm{cos}\varphi _{SW},$$
(5.17)
where
$$\mathrm{\Delta }\varphi \mathrm{"}\varphi _0\varphi _{SW}.$$
(5.18)
It is to be noted that exactly the same equations as Eq.(5.17) can be derived from Eqs.(4.5b) and (4.6a).
## 6 Numerical analysis (1) -Constraints from the <br>KTeV and NA48 data-
The data used as inputs in the numerical analysis given below are tabulated in Table 1. As the value of $`|\eta _{00}/\eta _+|`$ or $`\text{Re}(\epsilon ^{}/\eta _0)`$, we adopt those reported by KTeV Collaboration and NA48 Collaboration,<sup>e</sup><sup>e</sup>eNote that $`\eta _0`$ corresponds to $`\epsilon `$ used in . Since only $`\text{Re}(\epsilon ^{}/\epsilon )`$, but not $`|\eta _{00}/\eta _+|^2`$, is reported explicitly in , we take a weighted average of the two values of $`\text{Re}(\epsilon ^{}/\epsilon )`$ reported in and list this in Table 2 below. and as the values of $`\mathrm{\Delta }m`$, $`\tau _S`$ and $`\mathrm{\Delta }\varphi `$, we use those reported by KTeV Collaboration. As for $`\delta _2\delta _0`$, we use $`(42\pm 20)^{}`$, i.e., the Chell-Olsson value with the error arbitrarily extended by a factor of five to take account of its possible uncertainty . All the other data are from Particle Data Group (PDG) .
Our analysis consists of two parts:
The first half. We use Eq.(4.7b) to find $`\varphi _{SW}`$ from $`\mathrm{\Delta }m`$ and $`\gamma _S`$, use Eqs.(3.6) and (5.4) to find $`\text{Re}(\omega ^{})`$ from $`\gamma _S(\pi ^+\pi ^{})/\gamma _S`$ and $`\gamma _S(\pi ^0\pi ^0)/\gamma _S`$, and further use Eqs.(5.11a,b) and (5.9) to find $`\text{Re}(\epsilon ^{}/\eta _0)`$, $`\text{Im}(\epsilon ^{}/\eta _0)`$, $`|\eta _0|`$ and $`\varphi _0`$ from $`|\eta _{00}/\eta _+|`$, $`\mathrm{\Delta }\varphi `$, $`|\eta _+|`$ and $`\varphi _+`$. These results are shown as the intermediate outputs in Table 2.
The second half. The values of $`\eta _0`$, $`\epsilon ^{}/\eta _0`$, $`\varphi _{SW}`$ and $`\text{Re}(\omega ^{})`$ obtained, supplemented with the value of $`\delta _2\delta _0`$, are used as inputs to find $`\text{Re}(\epsilon _2\epsilon _0)`$ and $`\text{Im}(\epsilon _2\epsilon _0)`$ with the help of Eqs.(5.13) and (5.14), and to find $`\text{Re}(\epsilon )`$ and $`\text{Im}(\delta )`$ with the help of Eq.(5.17). The values of $`\text{Re}(\epsilon )`$ and $`\text{Im}(\delta )`$ are in turn used to constrain $`\text{Re}(\delta \epsilon _0)`$ and $`\text{Im}(\epsilon +\epsilon _0)`$ through Eq.(5.2) and constrain $`\text{Re}(\delta \epsilon _{\mathrm{}}+x_{\mathrm{}}^{()})`$ through Eq.(5.5). The numerical results obtained are shown in Table 3. The value of $`\text{Re}(\epsilon \delta +\epsilon _{\mathrm{}}x_{\mathrm{}}^{()})`$, which is nothing but the value of $`d_L^{\mathrm{}}/2`$, is also shown.
## 7 Numerical analysis (2) -Constraints from the <br>CPLEAR results-
Immediately after CPLEAR Collaboration reported their preliminary result on the asymptotic leptonic asymmetries, $`d_{1,2}^{\mathrm{}}(t\tau _S)`$, we showed that this result, combined with the other relevant data available, could be used with the help of the Bell-Steinberger relation to constrain many of the $`CP`$$`T`$ and/or $`CPT`$-violating parameters introduced. The analysis went as follows. Assuming $`\text{Re}(x_{\mathrm{}}^{()})=0`$,<sup>g</sup><sup>g</sup>gIn most of the experimental analyses prior to those by CPLEAR Collaboration, either $`CPT`$ symmetry is taken as granted or no distinction is made between $`x_\mathrm{}+`$ and $`x_{\mathrm{}}`$, which implies that $`x_{\mathrm{}}^{()}`$ is presupposed to be zero implicitly. Accordingly, we identified $`x`$ used in with our $`x_{\mathrm{}}^{(+)}`$. Eqs.(5.6a) and (5.15) were used to find the values of $`\text{Re}(\epsilon )`$ and $`\text{Re}(\epsilon _{\mathrm{}})`$. The value of $`\text{Re}(\epsilon _{\mathrm{}})`$ was then used to constrain $`\text{Re}(\delta )`$ and $`\text{Im}(\delta )`$ through Eqs.(5.6b) and (5.16) respectively and all these values were combined with the value of $`\eta _0`$ to determine or constrain $`\text{Im}(\epsilon +\epsilon _0)`$ and $`\text{Re}(\epsilon _0)`$.
In order to appreciate the results obtained under the $`2\pi `$ dominance and to separately constrain, as far as possible, the parameters not yet separately constrained in the previous section, we now proceed to perform an analysis similar to the one explained above , with the new results \[13-16\] reported by CPLEAR Collaboration taken into account. In \[14-16\], CPLEAR Collaboration have defined two kinds of experimental asymmetries $`A_T^{exp}(t)`$ and $`A_\delta ^{exp}(t)`$ which are related to $`d_{1,2}^{\mathrm{}}(t)`$ and behave as
$`A_T^{exp}(t\tau _S)`$ $``$ $`4\text{Re}(\epsilon +\epsilon _{\mathrm{}}x_{\mathrm{}}^{()}),`$ (7.1a)
$`A_\delta ^{exp}(t\tau _S)`$ $``$ $`8\text{Re}(\delta ),`$ (7.1b)
and, by performing
1. fit to $`A_T^{exp}`$ under the assumption of $`\text{Re}(\epsilon _{\mathrm{}})=0`$ and $`x_{\mathrm{}}^{()}=0`$ ,
2. fit to $`A_\delta ^{exp}`$ , and
3. fit to both $`A_T^{exp}`$ and $`A_\delta ^{exp}`$ using as constraints the Bell-Steinberger relation and the PDG value of $`d_L^{\mathrm{}}`$ ,
succeeded in determining $`\text{Re}(\epsilon )`$, $`\text{Re}(\delta )`$, $`\text{Im}(\delta )`$, $`\text{Re}(\epsilon _{\mathrm{}})`$, $`\text{Im}(x_{\mathrm{}}^{(+)})`$ and/or $`\text{Re}(x_{\mathrm{}}^{()})`$ simultaneously.
Among the numerical outputs obtained by CPLEAR Collaboration, there are two pieces, $`\text{Re}(\epsilon )=(1.55\pm 0.35)\times 10^3`$ from and $`\text{Re}(\delta )(0.30\pm 0.33)\times 10^3`$ from , are in fact determined predominantly by the asymptotic values of $`A_T^{exp}(t)`$ and $`A_\delta ^{exp}(t)`$.<sup>h</sup><sup>h</sup>hIn contrast, the values of $`\text{Im}(x_{\mathrm{}}^{(+)})`$, $`\text{Re}(x_{\mathrm{}}^{()})`$ and $`\text{Im}(\delta )`$ obtained are sensitive to the behavior of $`A_T^{exp}(t)`$ and of $`A_\delta ^{exp}(t)`$ at $`t`$ comparable to $`\tau _S`$. One may therefore interpret these outputs as giving the values of $`A_T^{exp}(t\tau _S)/4`$ and $`A_\delta ^{exp}(t\tau _S)/8`$ respectively. Replacing Eq.(5.6a) with Eq.(7.1a), using Eq.(5.5) instead of Eq.(5.6b), and with the data listed in Table 4 as well as in Table 1 used as inputs, we perform an analysis similar to the previous one , and obtain the result shown in Table 5.
A couple of remarks are in order.
1. The assumption of $`x_{\mathrm{}}^{()}=0`$ has little influence numerically on determination of $`\text{Re}(\epsilon )`$, $`\text{Im}(\delta )`$ and $`\text{Im}(\epsilon +\epsilon _0)`$ and the error of these parameters is dominated by that of $`\eta _{000}`$.
2. Our constraint to $`\text{Re}(\epsilon _{\mathrm{}})`$ is better to be interpreted as a constraint to $`\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()})`$, the error of which is controlled dominantly by that of $`A_T^{exp}`$.
3. The error of $`\text{Re}(\delta )`$ and $`\text{Re}(\epsilon _0)`$ is also controlled dominantly by that of $`A_T^{exp}`$.
4. The numerical results we have obtained are fairly in agreement with those obtained by CPLEAR Collaboration in , except that we have not been able to separate $`\text{Re}(\epsilon _{\mathrm{}})`$ from $`\text{Re}(x_{\mathrm{}}^{()})`$.
## 8 Case study. -$`T`$ or $`CPT`$ violation ?-
In the analyses given in the previous sections, we have taken account of the possibility that any of $`CP`$$`T`$ and $`CPT`$ symmetries might violated in the $`K^0\overline{K^0}`$ system. Our numerical results shown in Table 3 and Table 5 indicate that $`CPT`$ symmetry appears consistent with experiments while $`T`$ symmetry appears not consistent with experiments. To confirm these observations, we now go on to perform a case study.
Case A. $`CPT`$ is a good symmetry.
Putting
$$\delta =\text{Re}(\epsilon _I)=\text{Re}(\epsilon _{\mathrm{}})=x_{\mathrm{}}^{()}=0,$$
Eqs.(5.2), (5.5), (5.6a,b), (5.8) and (5.13) reduce respectively to
$$\eta _I\epsilon +i\text{Im}(\epsilon _I),$$
(8.1a)
$$d_L^{\mathrm{}}2\text{Re}(\epsilon ),$$
(8.1b)
$$d_1^{\mathrm{}}(t\tau _S)4\text{Re}(\epsilon ),$$
(8.1c)
$$d_2^{\mathrm{}}(t\tau _S)0,$$
(8.1d)
$$d_L^{\mathrm{}}d_1^{\mathrm{}}(t\tau _S)/2,$$
(8.1e)
$$\epsilon ^{}/\eta _0\text{Re}(\omega ^{})\text{Im}(\epsilon _2\epsilon _0)e^{i\mathrm{\Delta }\varphi ^{}}/[|\eta _0|\mathrm{cos}(\delta _2\delta _0)].$$
(8.1f)
Eq.(8.1f) gives<sup>i</sup><sup>i</sup>iIt is to be noted that, if and only if $`CPT`$ symmetry is supplemented with the very accidental empirical fact $`\varphi _{SW}\delta _2\delta _0+\pi /2`$, one would have $`\text{Im}(\epsilon ^{}/\eta _0)0`$; it is therefore, as emphasized in , not adequate to assume this in a phenomenological analysis.
$$\text{Im}(\epsilon ^{}/\eta _0)=\text{Re}(\epsilon ^{}/\eta _0)\mathrm{tan}\mathrm{\Delta }\varphi ^{},$$
(8.2)
and the simplified version of the Bell-Steinberger relation, Eq.(5.17), gives<sup>j</sup><sup>j</sup>jEq.(8.3a) states that deviation of $`\varphi _0`$ from $`\varphi _{SW}`$ measures $`CPT`$ violation. This is equivalent to the more familiar statement: deviation of $`(2/3)\varphi _++(1/3)\varphi _{00}`$ from $`\varphi _{SW}`$ measures $`CPT`$ violation, because Eq.(5.9), supplemented with the experimental observation $`|\eta _{00}/\eta _+|1`$ and $`\mathrm{\Delta }\varphi 0`$, gives $`\varphi _0(2/3)\varphi _++(1/3)\varphi _{00}`$.
$`\varphi _0`$ $``$ $`\varphi _{SW},`$ (8.3a)
$`\text{Re}(\epsilon )`$ $``$ $`|\eta _0|\mathrm{cos}\varphi _{SW}.`$ (8.3b)
¿From the input data (Table 1 and Table 4) and the intermediate output data (Table 2), we observe the following:
(1) The experimental values of $`d_L^{\mathrm{}}`$, $`d_1^{\mathrm{}}(t\tau _S)`$ and $`d_2^{\mathrm{}}(t\tau _S)`$ are compatible with Eqs.(8.1d,e).
(2) The values of $`\text{Re}(\epsilon ^{}/\eta _0)`$, $`\text{Im}(\epsilon ^{}/\eta _0)`$, $`\varphi _0`$ and $`\delta _2\delta _0`$ are, as illustrated in Fig. 1, compatible with Eq.(8.2).
(3) The values of $`\varphi _0`$ and $`\varphi _{SW}`$ are compatible wiht Eq.(8.3a).
(4) The values of $`\text{Re}(\epsilon )`$ determined from Eqs.(8.1a) and (8.1b), $`(1.652\pm 0.029)\times 10^3`$ and $`(1.635\pm 0.060)\times 10^3`$, are compatible with each other and, as a weighted average, give
$$\text{Re}(\epsilon )(1.649\pm 0.026)\times 10^3,$$
(8.4)
which is compatible with $`(1.656\pm 0.014)\times 10^3`$ determined with the aid of the Bell-Steinberger relation Eq.(8.3b).<sup>k</sup><sup>k</sup>kEq.(5.15), with $`\text{Re}(\epsilon _{\mathrm{}})=0`$, yields $`(1.667\pm 0.048)\times 10^3`$.
(5) Eqs.(8.1a) and (8.1f) give
$`\text{Im}(\epsilon +\epsilon _0)`$ $``$ $`(1.570\pm 0.030)\times 10^3,`$ (8.5a)
$`\text{Im}(\epsilon _2\epsilon _0)`$ $``$ $`(3.02\pm 1.09)\times 10^4.`$ (8.5b)
Case B. $`T`$ is a good symmetry.<sup>l</sup><sup>l</sup>lThe possibility of $`CP/CPT`$ violation in the framework of $`T`$ symmetry was examined before by one of the present authors (S.Y.T) when the experimental results which upset $`CPT`$ symmetry (e.g., $`|\eta _{00}|`$ is nerely twice as large as $`|\eta _+|`$ !) had been reported. The same possibility was recently reconsidered by Bigi and Sanda .
Putting
$$\epsilon =\text{Im}(\epsilon _I)=\text{Im}(x_{\mathrm{}}^{(+)})=\text{Im}(x_{\mathrm{}}^{()})=0,$$
Eqs.(5.2), (5.5), (5.6a) and (5.13) reduce respectively to
$$\eta _I\delta +\text{Re}(\epsilon _I),$$
(8.6a)
$$d_L^{\mathrm{}}2\text{Re}(\delta )+2\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()}),$$
(8.6b)
$$d_1^{\mathrm{}}(t\tau _S)2\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()}),$$
(8.6c)
$$\epsilon ^{}/\eta _0\text{Re}(\omega ^{})\text{Re}(\epsilon _2\epsilon _0)e^{i(\mathrm{\Delta }\varphi ^{}+\pi /2)}/[|\eta _0|\mathrm{cos}(\delta _2\delta _0)].$$
(8.6d)
Eq.(8.6d) gives
$$\text{Im}(\epsilon ^{}/\eta _0)=\text{Re}(\epsilon ^{}/\eta _0)\mathrm{cot}\mathrm{\Delta }\varphi ^{},$$
(8.7)
and the simplified version of the Bell-Steinberger relation, Eq.(5.17), gives
$`\varphi _0`$ $``$ $`\varphi _{SW}\pm \pi /2,`$ (8.8a)
$`\text{Im}(\delta )`$ $``$ $`\pm |\eta _0|\mathrm{cos}\varphi _{SW}.`$ (8.8b)
¿From the input data (Table 1 and Table 4) and the intermediate output data (Table 2), we observe the following:
(1) As illustrated also in Fig.1, the values of $`\text{Re}(\epsilon ^{}/\eta _0)`$, $`\text{Im}(\epsilon ^{}/\eta _0)`$, $`\varphi _0`$ and $`\delta _2\delta _0`$ are not compatible with Eq.(8.7).
(2) The values of $`\varphi _0`$ and $`\varphi _{SW}`$ are not compatible with Eq.(8.8a).
(3) Eq.(8.6a) gives
$$\text{Im}(\delta )(1.570\pm 0.030)\times 10^3,$$
(8.9)
to be compared with $`\pm (1.656\pm 0.014)\times 10^3`$ determined with the aid of Eq.(8.8b).
(4) Eq.(8.6d) gives
$$\text{Re}(\epsilon _2\epsilon _0)(0.61\pm 3.12)\times 10^4,$$
(8.10)
while Eqs.(8.6a,b,c) give in turn
$`\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()})`$ $``$ $`(3.14\pm 1.40)\times 10^3,`$ (8.11a)
$`\text{Re}(\delta )`$ $``$ $`(1.51\pm 1.40)\times 10^3,`$ (8.11b)
$`\text{Re}(\epsilon _0)`$ $``$ $`(3.16\pm 1.40)\times 10^3.`$ (8.11c)
The observation (1) establishes the existence of direct $`CP/T`$ violation in the $`K^0\overline{K^0}`$ system .<sup>m</sup><sup>m</sup>mWe like to mention that Eq.(8.7) would become consistent with experiments if, say, $`\varphi _{00}`$ would prove to be away from $`\varphi _+`$ roughly by $`6^{}`$ or more. The observation (2), though subject to the validity of the Bell-Steinberger relation, also implies that $`CP/T`$ symmetry is violated in the $`K^0\overline{K^0}`$ system.
## 9 Summary and concluding remarks
In order to identify or search for violation of $`CP`$$`T`$ and $`CPT`$ symmetries in the $`K^0\overline{K^0}`$ system, parametrizing the mixing parameters and the relevant decay amplitudes in a convenient and well-defined way, we have, with the aid of the Bell-Steinberger relation and with all the relevant experimental data used as inputs, performed numerical analyses to derive constraints to the symmetry-violating parameters in several ways. The analysis given in Sec.6 is based on the data on $`2\pi `$ decays as well as the well measured leptonic asymmetry $`d_L^{\mathrm{}}`$, while, in the analysis given in Sec.7, the data on $`3\pi `$ and $`\pi ^+\pi ^{}\gamma `$ decays and on the newly measured leptonic asymmetries are also taken into account.
The numerical outputs of our analyses are shown in Table 3 and Table 5, and the main results may be summarized as follows:
(1) The $`2\pi `$ data directly give $`\text{Im}(\epsilon _2\epsilon _0)=(2.95\pm 1.13)\times 10^4`$ in general, or $`(3.02\pm 1.09)\times 10^4`$ if $`CPT`$ symmetry is assumed, where possible large uncertainty associated with $`\delta _2\delta _0`$ has been fully taken into account. This result indicates that $`CP`$ and $`T`$ symmetries are definitively violated in decays of $`K^0`$ and $`\overline{K^0}`$ into $`2\pi `$ states.
(2) The well-measured leptonic asymmetry $`d_L^{\mathrm{}}`$ directly gives $`\text{Re}(\epsilon \delta +\epsilon _{\mathrm{}}x_{\mathrm{}}^{()})=(1.635\pm 0.060)\times 10^3`$, which implies presumably that $`CP`$ and $`T`$ violations are present also in the the $`K^0\overline{K^0}`$ mixing (i.e., $`\text{Re}(\epsilon )0`$).<sup>n</sup><sup>n</sup>nOf course, $`d_L^{\mathrm{}}0`$ does not exclude $`CPT`$ violation (i.e., $`\text{Re}(\delta \epsilon _{\mathrm{}}+x_{\mathrm{}}^{()})0`$.)
(3) The Bell-Steinberger relation, with the $`2\pi `$ intermediate states alone taken into account, gives $`\text{Re}(\epsilon )=(1.656\pm 0.014)\times 10^3`$ and $`\text{Im}(\epsilon +\epsilon _0)=(1.566\pm 0.014)\times 10^3`$. If $`CPT`$ symmetry is assumed, $`\text{Re}(\epsilon )`$ is determined without recourse to the Bell-Steinberger relation to be $`(1.649\pm 0.026)\times 10^3`$. All these indicate that $`CP`$ and $`T`$ violations are present in the mixing parameters.
(4) The parameters, nonvanishing of which signals $`CP`$ and $`CPT`$ violations, have also been constrained. The $`2\pi `$ data alone directly give $`\text{Re}(\epsilon _2\epsilon _0)=(0.084\pm 0.328)\times 10^3`$ and, with the aid of the Bell-Steinberger relation, give $`\text{Im}(\delta )=(0.004\pm 0.027)\times 10^3`$, $`\text{Re}(\delta \epsilon _0)=(0.004\pm 0.026)\times 10^3`$ and $`\text{Re}(\delta \epsilon _{\mathrm{}}+x_{\mathrm{}}^{()})=(0.021\pm 0.062)\times 10^3`$. These results imply that there is no evidence for $`CPT`$ violation on the one hand and that $`CPT`$ symmetry is tested at best to the level of a few $`\times 10^5`$ on the other hand.
(5) The Bell-Steinberger relation, even with the intermediate states other than the $`2\pi `$ states taken into account, still allows one to determine $`\text{Re}(\epsilon )`$ and $`\text{Im}(\epsilon +\epsilon _0)`$ and to constrain $`\text{Im}(\delta )`$ to a level better than $`10^4`$. On the other hand, the constraint to $`\text{Re}(\delta )`$, $`\text{Re}(\epsilon _0)`$ and $`\text{Re}(\epsilon _{\mathrm{}}x_{\mathrm{}}^{()})`$ is a little loose and is at the level of a few $`\times 10^4`$.
The recent data reported by KTeV Collabotration and NA48 Collaboration are extremely remarkable in that they play a vital role in establishing $`\text{Im}(\epsilon _2\epsilon _0)0`$, and that this is at present the only piece which indicates ”direct violation” (in the sense defined in Sec.4) of $`CP`$ and $`T`$ symmetries and thereby unambiguously rules out superweak (or superweak-like) models of $`CP`$ violation.
The analyses done by CPLEAR Collaboration \[14-16\] are also very remarkable in particular in that they have succeeded in deriving constraint to $`\text{Re}(x_{\mathrm{}}^{()})`$, and in that they have determined $`\text{Re}(\epsilon +\epsilon _{\mathrm{}}x_{\mathrm{}}^{()})`$ and $`\text{Re}(\delta )`$ directly (i.e., without invoking the Bell-Steinberger relation) with accuracy down to the level of a few $`\times 10^4`$.<sup>o</sup><sup>o</sup>oAlvarez-Gaume, Kounnas and Lola further claim that the CPLEAR results allow one to conclude, without invoking Bell-Steinberger relation, that $`T`$ symmetry is violated independent of whether $`CP`$ and/or $`CPT`$ symmetries are violated or not.
It is expected that the new experiments at the facilities such as DA$`\mathrm{\Phi }`$NE, Frascati, will be providing data with such precision and quality that a more precise and thorough test of $`CP`$, $`T`$ and $`CPT`$ symmetries, and a test of the Bell-Steinberger relation as well, become possible . |
warning/0002/cond-mat0002387.html | ar5iv | text | # Evidence for a Goldstone Mode in a Double Layer Quantum Hall System
## Abstract
The tunneling conductance between two parallel 2D electron systems has been measured in a regime of strong interlayer Coulomb correlations. At total Landau level filling $`\nu _T=1`$ the tunnel spectrum changes qualitatively when the boundary separating the compressible phase from the ferromagnetic quantized Hall state is crossed. A huge resonant enhancement replaces the strongly suppressed equilibrium tunneling characteristic of weakly coupled layers. The possible relationship of this enhancement to the Goldstone mode of the broken symmetry ground state is discussed.
When two parallel two-dimensional electron systems (2DES) are sufficiently close together, interlayer Coulomb interactions can produce collective states which have no counterpart in the individual 2D systems . One of the simplest, yet most interesting, examples occurs when the total electron density, $`N_T`$, equals the degeneracy $`eB/h`$ of a single spin-resolved Landau level produced by a magnetic field $`B`$. In the balanced case (i.e. with layer densities $`N_1`$=$`N_2`$=$`N_T/2`$), the Landau level filling factor of each layer viewed separately is $`\nu =hN_T/2eB`$=$`1/2`$. If the separation $`d`$ between the layers is large, they behave independently and are well described as gapless composite fermion liquids. No quantized Hall effect (QHE) is seen. On the other hand, as $`d`$ is reduced, the system undergoes a quantum phase transition to an incompressible state best described by the total filling factor $`\nu _T`$=$`1/2+1/2`$=$`1`$. A quantized Hall plateau now appears at $`\rho _{xy}=h/e^2`$. Both Coulomb interactions and interlayer tunneling contribute to the strength of this QHE but there is strong evidence from experiment and theory that the incompressibility survives in the limit of zero tunneling. This remarkable collective state exhibits a broken symmetry, spontanteous interlayer phase coherence, and may be viewed as a new kind of easy-plane ferromagnet. The magnetization of this ferromagnet exists in a pseudospin space; electrons in one layer are pseudospin up, while those in the other layer are pseudospin down. Numerous interesting properties are anticipated, including linearly dispersing Goldstone collective modes (i.e. pseudospin waves), a finite temperature Kosterlitz-Thouless (K-T) transition, dissipationless transport for currents directed oppositely in the two layers, and bizarre topological defects in the pseudospin field. To date, most experimental results on this system have derived from electrical transport measurements although recently an optical study has been reported.
In this paper we report a new study of the double layer $`\nu _T`$=1 ferromagnetic quantum Hall state, and its transition at large layer separation to a compressible phase, using the method of tunneling spectroscopy. Earlier experiments have shown that there is a very strong suppression of the equilibrium tunneling between two widely separated parallel 2DESs at high magnetic field. This suppression is a result of the energetic penalty accompanying the rapid injection (or extraction) of an electron into the strongly correlated electron system produced by Landau quantization. Other than a small downward shift in energy produced by the excitonic attraction of a tunneled electron and the hole it leaves behind in the source layer, the measured tunneling spectrum is simply a convolution of the spectral functions in the individual layers. Here we show that when the layers are close enough together to support the bilayer $`\nu _T`$=1 QHE state, this is no longer the case. The strong suppression is replaced by a huge resonant enhancement of the tunneling. The appearance of this resonance suggests the existence of a soft collective mode of the double layer system which enhances the ability of electrons to tunnel. This may well be the predicted Goldstone mode of the broken symmetry ferromagnetic ground state at $`\nu _T`$=1.
The samples used in this experiment are $`\mathrm{GaAs}/\mathrm{Al}_\mathrm{x}\mathrm{Ga}_{1\mathrm{x}}\mathrm{As}`$ double quantum well (DQW) heterostructures grown by molecular beam epitaxy (MBE). Two 180Å GaAs quantum wells are separated by a 99Å $`\mathrm{Al}_{0.9}\mathrm{Ga}_{0.1}\mathrm{As}`$ barrier layer. The DQW is embedded in thick $`\mathrm{Al}_{0.3}\mathrm{Ga}_{0.7}\mathrm{As}`$ layers which contain Si doping sheets set back from the GaAs wells sufficiently far to produce 2D electron gases in each well with nearly equal electron densities of $`5.4\times 10^{10}cm^2`$. A square mesa, $`250\mu m`$ on a side, is patterned using standard photolithography. Gate electrodes deposited above and below this mesa allow control over the densities $`N_{1,2}`$ of the two 2D layers. The low temperature mobility of the as-grown sample is $`7.5\times 10^5cm^2/Vs`$ but this drops to $`2.5\times 10^5cm^2/Vs`$ when the layer densities are reduced to $`2\times 10^{10}cm^2`$. Ohmic contacts are placed at the ends of four arms extending outward from the central mesa. Using a selective depletion scheme, these contacts can be connected to both 2D layers in parallel or to either layer individually. Consequently, both conventional resistivity and interlayer tunneling measurements can be made on the same sample during a single cooldown from room temperature. Qualitatively identical tunneling data were obtained from three distinct samples taken from the same MBE wafer.
At zero magnetic field and low temperature the tunneling current-voltage ($`I`$-$`V`$) characteristics of our samples exhibit the simple sharp resonances characteristic of a high degree of momentum and energy conservation upon tunneling. The tunneling conductance $`dI/dV`$ exhibits a sharp peak which, for equal layer densities, is centered at zero interlayer bias voltage $`V`$. The width of this peak ranges from 0.15meV for $`N_1=N_2`$=$`5.4\times 10^{10}cm^2`$ to about 0.2meV at $`2.1\times 10^{10}cm^2`$ and reflects the static disorder in the system. The peak conductance is typically only about $`3\times 10^8\mathrm{\Omega }^1`$. The samples were designed to be weakly tunneling in order to avoid problems arising from the small sheet conductivities of the 2D layers which develop at high magnetic field.
Figure 1 displays the central result of this investigation. Four low temperature (T=40mK) tunneling conductance spectra are shown, each at a different total density $`N_T`$ in the bilayer system. In each case the individual layer densities were carefully matched by adjusting (via the gates) the symmetry and voltage location of the tunnel resonance at zero magnetic field. The traces were taken at different magnetic fields but each represents total Landau level filling factor $`\nu _T`$=1. For trace A the density is relatively high, $`N_T=10.9\times 10^{10}cm^2`$, and the well-known coulombic suppression of tunneling at the Fermi level (i.e. at $`V`$=0) is clearly evident. Trace B, taken at $`N_T=6.9\times 10^{10}cm^2`$, is similar to A except that the suppression effect is weaker and the overall spread in energy of the tunneling is less. This is the expected behavior since the inter-particle Coulomb energy falls with density. Trace C, at $`N_T=6.4\times 10^{10}cm^2`$, reveals a qualitatively new feature: a small yet sharp peak in $`dI/dV`$ at $`V=0`$. Finally, at the still lower density $`N_T=5.4\times 10^{10}cm^2`$, trace D shows that the peak has become enormous and dwarfs all other features in the tunnel spectrum. The height of this peak continues to grow as the density is reduced to $`3.2\times 10^{10}cm^2`$ where it exceeds even the zero magnetic field tunneling conductance by more than a factor of 10.
It is important to note that while the two layers have equal densities at $`V=0`$, the finite interlayer capacitance disrupts this balance at non-zero $`V`$. For the data in Fig. 1, this effect preserves $`\nu _T=1/2+1/2=1`$ since carriers are merely shifted from one layer to the other. We have found it possible to compensate for this effect, i.e. maintain the individual layers at fixed density, by adjusting the top and bottom gate voltage in linear proportion to the interlayer voltage $`V`$. This compensation alters the details of the tunneling characteristic at large $`V`$, but it has a negligible effect near $`V=0`$. Since our primary focus is the resonant peak at $`V=0`$, we shall for simplicity restrict the subsequent discussion to data obtained without compensating for the capacitive charge transfer.
Figure 2 shows the magnetic field dependence of the zero bias (i.e. $`V=0`$) tunneling conductance $`G_0`$, at T=40mK, for $`N_T=10.9\times 10^{10}`$ and $`4.2\times 10^{10}cm^2`$. In order to directly compare the two curves, the data is plotted against the inverse total filling factor $`\nu _T^1=eB/hN_T`$, instead of magnetic field. As expected, both curves exhibit quantum oscillations of the tunneling conductance at small magnetic field. These oscillations are less pronounced in the low density data owing to the reduced electron mobility. At high magnetic field, in the vicinity of $`\nu _T`$=1, the two data sets differ qualitatively. The high density data show that for magnetic fields above about $`\nu _T^1`$$``$$`0.3`$, the tunneling conductance is near zero. This again is the coulombic suppression characteristic of tunneling between two weakly coupled 2D electron systems at high magnetic field. By contrast, the low density data show an enormous enhancement of the tunneling around $`\nu _T=1`$. The enhancement appears strongest at, or very near, $`\nu _T=1`$, but it is clearly a very substantial effect over a wide range of filling factors.
The temperature dependence of the zero bias tunneling conductance is displayed in Fig. 3. Two sets of data are shown, one for $`N_T=10.9\times 10^{10}cm^2`$ and one for $`N_T=4.2\times 10^{10}cm^2`$. There is again a qualitative difference between the tunneling at high and low density. At high density the conductance falls with decreasing temperature. As reported previously, this dependence is consistent with simple thermal activation. The low density data behave in the opposite fashion, rising as the temperature falls. The rise becomes fairly steep around T=0.4K and then levels off below about 40mK.
Magneto-transport measurements on similar double quantum well samples have established the approximate location of the phase boundary between the incompressible $`\nu _T=1`$ quantized Hall phase and the compressible non-quantized Hall phase at large layer separation. In the limit of weak tunneling this boundary was found to be near $`d/\mathrm{}2`$. In this ratio $`d`$ is the center-to-center distance between the quantum wells and $`\mathrm{}=(\mathrm{}/eB)^{1/2}`$ is the magnetic length. In its as-grown state, i.e. when $`N_T=10.9\times 10^{10}cm^2`$, the present sample has $`d/\mathrm{}=2.4`$ at $`\nu _T=1`$. Reducing the density via gating to $`N_T=4.2\times 10^{10}cm^2`$, gives $`d/\mathrm{}=1.45`$. Control over the layer densities thus allows us to span the expected phase boundary using a single sample. The inset to Fig. 2 shows the zero bias tunneling conductance $`G_0`$ at $`\nu _T=1`$, at T=40mK, versus $`d/\mathrm{}`$. There is a sharp transition near $`d/\mathrm{}1.8`$ separating two very different tunneling regimes. For $`d/\mathrm{}>1.8`$ the zero bias conductance is suppressed and the tunneling spectra are qualitatively the same as seen in samples having negligible interlayer correlations. On the other hand, as $`d/\mathrm{}`$ falls below this critical value a resonant enhancement of the tunneling appears at zero bias. The magnitude of this peak grows continuously as $`d/\mathrm{}`$ falls.
The rough agreement between the critical $`d/\mathrm{}`$ value found in transport experiments and that reported here in tunneling studies suggests that they reflect the same phase transition. To investigate this, resistivity measurements were performed on the sample. As anticipated, a quantized Hall effect does develop at low density. The dotted trace in Fig. 2 shows the observed minimum in $`\rho _{xx}`$ at $`N_T=4.2\times 10^{10}cm^2`$ and T=40mK. The relative weakness of this many-body integer QHE may simply be due to the low mobility of the sample ($`2\times 10^5cm^2/Vs`$) at this reduced density. A second possibility is more interesting: It is well known that for some interaction-driven quantized Hall states Arrhenius behavior of the resistivity $`\rho _{xx}`$ is not observed until the temperature is reduced well below the measured activation energy. For example, in conventional single layer 2D systems at $`\nu =1`$, spin flip activation energies of order 40K are common, but can only be observed at temperatures below a few Kelvin. This effect originates in the low energy of long wavelength spin waves in the system. The activation energy measured in transport reflects the large exchange contribution to the energy required to create a well-separated quasiparticle-quasihole pair. By contrast, at long wavelengths the collective mode energy is the much smaller bare Zeeman energy. Thermal population of large numbers of these modes is apparently effective in suppressing the onset of a thermally activated resistivity. The same phenomenon is observed in the $`\nu _T=1`$ QHE in double layer systems. Now, however, the collective modes are the predicted pseudospin waves. At $`q=0`$ these modes are gapless, but only in the non-physical limit of zero tunneling. In real samples a gap at $`q=0`$ is expected, the size of which is determined by both the single particle tunneling energy $`\mathrm{\Delta }_{SAS}`$ and an interlayer capacitive charging energy. Using our estimate of $`\mathrm{\Delta }_{SAS}=90\mu K`$ and published estimates of the charging energy, the long wavelength pseudospin wave energy is only about 70mK. It is therefore not so surprising that the $`\rho _{xx}`$ minimum at $`\nu _T=1`$ is weak at 40mK and absent above about 200mK. We note in passing that the estimated $`\mathrm{\Delta }_{SAS}`$ in our sample is only $`1.3\times 10^6`$ of the typical Coulomb energy $`e^2/ϵ\mathrm{}`$ at $`N_T=4.2\times 10^{10}cm^2`$. This is by far the smallest tunneling strength of any sample reported exhibiting the $`\nu _T=1`$ bilayer QHE state.
A qualitative explanation for the tunneling enhancement reported here can be constructed from known aspects of double layer systems at $`\nu _T=1`$. At high densities, when the layers are weakly coupled, the zero bias tunneling is heavily suppressed. The energetic penalty associated with tunneling in this case arises because an electron attempting to tunnel is totally ”unaware” of the strong correlations present within the layer it is about to enter. In the $`\nu _T=1`$ bilayer QHE state just the opposite is true; the very essence of the state lies in the strong interlayer correlation it contains. Indeed, the so-called $`\mathrm{\Psi }_{111}`$ wavefunction believed to represent this QHE state may be viewed as a collection of correlated interlayer excitons. An electron in one layer is always opposite a hole in the other layer. It is not implausible that this interlayer correlation enhances tunneling. This simplistic view, while appealing, does not address the sharply resonant nature of the enhancement. The experimental observation of a narrow peak in the tunneling conductance at zero interlayer voltage suggests that there is a collective mode near zero energy which can transfer charge between the two layers. The predicted pseudospin Goldstone mode does exactly this. This mode involves oscillations of the pseudospin magnetization both in the $`xy`$-plane and in the $`z`$-direction of pseudospin space. This latter component resembles an antisymmetric interlayer plasma oscillation. In the absence of tunneling no charge is transferred between the layers and the mode energy vanishes in the long wavelength, $`q=0`$ limit. In a real sample however, the finite tunneling amplitude leads to an energy gap at $`q=0`$ and, most importantly, allows the collective mode to transfer charge. The estimated gap ($`70`$mK=6$`\mu `$eV) is much smaller than the observed width of the tunnel resonance (typically 150$`\mu `$eV at low temperature) so it is not surprising that a single peak in dI/dV appears at V=0. From this point of view, tunneling appears to provide direct spectroscopic evidence for the Goldstone mode of the broken symmetry QHE state at $`\nu _T=1`$.
In conclusion, we have examined tunneling in double layer 2D electron systems in which interlayer correlations are very strong. A dramatic resonant enhancement of the zero bias tunneling conductance is observed when the system crosses the phase boundary from a compressible fluid into the $`\nu _T=1`$ ferromagnetic quantized Hall state. Although detailed understanding of the tunnel spectra is lacking, it seems likely that the observed resonance is intimately connected with the Goldstone mode of the broken symmetry state.
We are indebted to A.H. MacDonald, S.M. Girvin, and A.C. Gossard for numerous useful discussions and to the National Science Foundation and the Department of Energy for financial support. One of us (I.B.S.) acknowledges the Department of Defense for a National Defense Science and Engineering Graduate Fellowship. |
warning/0002/hep-th0002113.html | ar5iv | text | # Chaotic behavior in a 𝑍₂×𝑍₂ field theoryThis work is supported in part by funds provided by the U.S. Department of Energy (D.O.E.) under cooperative research agreement #DF-FC02-94ER40818.
## I Introduction
Nonlinearity plays an important role in field theory. It is responsable for the presence of interactions and may enter the game allowing interesting situations, as is the case for instance when the system engenders spontaneous symmetry breaking. In this case the classical equations of motion may give rise to interesting field configurations such as kinks, vortices and monopoles, depending on the particular model in consideration. See for instance Refs. for details.
In the standard route to defect formation, one usually search for static field configurations that solve the equations of motion and present finite energy. In the simplest case of a single real scalar field, a system given in terms of a potential that presents $`Z_2`$ symmetry may support topological solutions in $`1+1`$ dimensions when the potential has at least two degenerate minima. In the case of two real scalar fields, the system is richer and may support distinct types of topological solutions, as we illustrate below.
Nonlinearity also plays an important role on chaotic behavior of systems. For linear systems, the qualitative nature of the behavior does not change when one changes their parameters. However, for systems governed by nonlinear dynamics one may find examples where small change in a parameter can lead to dramatic changes in both the quantitative and qualitative behavior of the system. See for instance Refs. for details.
There are distinct routes to chaos in field theory, and here we shall follow the point of view introduced in Ref. . In this case one investigates field theoretical models under the assumption of spatially homogeneous field configurations. We think of field configurations whose space variations are much smaller than the corresponding time variations, and so we treat the fields as depending only on time. This point of view has been explored in several different works, as for instance in Refs. and in some works therein. We notice that these works deal with different models, describing Abelian and non Abelian gauge fields in the presence of spontaneous symmetry breaking. In these cases the symmetry to be broken is continuum, local, and the system supports vortices or monopoles. In the present work, however, we shall explore another system, simpler, that presents discrete symmetry and is described by a couple of real scalar fields. Recent investigations have shown that this system presents interesting properties and some soliton solutions have also been found, as we comment on below.
The main motivation of the present work is to investigate the presence of chaos in a system that describes two real scalar fields engendering the discrete $`Z_2\times Z_2`$ symmetry. We shall do this by comparing the case where one supposes the field configurations to be spatially homogeneous with the more familiar situation which considers static configurations. The last case is appropriate for searching for soliton solutions, and we comment on that in the next Sec. II. We consider the case of spatially homogeneous fields in Sec III where we search for chaotic behavior. As we also comment on, the subject of this paper is related to recent issues , which ask for instance whether the chaotic behavior found in flat spacetime persists during the cosmological expansion. In this paper we deal with a model which is defined in flat spacetime and may be seen as the results of an expanding FRW universe, valid under the approximation of very slow expansion rate. Furthermore, the present investigation lands very naturally to the context of hybrid inflation, where one requires real scalars such as the inflaton field and at least another scalar field, which couples to the inflaton field in a way similar to the one we consider in this work. We end the paper by summarizing the results in Sec. IV.
## II General considerations
The system we study in this paper is a $`Z_2\times Z_2`$ model of two real fields $`\varphi (x,t)`$ and $`\chi (x,t)`$. The potential that specifies the system is defined in terms of a single real parameter, $`r`$, which controls the number of minimum energy states of that potential. As we are going to show, when $`r`$ changes sign the potential changes form, from the case with four minima $`(r>0)`$ to the case where only two minima are present $`(r<0)`$. For $`r>0`$, spontaneous symmetry breaking appears in both the $`\varphi `$ and $`\chi `$ directions in the $`(\varphi ,\chi )`$ plane. For $`r<0`$ we have spontaneous symmetry breaking only in the $`\varphi `$-field direction, and this last situation is similar to the case of hybrid inflation where a first-order nonthermal phase transition after preheating seems to be present, as recently considered in Ref. .
The field theory that we consider is described by the Lagrangian density
$$=\frac{1}{2}_\alpha \varphi ^\alpha \varphi +\frac{1}{2}_\alpha \chi ^\alpha \chi V(\varphi ,\chi )$$
(1)
We are using standard notation, with $`\mathrm{}=c=1`$, and $`V(\varphi ,\chi )`$ is the potential, which one supposes to be given by
$$V(\varphi ,\chi )=\frac{1}{2}\left(\frac{h}{\varphi }\right)^2+\frac{1}{2}\left(\frac{h}{\chi }\right)^2$$
(2)
where $`h=h(\varphi ,\chi )`$ is a smooth function of the two fields $`\varphi `$ and $`\chi `$. Here it obeys
$$\frac{1}{\mu }h(\varphi ,\chi )=\stackrel{~}{h}(\varphi ,\chi )=r\left(\frac{1}{3}\varphi ^3\varphi \right)+\varphi \chi ^2$$
(3)
The parameter $`\mu `$ is real, with dimension of energy, and $`r`$ is another real parameter, dimensionless.
In $`1+1`$ dimensions this system presents interesting soliton solutions , which has been used in applications in condensed matter and in field theory . The equations of motion for $`\varphi =\varphi (x,t)`$ and $`\chi =\chi (x,t)`$ are given by
$`\varphi _{tt}\varphi _{xx}+h_\varphi h_{\varphi \varphi }+h_\chi h_{\chi \varphi }`$ $`=`$ $`0`$ (4)
$`\chi _{tt}\chi _{xx}+h_\varphi h_{\varphi \chi }+h_\chi h_{\chi \chi }`$ $`=`$ $`0`$ (5)
These are the equations we shall deal with in the following. For simplicity, however, we rewrite them in terms of dimensionless variables $`\stackrel{~}{t}=\mu t`$ and $`\stackrel{~}{x}=\mu x`$ to get
$`\varphi _{\stackrel{~}{t}\stackrel{~}{t}}\varphi _{\stackrel{~}{x}\stackrel{~}{x}}+\stackrel{~}{h}_\varphi \stackrel{~}{h}_{\varphi \varphi }+\stackrel{~}{h}_\chi \stackrel{~}{h}_{\chi \varphi }`$ $`=`$ $`0`$ (6)
$`\chi _{\stackrel{~}{t}\stackrel{~}{t}}\chi _{\stackrel{~}{x}\stackrel{~}{x}}+\stackrel{~}{h}_\varphi \stackrel{~}{h}_{\varphi \chi }+\stackrel{~}{h}_\chi \stackrel{~}{h}_{\chi \chi }`$ $`=`$ $`0`$ (7)
In the case of static fields we have $`\varphi =\varphi (\stackrel{~}{x})`$ and $`\chi =\chi (\stackrel{~}{x})`$. We change $`\stackrel{~}{x}y`$, for simplicity, and now the equations of motion become
$`{\displaystyle \frac{d^2\varphi }{dy^2}}`$ $`=`$ $`2r^2(\varphi ^21)\varphi +2(r+2)\varphi \chi ^2`$ (8)
$`{\displaystyle \frac{d^2\chi }{dy^2}}`$ $`=`$ $`2(\chi ^2r)\chi +2(r+2)\varphi ^2\chi `$ (9)
which are the equations we deal with when searching for soliton solutions. It is interesting to see that these equations of motion (8) and (9) can be solved by configurations that obey the pair of first-order differential equations
$`{\displaystyle \frac{d\varphi }{dy}}`$ $`=`$ $`r(\varphi ^21)+\chi ^2`$ (10)
$`{\displaystyle \frac{d\chi }{dy}}`$ $`=`$ $`2\varphi \chi `$ (11)
Solutions that obey this pair of first-order equations are BPS solutions . They are stable configurations that minimize the energy, as explicitly shown in . See Ref. for further comments on BPS solutions, and for showing explicitly that the system defined by the potential (2) is the bosonic portion of a supersymmetric theory.
In the case of spatially homogeneous fields we have $`\varphi =\varphi (\stackrel{~}{t})`$ and $`\chi =\chi (\stackrel{~}{t})`$. We change $`\stackrel{~}{t}t`$, for simplicity, and here we get the equations of motion
$`{\displaystyle \frac{d^2\varphi }{dt^2}}`$ $`=`$ $`2r^2(1\varphi ^2)\varphi 2(r+2)\varphi \chi ^2`$ (12)
$`{\displaystyle \frac{d^2\chi }{dt^2}}`$ $`=`$ $`2(r\chi ^2)\chi 2(r+2)\varphi ^2\chi `$ (13)
These are the equations we have to deal with when searching for a chaotic behavior in the time evolution. It is interesting to see that the above equations have some analogies with the equations investigated in Ref. . However, in Ref. the system under consideration is the Abelian-Higgs model. This is the relativistic generalization of the Ginzburg-Landau theory of superconductivity, and it is well known that it presents vortex solutions . The reason for the similarity between eqs. (12) and (13) and the equations investigated in seems to rely on the assumptions introduced in Ref. .
We follow the aim of this paper, which is to investigate the chaotic behavior of eqs. (12) and (13), together with the study of the static solutions of eqs. (8) and (9). We start presenting some general considerations regarding the static solutions. Our system is identified by the potential $`V(\varphi ,\chi )=\mu ^2\stackrel{~}{V}(\varphi ,\chi )`$, where
$$\stackrel{~}{V}(\varphi ,\chi )=\frac{1}{2}r^2(\varphi ^21)^2+r(\varphi ^21)\chi ^2+\frac{1}{2}\chi ^4+2\varphi ^2\chi ^2$$
(14)
The potential depends on the real parameter $`r`$, and then the behavior of the system is directly related to such parameter. The main characteristic of this potential, which certainly strongly affect both the static solutions and the chaotic time behavior, is its different shape for $`r`$ positive and negative. In fact, when $`r>0`$ the potential contains four vacuum states in the $`(\varphi ,\chi )`$ plane, while for $`r<0`$ there are only two vacuum states, and in principle we expect a different chaotic behavior for the dynamical trajectories that follows in accordance with eqs. (12) and (13), as well as distinct features for the static solutions. It is important to say that in the two particular cases when $`r=1`$ and when $`r=2`$ the two fields decouples and so we expect no chaotic behavior in these two cases. First we show that, from the point of view of the topological soliton solutions admitted by Eqs. (8) and (9), the case $`r>0,r1`$ is richer than the case of $`r<0,r2`$. When $`r`$ is positive, $`r1`$, the four vacuum states are given by
$`v_1`$ $`=`$ $`(1,\mathrm{\hspace{0.17em}0}),v_2=(1,\mathrm{\hspace{0.17em}0})`$ (15)
$`v_3`$ $`=`$ $`(0,\sqrt{r}),v_4=(0,\sqrt{r})`$ (16)
For $`r`$ negative, $`r2`$, we have the two vacuum states
$$\overline{v}_1=(1,\mathrm{\hspace{0.17em}0}),\overline{v}_2=(1,\mathrm{\hspace{0.17em}0})$$
(17)
As it was already shown , the energy of the stable BPS solutions are obtained according to the values of $`\stackrel{~}{h}(\varphi ,\chi )`$ in the vacuum states.
For $`r`$ positive, $`r1`$ we name $`h_i`$ as the value of $`h`$ at the vacuum state $`v_i`$, and so we have
$`h_1`$ $`=`$ $`(2/3)\mu r`$ (18)
$`h_2`$ $`=`$ $`(2/3)\mu r`$ (19)
$`h_3`$ $`=`$ $`h_4=0`$ (20)
This means that we have two type of BPS solutions, with energies $`E_i=|\mu |\epsilon _i,i=1,2`$, where
$`\epsilon _1`$ $`=`$ $`{\displaystyle \frac{4}{3}}r,(r>0,r1)`$ (21)
$`\epsilon _2`$ $`=`$ $`{\displaystyle \frac{2}{3}}r,(r>0,r1)`$ (22)
In the first case the solutions connect the vacuum states $`(\pm 1,0)`$ and in the second case pair of vacuum states belonging to different axes. In the topological sector defined by the two vacuum states $`(1,0)`$ and $`(1,0)`$ there are pairs of analytical solutions given by
$$\varphi (y)=\mathrm{tanh}[r(y\overline{y})],\chi (y)=0$$
(23)
and by
$`\varphi (y)`$ $`=`$ $`\mathrm{tanh}[2(y\overline{y})]`$ (24)
$`\chi (y)`$ $`=`$ $`\pm {\displaystyle \frac{\sqrt{r2}}{\mathrm{cosh}[2(y\overline{y})]}}`$ (25)
Here $`\overline{y}`$ stands for the center of the topological solution. The first pair (23) is valid for $`r>0`$ and the second one for $`r>2`$. We see that the second pair of solutions (24) and (25) gets to the first pair (23) in the limit $`r2`$.
For $`r`$ negative, $`r2`$, there is just one type of BPS solution, with energy $`\overline{E}_1=|\mu |\overline{\epsilon }_1`$, where
$$\overline{\epsilon }_1=\frac{4}{3}|r|,(r<0,r2)$$
(26)
The first pair of solutions given by Eq. (23) is also a pair of solutions in the case $`r<0`$.
It is clear that for $`r>0,r1`$ the situation is richer than in the case $`r<0,r2`$, at least from the point of view of the above topological soliton solutions. In the first case the system admits two type of topological BPS solutions, while in the second case it has only one. We also expect a difference between the two $`r>0`$ and $`r<0`$ regimes to be present in the chaotic behavior of the eqs. (12) and (13).
Another interesting feature of this model concerns the masses of the $`\varphi `$ and $`\chi `$ fields. From the potential $`V(\varphi ,\chi )`$ we can introduce the matrix
$$\left(\genfrac{}{}{0pt}{}{V_{\varphi \varphi }^vV_{\varphi \chi }^v}{V_{\chi \varphi }^vV_{\chi \chi }^v}\right)$$
(27)
where $`V_{\varphi \varphi }=^2V/\varphi \varphi `$ and so forth. The superscript $`v`$ stands for substituting the fields for their vacuum values once the derivatives are done. For the potential introduced in Eq. (14) we see that
$$V_{\varphi \chi }=V_{\chi \varphi }=4\mu ^2(r+2)\varphi \chi $$
(28)
and so $`V_{\varphi \chi }^v=V_{\chi \varphi }^v`$ vanishes for every vacuum state, despite the sign of $`r`$. This means that we can read the (square) mass of each field directly from the matrix (27). We have at most two different mass values. They are
$`m_\varphi ^2`$ $`=`$ $`4\mu ^2r^2m_\chi ^2=\mathrm{\hspace{0.17em}4}\mu ^2(r>0,\varphi \mathrm{axis})`$ (29)
$`m_\varphi ^2`$ $`=`$ $`4\mu ^2rm_\chi ^2=4\mu ^2r(r>0,\chi \mathrm{axis})`$ (30)
and also
$$\overline{m}_\varphi ^2=4\mu ^2r^2\overline{m}_\chi ^2=4\mu ^2(r<0,\varphi \mathrm{axis})$$
(31)
where we inform the axis of the vacuum state used to obtain the respective masses. We note that for $`r<0`$ the $`\chi `$ field develops no symmetry breaking, and its mass does not depend on $`r`$. This is the way this specific $`Z_2\times Z_2`$ system behaves, and this behavior may have interesting connections with models of hybrid inflation . We remark that there are other $`Z_2\times Z_2`$ systems, as for instance the ones investigated recently in Ref. , which are also defined with a single real parameter, but that present distinct behavior, with different sets of masses and vacuum states. They certainly lead to richer scenarios, and may give other informations on the role the parameter $`r`$ plays in such models.
## III Chaotic behavior
In this Section we investigate the chaotic behavior of the system introduced in Sec. II. We focus our attention to the pair of Eqs. (12) and (13). These equations can be seen as the equations of motion that follow from the Hamiltonian
$$H=\frac{1}{2}\dot{p}_\varphi ^2+\frac{1}{2}\dot{p}_\chi ^2+\stackrel{~}{V}(\varphi ,\chi )$$
(32)
They describe the motion of a classical particle in the bidimensional potential $`\stackrel{~}{V}(\varphi ,\chi )`$. We remark that this choice of Hamiltonian makes the energy dimensionless. The equations of motion can be written in first order form, in terms of the canonical variables $`(\varphi ,p_\varphi )`$ and $`(\chi ,p_\chi )`$,
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{H}{p_\varphi }}=p_\varphi `$ (33)
$`\dot{p}_\varphi `$ $`=`$ $`{\displaystyle \frac{H}{\varphi }}=2r^2(1\varphi ^2)\varphi 2(r+2)\varphi \chi ^2`$ (34)
$`\dot{\chi }`$ $`=`$ $`{\displaystyle \frac{H}{p_\chi }}=p_\chi `$ (35)
$`\dot{p}_\chi `$ $`=`$ $`{\displaystyle \frac{H}{\chi }}=2(r\chi ^2)\chi 2(r+2)\varphi ^2\chi `$ (36)
In the rest of the paper we study for which values of the parameter $`r`$ and of the energy $`E`$ the trajectories $`\varphi (t),\chi (t),p_\varphi (t)`$, and $`p_\chi (t)`$ show a regular or a chaotic behavior.
The motion is always bounded, i.e. confined to a finite region whose size changes according to the energy $`E`$ of the particle. In Fig. 1 we report the potential contour plot defined by $`\stackrel{~}{V}(\varphi ,\chi )=E`$; the different values of the energy $`E`$ are reported in the figure. Fig. 1a refers to $`r=5`$ while Fig. 1b and Fig. 1c refer to $`r=1`$ and $`r=3`$, respectively.
For $`r<0`$ we note the presence of two minima $`(\pm 1,0)`$, corresponding to the two vacuum states previously discussed, and a saddle point at the origin. The potential vanishes at the two minima, while at the origin $`(0,0)`$ it assumes the value $`r^2/2`$ (12.5 for case $`r=5`$ in Fig. 1a). Dynamical trajectories, e.g. solutions of Eqs. (33)-(36) are confined to one of the two minima regions if $`E<12.5`$, while they can jump from one side to the other when $`E>12.5`$. For $`r=2`$ the two fields decouple and the numerical investigation shows no chaotic behavior.
The situation for $`r>0`$ is qualitatively different because we have four minima $`\overline{V}=0`$ at $`(\pm 1,0)`$ and at $`(0,\pm \sqrt{r})`$, corresponding to the four vacuum states previously discussed, and a local maximum $`\overline{V}=r^2/2`$ at the origin $`(0,0)`$. There are also four saddle points with $`\varphi 0`$ and $`\chi 0`$; they are given by
$$\varphi ^2=\frac{r}{2(r+1)},\chi ^2=\frac{r^2}{2(r+1)}$$
(37)
In the saddle points the potential gets the value
$$\overline{V}=\frac{r^2}{2(r+1)}$$
(38)
The case $`r=1`$ is a particular case because the potential presents $`Z_4`$ symmetry (see Fig. 1b); the two fields decouple and the numerical simulations show no chaotic behavior. When $`r1`$ trajectories can have different behaviors according to the energy $`E`$. We discuss the case $`r=3`$ in Fig. 1c: when $`E1.125`$ the four minima are unconnected regions and the motion remains confined around the minimum in which we start the trajectory. As soon as $`E>1.125`$ (see Fig. 1c) the trajectory crosses the saddle point \[see Eq. (38)\] between the minima and wonders from one minimum to another. The region around the relative maximum at the origin $`(0,0)`$ is still not allowed. When $`E>r^2/2=4.5`$ the dynamical trajectories can cross the origin and the four mimina are all directly interconnected.
In order to study the dynamical behavior of our system for different values of $`r`$ and energy $`E`$, we integrate numerically Eqs. (33)-(36) using a fourth order symplectic algorithm with a time step $`\mathrm{\Delta }t=0.001`$ . The time step has been determined in order to keep the error in energy conservation below $`\mathrm{\Delta }E/E=10^8`$ for any value of $`r`$ and $`E`$; the results are stable respect to a further reduction of $`\mathrm{\Delta }t`$. For any conservative system such as the one we are considering, the Hamiltonian is an integral of motion, and the energy conservation restricts trajectories to lie on a three-dimensional surface $`H(\varphi ,\chi ,p_\varphi ,p_\chi )=E`$ in the four-dimensional phase space $`(\varphi ,\chi ,p_\varphi ,p_\chi )`$. We can obtain a graphical information about the system by plotting the intersection of this three-dimensional surface with a plane. These kind of plot is called Poincaré surface of section, and gives an indication about the dynamical behavior of the system. To define the surface of section in our case we follow a trajectory and we plot $`\varphi `$ and $`\chi `$ each time $`p_\varphi =0`$. Regular regions will appear as a series of points (a mapping) which lie on a one dimensional curve (invariant KAM curve), while chaotic regions will appear as a scatter of points limited to a finite area due to energy conservation .
In Fig. 2 and in Fig. 3 we show the surface of section in the plane $`(\varphi ,\chi )`$ for $`r=5`$ and for $`r=3`$, respectively. Each figure is obtained from $`100`$ different trajectories followed for a time $`t=500`$. The trajectories are generated from random uniformly-distributed initial conditions $`\varphi _0,p_{\varphi }^{}{}_{0}{}^{},\chi _0,p_{\chi }^{}{}_{0}{}^{}`$.
For $`r=5`$ we plot the two cases $`E=10`$ and $`E=20`$ (Fig. 2a and Fig. 2b respectively). When $`E=10`$ the two minima are separated and the system exhibits only regular behavior (the surface of section shows only invariant KAM curves). When $`E=20`$ the two minima are connected and the surface of section contains both regular and stochastic regions. The volume of the stochastic region increases with increasing energy, as it will be clear when we calculate Lyapunov exponents.
For $`r=3`$ the situation is richer. We report the surface of section for three different energies $`E=0.5`$, $`E=2`$, and $`E=6`$. When $`E=0.5`$ (Fig. 3a) the four mimina are not interconnected and we only have invariant KAM curves. When $`E=2`$ (Fig. 3b) neighbour minima regions are connected. Regular and stochastic regions cohexist in the surface of section and are interwingled in a very complicated way. The scatter of points representing the chaotic region fill the area between regular curves, and then smaller regular islands are imbedded in the chaotic sea. The white regions correspond to regular region, and these would be filled by regular curves if we could increase even more the number of initial conditions considered in the figure. For energy $`E=6`$ (Fig. 3c) the trajectory can cross the center and the resulting surface of section shows an increase of the chaotic area respect to the case $`E=2`$.
The cases $`r=1`$ and $`r=2`$ (which we do not report in figure) are two very particular cases. The two equations decouples, and indeed the Poincaré surface of section show only KAM curves for every value of the energy $`E`$.
The information we can get from Poincaré surface of section is qualitative. A way to quantify the chaotic behavior of a system is by calculating its Lyapunov exponents. Chaos is defined in terms of the dyamical behavior of pairs of orbits which initially are close together in the phase space. The Lyapunov exponents are given by the rate of exponential divergence of close orbits and are defined from the long term evolution of an initial infinitesimal volume $`\mathrm{\Omega }`$ in the phase space from the following formula
$$\lambda _i=\underset{t\mathrm{}}{lim}\lambda _i(t)=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{ln}\frac{d_i(t)}{d_i(0)}$$
(39)
where $`d_i(t)`$ are the lenght of the principal axes of the ellipsoid $`\mathrm{\Omega }`$ at time $`t`$. There is one Lyapunov exponent for each dimension of the phase space and the flow is chaotic if at least one Lyapunov exponent is positive. For a two-dimensional Hamiltonian system we expect two positive $`\lambda _1,\lambda _2`$ and two negative $`\lambda _3,\lambda _4`$ Lyapunov exponents with the following symmetry $`\lambda _3=\lambda _2`$ and $`\lambda _4=\lambda _1`$, because of the symplectic structure of Hamiltonian systems.
To compute Lyapunov exponents for different values of $`r`$ and energy $`E`$ we use the standard method developped in Ref. . The average number of time steps to use in order to have a good convergence of the results is about $`10^710^8`$, which means $`t=10^410^5`$. The symmetries $`\lambda _3=\lambda _2`$ and $`\lambda _4=\lambda _1`$ are always perfectly verified in the numerical results; this is a good check for the integration codes. Fig. 4 shows the numerical calculation of the two largest Lyapunov exponents for case $`r=3`$ and $`E=5.5`$; we report the quantity $`\lambda _i(t)`$ for $`i=1,2`$ as a function of time. The top panel in Fig. 4 shows two different trajectories both starting from initial conditions in the chaotic region. After a transient time, $`\lambda _1(t)`$ converges for both trajectories to a constant value of about $`0.3`$, while $`\lambda _2(t)`$ converges to $`0.01`$. To obtain $`\lambda _i`$ we evaluate numerically the limit in formula (39) considering a time average of $`\lambda _i(t)`$ over the time interval $`t=210^4310^4`$. The bottom panel of Fig. 4 shows the behavior of $`\lambda _1(t)`$ and $`\lambda _2(t)`$ for a different initial condition, this time chosen to be in the regular region of the phase space; $`\lambda _1(t)`$ and $`\lambda _2(t)`$ do not reach a finite asymptotic value but tend to zero for large $`t`$, showing that the Lyapunov exponents are equal to zero.
For each value of $`r`$ and energy we consider $`N=100`$ different trajectories with random initial conditions $`\varphi _0,p_{\varphi }^{}{}_{0}{}^{},\chi _0,p_{\chi }^{}{}_{0}{}^{}`$. We compute $`\lambda _i`$ for each trajectory and we count the number $`N_c`$ of chaotic trajectories (i.e. with positive $`\lambda _1`$, as in Fig. 4a). In the following we report the ratio $`R=N_c/N`$ as an indicator of the measure of volume of phase space which is chaotic, and the value of $`\lambda _1`$ averaged over the $`N_c`$ chaotic trajectories. We do not report $`\lambda _2`$ which is, for each value of $`r`$ and energy, one order of magnitude smaller than $`\lambda _1`$.
We consider separately the two cases $`r<0`$ and $`r>0`$ for which we expect to find a qualitative difference in the chaotic properties of Eqs. (12) and (13). We discussed above the topological soliton solutions of Eqs. (8) and (9) for these two cases and we found that case $`r>0`$ is richer than case $`r<0`$. In fact, for $`r>0,r1`$ the system supports two topological sectors, admitting two different type of topological BPS solutions; for $`r<0,r2`$ there is just one topological sector. This difference is due to the presence of four or two minima in the potential $`V(\varphi ,\chi )`$, for $`r>0`$ or $`r<0`$ respectively. We expect the difference in the shape of the potential to have important implications also on the chaotic behavior of our system. In Fig. 5 and Fig. 6 we plot the ratio $`R=N_c/N`$ and $`\lambda _1`$, as function of energy for the two cases $`r<0,r2`$ and $`r>0,r1`$, respectively. The system exhibits an order to chaos transition as a function of energy for any value of $`r`$, but for $`r=2`$ and $`r=1`$. These two cases are particular cases because the two field decouple and the system degenerates to two systems of a single field each one; $`R`$ and $`\lambda _1`$ are identically equal to zero in the whole energy range and are not reported in figures.
We show the typical behavior for $`r`$ negative in Fig. 5. We report $`R`$ and $`\lambda _1`$ vs. energy for $`r=5`$; other values of $`r`$ have the same qualitative behavior. The onset of chaos is for energy equal to $`r^2/2`$ ($`12.5`$ in figure), when both minima can be visited by a single trajectory and $`R`$ and $`\lambda _1`$ start to be different from zero. For larger values of energy both $`R`$ and $`\lambda _1`$ increase with the energy.
The situation is different for $`r`$ positive. In Fig. 6 we plot $`R`$ and $`\lambda _1`$ vs. $`E`$ for two different values of $`r`$, $`r=1.5`$ (dashed line) and $`r=3`$ (dashed-dotted line). The behavior for $`r=1.5`$ and $`r=3`$ is qualitatively similar, although shifted in energy. Chaos starts when the energy overcomes the barrier between two minima \[which is given in accordance with Eq. (38)\], and it keep increasing as function of the energy until a new back bending appears, both in $`R`$ and in $`\lambda _1`$. This behavior is to be connected to a stabilization effect occurring when the trajectory can cross the origin since both $`R`$ and $`\lambda _1`$ have local minima at $`E=r^2/2`$. Although the chaotic behavior for $`r=1.5`$ and for $`r=3`$ are qualitatively similar, the absolute value of $`R`$ and $`\lambda _1`$ is as larger as closer $`r`$ gets to unity. This result can be of interest to application since the absolute value of $`r`$ also controls the masses of the two fields and may be used in models of hybrid inflation.
The behavior at $`r>0`$ is therefore different from the behavior at $`r<0`$. The most important qualitative difference between these two cases is that for $`r>0`$ we can find localized regions in energy where suppression of chaotic behavior appears for increasing energies. Such behavior is not present for $`r`$ negative, and this is directly related to the reduction of the number of minima in this last case. Similarly, we have shown that for static solutions the number of topological sectors changes from two to one when $`r`$ changes from positive to negative values, respectively.
## IV Comments and conclusions
In this work we have investigated the presence of chaotic behavior in a system of two real scalar fields that engenders the $`Z_2\times Z_2`$ symmetry. The system is defined by a potential that is controlled by a single real parameter, $`r`$. This parameter can be positive or negative, and each case leads to different behavior. For static field configurations for instance, the system suports two dintinct topological sectors for $`r`$ positive, $`r1`$, and only one when $`r`$ is negative, $`r2`$.
Within the context of chaotic behavior, we have shown that chaos is present almost everywhere in parameter space, with distinct qualitative behavior for $`r>0`$ and for $`r<0`$. The system is richer for $`r`$ positive, and this is directly related to the presence of the four minima when $`r>0`$. For $`r>0`$ we can find localized regions in energy where suppression of chaotic behavior appears for increasing energy values, a fact that is absent for $`r`$ negative. On the other hand, in the case of static solution there are two topological sectors for $`r>0`$, in contraposition with the single sector that the system supports for $`r<0`$.
There are other models that engender the $`Z_2\times Z_2`$ symmetry, and that are also governed by a similar real parameter $`r`$. Some examples appear in Ref. , and they can also be studied within the context of chaotic behavior, under the assumption of spatially homogeneous field configurations. Such investigations can introduce further light toward a better understanding of the connection between chaotic behaviors and the parameter $`r`$ that specifies the physical properties of the system.
As we have already commented on, the present investigations may be of interest to inflationary cosmology, in connection with issues raised in the recent works . The present results seem to be valid in an expanding FRW universe when the rate of expansion is very small. The chaotic behavior that we have found may have been present in early times and may have played some role in the cosmic evolution. A study in which the coupling to gravity is fully considered, and other related issues are presently under cosideration.
###### Acknowledgements.
We would like to thank Michel Baranger for many interesting comments, and Giuseppe Politi for a critical reading of the manuscript. V.L. thanks BLANCEFLOR Foundation and CNR for financial support. We also thank Center for Theoretical Physics, MIT, and Department of Physics, Harvard University, for the kind hospitality.
FIGURE CAPTIONS
Fig. 1. Potential contours plot for different values of E. a), b), and c) refer to $`r=5,1`$ and $`3`$, respectively.
Fig. 2. Poincaré surface of section for $`r=5`$ in $`(\varphi ,\chi )`$ plane. $`E=10`$ in a) and $`E=20`$ in b).
Fig. 3. Poincaré surface of section for $`r=3`$ in $`(\varphi ,\chi )`$ plane. $`E=0.5`$ in a), $`E=2`$ in b), and $`E=6`$ in c).
Fig. 4. Lyapunov Exponents $`\lambda _1(t)`$ and $`\lambda _2(t)`$ vs. time for $`r=3`$ and $`E=5.5`$. Two chaotic trajectories are plotted in the top panel, while a regular trajectory is plotted in the bottom panel.
Fig. 5. Fraction R of the phase space which is chaotic and largest Lyapunov exponent $`\lambda _1`$ vs. $`E`$ for $`r=5`$
Fig. 6. Same as in Fig.5 for two positive values of $`r`$, namely $`r=1.5`$ (dashed) and $`r=3`$ (dotted-dashed). |
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