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warning/0002/hep-ph0002292.html	ar5iv	text	# 1 Introduction ## 1 Introduction The goal of superstring phenomenology at its present stage is to develop the tools and methodology to connect between string theory and experimental data. It is clear that to understand the mechanism which selects the string vacuum a non–perturbative formulation is needed. However, it is rather likely that detailed confrontation with the experimental data will have to rely on perturbative means. For this purpose the tools to construct realistic string models and the methodology to extract their phenomenological implications must be further developed. The first mandatory task of superstring phenomenology is to produce string solutions which are as realistic as possible with present day technology. The goal in this regard is to construct string models that aim to reproduce the phenomenological data provided by the Standard Model spectrum. Moreover, with the lack of substantial experimental evidence for any extension of the Standard Model, the most desired solution would be one that reproduces solely the Standard Model. Once the first goal is achieved the subsequent goal is to extract possible experimental signatures, beyond the Standard Model, which may provide further evidence for specific string models, in particular, and for string theory, in general. In practice, of course, it makes sense to try to extract the experimental consequences of the theory at every stage of its development. This, for example, was the drive behind much of the superstring inspired activity that followed the seminal Candelas et al. paper . Pursuing the minimalist approach, and taking the Standard Model data as the guide toward understanding the basic building blocks of nature, one is compelled to assess that grand unification structures are relevant in nature. Proton decay constraints then imply the big desert scenario, and gravity becomes important only near the Planck scale. In this eventuality the perturbative heterotic string, which naturally accommodates the grand unification structures with chiral matter, is the relevant framework. Over the past decade the free fermionic formulation of the heterotic string has been utilized to derive the most realistic string models to date . A large number of three generation models have been constructed, which differ in their detailed phenomenological characteristics. All these models share an underlying $`Z_2\times Z_2`$ orbifold structure, which naturally gives rise to three generation models with the standard $`SO(10)`$ embedding of the Standard Model spectrum . In this respect the phenomenological success of the free fermionic models can be regarded as indicating the relevance of the $`Z_2\times Z_2`$ orbifold structure in nature. Furthermore, recently, and for the first time since the advent of superstring phenomenology, it was demonstrated that free fermionic models also produce Minimal Standard Heterotic String Models (MSHSM) . In such models the low energy spectrum, which is charged under the Standard Model gauge group, consists solely of the spectrum of the Minimal Supersymmetric Standard Model. It should be emphasized that it is not suggested that one of the free fermionic models that has been constructed to date is the correct string vacuum. Indeed, such a claim would at least require a derivation of the detailed fermion mass spectrum, as well as an understanding of the dynamics which break supersymmetry. However, a plausible interpretation of the phenomenological success of the free fermionic models is that the true string vacuum is in the neighborhood of these models. This interpretation would then single out, for example, the $`Z_2\times Z_2`$ orbifold as the relevant string compactification. The general lesson in this respect is to extract the string structures which are relevant for the Standard Model phenomenological data. Subsequent to achieving the mandatory task of demonstrating the phenomenological viability of a particular class of heterotic string compactification, trying to extract possible experimental signatures beyond the Standard Model becomes more compelling. Various such possible experimental signatures, inspired from string theory, have been discussed in the past. Among them: additional gauge bosons, exotic matter, and specific patterns of supersymmetry breaking. Our aim here is to use the same tools, that have been used to derive realistic string models, to try to extract the experimental signatures beyond the Standard Model. Such experimental signatures then have the advantage of being “derived” rather than “inspired” from string theory. In this paper, we focus on the possibility of additional gauge bosons, beyond the Standard Model. The possibility of such additional gauge structure has of course been discussed extensively in the past, mainly in the context of additional symmetries which arise in $`SO(10)`$ and $`E_6`$ grand unification . However, as our analysis demonstrates the most likely additional gauge bosons to arise from realistic string models are not of this origin. Therefore, of particular interest in our discussion will be the additional gauge bosons which are of particular string origin, i.e., those that do not arise in Grand Unified theories. ## 2 Free fermionic phenomenology The analysis of the free fermionic models is conducted in two steps. In the first step the free fermionic model building rules are used to construct a consistent three generation string vacuum. Subsequently one extracts the full massless spectrum as well as the cubic level and higher order non–renormalizable terms in the superpotential. At this stage the tools used are perturbative heterotic string theory techniques. The superstring derived three generation models contain numerous massless vector–like states, some of which carry fractional electric charge. The models typically also contain a number of additional $`U(1)`$ symmetries in the observable sector, plus a hidden gauge group which is a subgroup of the original hidden $`E_8`$ of the heterotic string. One additionally finds that many of these three generation models contain an anomalous $`U(1)_A`$ symmetry, which generates a Fayet–Iliopoulos term, $`ϵ{\displaystyle \frac{g_s^2M_P^2}{192\pi ^2}}\mathrm{Tr}Q^{(A)},`$ (2.1) where $`\mathrm{Tr}Q^{(A)}0,`$ is the trace of the $`U(1)_A`$ charge over all the massless fields. The Fayet–Iliopoulos term breaks supersymmetry near the Planck scale, and destabilizes the string vacuum. Supersymmetry is restored and the vacuum is stabilized if there exists a direction in the scalar potential $`\varphi =_i\alpha _i\varphi _i`$ which is $`F`$–flat and also $`D`$–flat with respect to the non–anomalous gauge symmetries and in which $`_iQ_i^A\|\alpha _i\|^2`$ and $`ϵ`$ are of opposite sign. If such a direction exists it will acquire a vacuum expectation value (VEV) cancelling the anomalous $`D`$–term, restoring supersymmetry and stabilizing the string vacuum. Since the fields corresponding to such a flat direction typically also carry charges for the non–anomalous $`D`$–terms, a non–trivial set of constraints, $`D_A`$ $`=`$ $`{\displaystyle \underset{m}{}}Q_m^{(A)}\|\phi _m\|^2+ϵ=0,`$ (2.2) $`D_i`$ $`=`$ $`{\displaystyle \underset{m}{}}Q_m^{(i)}\|\phi _m\|^2=0.`$ (2.3) on the possible choices of VEVs is imposed. These scalar VEVs will in general break some, or all, of the additional symmetries spontaneously. Additionally one must insure that the supersymmetric vacuum is also $`F`$–flat. Each superfield $`\mathrm{\Phi }_m`$ (containing a scalar field $`\phi _m`$ and chiral spin–$`\frac{1}{2}`$ superpartner $`\psi _m`$) that appears in the superpotential imposes further constraints on the scalar VEVs. $`F`$–flatness will be broken (thereby destroying spacetime supersymmetry) at the scale of the VEVs unless, $$F_m\frac{W}{\mathrm{\Phi }_m}=0;W=0,$$ (2.4) where $`W`$ is the superpotential which contains cubic level and higher order non–renormalizable terms. The higher order terms have the generic form $$<\mathrm{\Phi }_1^f\mathrm{\Phi }_2^f\mathrm{\Phi }_3^b\mathrm{}\mathrm{\Phi }_N^b>.$$ Some of the fields appearing in the non–renormalizable terms will in general acquire a non–vanishing VEV by the anomalous $`U(1)`$ cancellation mechanism. Thus, in this process some of the non–renormalizable terms induce effective renormalizable operators in the effective low energy field theory wherein either all fields or all fields but one are replaced with VEVs. One must insure that such terms do not violate supersymmetry at an unacceptable level. In practice, however, the studies performed to date have been restricted to the case in which supersymmetry remains unbroken to all orders of non–renormalizable terms. Thus, the second stage of the string model building analysis is conducted by analyzing the $`F`$– and $`D`$–flat directions. An important advance of the last few years has been the development of systematic techniques for the analysis of exact $`F`$– and $`D`$–flat directions . Additionally one may impose other phenomenological constraints on the scalar VEVs. For example, in demonstrating the existence of a free fermionic MSHSM, we required that a set of fields which induces the decoupling of all non–MSSM states, acquire a non–vanishing VEV along the $`F`$– and $`D`$–flat directions. The string model building analysis outlined above can be regarded as aiming to achieve the first goal of superstring phenomenology. Namely, to reproduce the data provided by the Standard Model. ## 3 Additional gauge symmetries in free fermionic models In this section we discuss the different classes of additional gauge symmetries that are obtained in the free fermionic models prior to the analysis of flat directions. For a given three generation string model, and a specific flat direction, the resulting string vacuum may give rise to additional matter and gauge bosons, which are beyond the Minimal Supersymmetric Standard Model. For example, the hidden sector may give rise to massless matter states, which are not charged with respect to the Standard Model gauge group, and which interact with the Standard Model states only via horizontal $`U(1)`$ symmetries. Such hidden matter states may have interesting cosmological implications and may serve as dark matter candidates. Similarly, for a given flat direction the string vacuum may contain a combination of the horizontal $`U(1)`$ symmetries, which remains unbroken. It is this type of unbroken $`U(1)`$ symmetry that we aim to study in this paper. Thus, for a given flat direction the first task is to extract the combinations of the $`U(1)`$ symmetries which remain unbroken. Some of the resulting combinations may be entirely hidden. Namely, the Standard Model states will not be charged under them. Such hidden combinations may therefore be less interesting from an experimental point of view. However, there may also exist unbroken combinations of the horizontal $`U(1)`$ symmetries, under which the Standard Model states are charged. It is precisely this type of unbroken $`U(1)`$ symmetries that are of enormous experimental and phenomenological interest. Furthermore, in a given string model the charges of the Standard Model states under such an unbroken $`U(1)`$ symmetry are completely specified. Consequently, the phenomenology of the additional $`Z^{}`$ in a given string model is specified, up to some educated assumptions on the scale of $`Z^{}`$ breaking and the strength of its coupling. Both of these assumptions will of course eventually be relaxed. The free fermionic models are constructed by specifying a set of boundary conditions basis vectors and the one–loop GSO projection coefficients . The basis vectors, $`𝐛_k`$, span a finite additive group $`\mathrm{\Xi }=_kn_k𝐛_k`$ where $`n_k=0,\mathrm{},N_{z_k}1`$, with $`N_{z_k}`$ the smallest positive integer such that $`N_{z_k}𝐛_k=\stackrel{}{0}`$ (mod 2). The physical massless states in the Hilbert space of a given sector $`\alpha \mathrm{\Xi }`$, are obtained by acting on the vacuum with bosonic and fermionic operators and by applying the generalized GSO projections. The $`U(1)`$ charges, $`Q(f)`$, with respect to the unbroken Cartan generators of the four dimensional gauge group, which are in one to one correspondence with the $`U(1)`$ currents $`f^{}f`$ for each complex fermion f, are given by: $`Q(f)={\displaystyle \frac{1}{2}}\alpha (f)+F(f),`$ (3.1) where $`\alpha (f)`$ is the boundary condition of the world–sheet fermion $`f`$ in the sector $`\alpha `$, and $`F_\alpha (f)`$ is a fermion number operator counting each mode of $`f`$ once (and if $`f`$ is complex, $`f^{}`$ minus once). For periodic fermions, $`\alpha (f)=1`$, the vacuum is a spinor in order to represent the Clifford algebra of the corresponding zero modes. For each periodic complex fermion $`f`$ there are two degenerate vacua $`\|+,\|`$ , annihilated by the zero modes $`f_0`$ and $`f_{0}^{}{}_{}{}^{}`$ and with fermion numbers $`F(f)=0,1`$, respectively. The four dimensional gauge group in the three generation free fermionic models arises as follows. The models can in general be regarded as constructed in two stages. The first stage consists of the NAHE set of boundary conditions basis vectors, which is a set of five boundary condition basis vectors, $`\{\mathrm{𝟏},𝐒,𝐛_1,𝐛_2,𝐛_3\}`$ . The NAHE set is a common set in the three generation models that we discuss here. The gauge group after imposing the GSO projections induced by the NAHE set basis vectors is $$SO(10)\times SO(6)^3\times E_8$$ with $`N=1`$ supersymmetry. The space–time vector bosons that generate the gauge group arise from the Neveu–Schwarz sector and from the sector $`\mathrm{𝟏}+𝐛_1+𝐛_2+𝐛_3`$. The Neveu–Schwarz sector produces the generators of $`SO(10)\times SO(6)^3\times SO(16)`$. The sector $`\zeta \mathrm{𝟏}+𝐛_1+𝐛_2+𝐛_3`$ produces the spinorial of $`\mathrm{𝟏}28`$ of $`SO(16)`$ and completes the hidden gauge group to $`E_8`$. At the level of the NAHE set the sectors $`𝐛_1`$, $`𝐛_2`$ and $`𝐛_3`$ produce 48 multiplets, 16 from each, in the $`\mathrm{𝟏}6`$ representation of $`SO(10)`$. We remark that in order to understand the origin of the various additional $`U(1)`$ symmetries that may appear in free fermionic models it is often useful to consider a version of the NAHE set, which is extended by adding the basis vector $`X`$ with periodic boundary conditions for the right–moving complex fermions $`\{\overline{\psi }^{1,\mathrm{},5},\overline{\eta }^{1,2,3}\}`$ . With this additional boundary basis vector the four dimensional gauge symmetry is extended to $$E_6\times U(1)^2\times SO(4)^3\times E_8.$$ One can regard this set as starting with a toroidally compactified model generated by the set of basis vectors $`\{\mathrm{𝟏},𝐒,𝐗,\zeta \}`$. The right–moving gauge group with this set is $`SO(12)\times E_8\times E_8`$, with $`N=4`$ space–time supersymmetry. The basis vectors $`𝐛_1`$ and $`𝐛_2`$ are then used to break $`N=4N=1`$ supersymmetry, and to reduce the gauge symmetry to $`E_6\times U(1)^2\times SO(4)^3\times E_8`$. The $`U(1)`$ combination produced by the world–sheet currents $`\overline{\eta }^1\overline{\eta }^1^{}+\overline{\eta }^2\overline{\eta }^2^{}+\overline{\eta }^3\overline{\eta }^3^{}`$ becomes the $`U(1)`$ symmetry in the decomposition of $`E_6`$ under $`SO(10)\times U(1)`$. The realistic free fermionic models can be regarded as starting with this set, but changing the sign of the GSO phase $`c(\genfrac{}{}{0pt}{}{\zeta }{𝐗})`$. With this GSO phase change the $`\mathrm{𝟏}6+\overline{\mathrm{𝟏}6}`$ generators in the adjoint of $`E_6`$ are projected out. The right–moving gauge group in this case becomes $`SO(10)\times U(1)_A\times U(1)^2\times SO(4)^3\times E_8`$, with $`U(1)_A`$ being anomalous . The second stage of the free fermionic basis construction consists of adding to the NAHE set three (or four) additional boundary condition basis vectors. These additional basis vectors reduce the number of generations to three chiral generation, one from each of the basis vectors $`𝐛_1`$, $`𝐛_2`$ and $`𝐛_3`$, and simultaneously break the four dimensional gauge group. The $`SO(10)`$ is broken to one of its subgroups $`SU(5)\times U(1)`$, $`SO(6)\times SO(4)`$ or $`SU(3)\times SU(2)\times U(1)^2`$. Similarly, the hidden $`E_8`$ symmetry is broken to one of its subgroups by the basis vectors which extend the NAHE set. This hidden $`E_8`$ subgroup may, or may not, contain $`U(1)`$ factors which are not enhanced to a non–Abelian symmetry. On the other hand, the flavor $`SO(6)^3`$ symmetries in the NAHE–based models are always broken to flavor $`U(1)`$ symmetries, as the breaking of these symmetries is correlated with the number of chiral generations. Three such $`U(1)_j`$ symmetries are always obtained in the NAHE based free fermionic models, from the subgroup of the observable $`E_8`$, which is orthogonal to $`SO(10)`$. These are produced by the world–sheet currents $`\overline{\eta }\overline{\eta }^{}`$ ($`j=1,2,3`$), which are part of the Cartan sub–algebra of the observable $`E_8`$. Additional unbroken $`U(1)`$ symmetries, denoted typically by $`U(1)_j`$ ($`j=4,5,\mathrm{}`$), arise by pairing two real fermions from the sets $`\{\overline{y}^{3,\mathrm{},6}\}`$, $`\{\overline{y}^{1,2},\overline{\omega }^{5,6}\}`$ and $`\{\overline{\omega }^{1,\mathrm{},4}\}`$. The final observable gauge group depends on the number of such pairings. Our interest in this paper is in additional gauge bosons that arise from the horizontal flavor symmetries. That is, in additional vector bosons which arise from combinations of world–sheet $`U(1)`$ currents of the Cartan subalgebra. The generators of such additional $`U(1)`$’s all arise from the Neveu–Schwarz sector. Before proceeding, however, we briefly discuss possible non–Abelian extensions of the Standard Model in these models, and postpone detailed analysis on these possibilities to future work. It is already clear that from the unbroken subgroup of $`SO(10)`$, we can obtain the traditional left–right symmetric extensions of the Standard Model. These originate from the $`SO(6)\times SO(4)`$ type models or from left–right symmetric models in which $`SO(10)`$ is broken to $`SU(3)\times U(1)\times SO(4)`$ at the string level. Additional sources of possible non–Abelian enhancement may arise from combinations of the boundary conditions basis vectors which extend the NAHE set. In some of the three generation model one finds combinations of the additional basis vectors $`Y=n_\alpha \alpha +n_\beta \beta +n_\gamma \gamma ,`$ (3.2) for which $`Y_LY_L=0`$ and $`Y_RY_R8`$. Such a combination may produce additional space–time vectors bosons, depending on the GSO projections. In these cases, some combination of the $`U(1)`$ generators of the four dimensional Cartan sub–algebra is enhanced to a non–Abelian gauge symmetry. Often it is found that this is a combination of the flavor symmetries, which is family universal, and combines with the $`U(1)_{BL}`$ generator to produce a baryonic, or leptonic, non–Abelian gauged symmetry . Such symmetries may therefore play an important role in insuring proton stability, but their phenomenological viability still needs to be studied. We will not discuss this type of enhanced symmetries further here and delegate more detailed studies to future work. To summarize, in the spirit of the minimalist approach pursued here, in this paper we are interested in additional gauge bosons that arise from the unbroken Cartan generators of the four dimensional gauge group. In particular, we are interested in possible combinations of the $`U(1)`$ symmetries, which remain unbroken by a set of flat directions that cancels the anomalous $`U(1)`$ $`D`$–term. ## 4 $`Z^{}`$s in free fermionic models We now turn to discuss the extra $`Z^{}`$ symmetries that appear in free fermionic models. The first objective is to study the additional $`U(1)`$ symmetries which appear prior to the analysis of flat directions. The subsequent objective is to determine which combinations of $`U(1)`$ symmetries remain unbroken after the analysis of flat directions, and possibly after imposing additional phenomenological constraints that are required by the Standard Model data. Such unbroken $`U(1)`$ combinations then come closer to being a prediction of the string models. The final goal is of course to extract which possible $`U(1)`$ combinations remain unbroken after the wealth of Standard Model experimental data is satisfied. Such an extra $`U(1)`$ combination then is truly a prediction of a specific string vacua. However, short of this ambitious and still unachievable goal, we can already at this stage extract the general characteristics of $`U(1)`$ combinations that may remain unbroken in detailed $`F`$– and $`D`$–flat solutions. The first type of $`Z^{}`$ symmetry that has been considered in the context of free fermionic models has been the $`U(1)`$ combination $$Q_Z^{}=\frac{BL}{2}\frac{2}{3}T_{3_R},$$ (4.1) which is embedded in $`SO(10)`$ and is orthogonal to the Standard Model weak–hypercharge. The phenomenology of this class of extra $`U(1)`$’s, as well as its family–universal extensions in the context of $`E_6`$ string inspired phenomenology, have been discussed extensively in the past . As we discussed above, in the free fermionic models the additional $`U(1)`$ symmetry (aside from (4.1) which is embedded in $`E_6`$ is given by the family universal combination of the horizontal symmetries, given by $`U(1)_{E_6}=U(1)_1+U(1)_2+U(1)_3.`$ (4.2) However, there are several reasons to argue that these particular $`U(1)`$ combinations, (4.1) and (4.2), in the free fermionic models cannot remain unbroken to low energies. In the first place, one often finds (although not always ) that the family universal $`U(1)`$ which is embedded in $`E_6`$ is anomalous and is therefore broken by the flat direction VEVs. Second, the scale of the breaking of the $`U(1)`$ symmetry which is embedded in $`SO(10)`$ is associated with the see–saw scale, which is needed to suppress the left–handed neutrino masses. Thus, the requirement of sufficiently small neutrino masses implies that this particular $`U(1)`$ symmetry cannot remain unbroken to low energies. The natural question is then which additional $`U(1)`$ symmetries, beyond the weak–hypercharge of the Standard Model, can remain unbroken to low energies. As discussed in the introduction the string models under considerations often contain an anomalous $`U(1)`$ symmetry. In those cases most, or all, of the horizontal $`U(1)`$ symmetries in the observable sector are broken by the choices of flat directions. Additionally, one has to impose plausible phenomenological constraints, like the decoupling of exotic fractionally charged states and quasi–realistic fermion mass spectrum, which may further result in the breaking of the observable horizontal symmetries. The choice of flat directions may leave unbroken $`U(1)`$ symmetries in the hidden sector, but those are of less interest from an experimental and phenomenological perspective. We further remark that in left–right symmetric models, with $`SU(3)\times U(1)\times SU(2)_L\times SU(2)_R`$ as the unbroken subgroup of $`SO(10)`$ at the string scale, one finds models in which all the horizontal $`U(1)`$ symmetries are anomaly free . That is, in these models there is no anomalous $`U(1)`$ symmetry. Consequently these semi–realistic string vacua are supersymmetric and anomaly free without the need for scalar VEVs which break some of the horizontal $`U(1)`$ symmetries. However, the phenomenology of this class of string models has not been studied extensively and one may expect that imposing plausible phenomenological constraints will necessitate some Planck scale VEVs. Therefore, in this paper we focus on string models that do contain an anomalous $`U(1)`$ symmetry. ## 5 $`Z^{}`$ in the FNY model As our concrete illustrative example of a $`Z^{}`$ appearing in a string model we consider the string derived model of ref. . We will refer to this model as the FNY model. The $`F`$– and $`D`$–flat directions in this model were studied in detail in refs. . There it was shown that there exist for this model flat directions which result in the decoupling of all the massless exotic fractionally charged states by the scalar VEVs. This is achieved due to the fact that in this model there exist cubic level superpotential terms, in which the exotic fractionally charged states are coupled to a set of $`SO(10)`$ singlets. Thus, assigning non–vanishing VEVs to this set of $`SO(10)`$ singlets results in all of the fractionally charged exotic states receiving mass of the order of the Fayet–Iliopoulos term. It was further shown that for these flat directions all the additional states beyond the spectrum of the Minimal Supersymmetric Standard Model receive masses from up to quintic order terms in the superpotential. Therefore, in this model all the states that are beyond the MSSM and which are charged with respect to the Standard Model gauge group decouple from the massless spectrum at or slightly below the string scale. This string model therefore provides the first known example of a Minimal Standard Heterotic–String Model (MSHSM). It should be emphasized that this does not indicate that the FNY string model is the correct string vacuum, nor is it our intention to claim that the FNY model passes all of the phenomenological constraints imposed by the Standard Model data. However, what we think is a reasonable lesson to extract is that the success of producing a MSHSM, as well as the other unique phenomenological characteristics of the free fermionic models, like the standard $`SO(10)`$ embedding of the weak–hypercharge, may be taken as suggesting that the correct string vacuum may indeed exist in the vicinity of the free fermionic point in the string moduli space. The details of the FNY string model, its massless spectrum, and superpotential terms up to sixth order are given in ref. . Here, for completeness, we only discuss the features of the model which are relevant for our discussion, and give in Table 1 the relevant states and charges in the effective low energy field theory. Prior to the analysis of flat directions the observable gauge group of the FNY model is: $`SU(3)_C\times SU(2)_L\times U(1)_C\times U(1)_L\times U(1)^6`$ <sup></sup><sup></sup>*$`U(1)_C=\frac{3}{2}U(1)_{BL}`$; $`U(1)_L=2U(1)_{T_{3_R}}`$., and the hidden gauge group is: $`SO(4)\times SU(2)\times U(1)^4`$. The Standard Model weak–hypercharge is given by $`U(1)_Y=\frac{1}{3}U(1)_C+\frac{1}{2}U(1)_L`$. The sectors $`𝐛_1`$, $`𝐛_2`$ and $`𝐛_3`$ produce the three light generations. Electroweak Higgs doublets $`\{h_{1,2,3},\overline{h}_{1,2,3}\}`$ arise from the Neveu–Schwarz sector, and $`H_{34}`$, $`H_{41}`$ from the sectors $`𝐛_3+\alpha \pm \beta `$ and $`𝐛_1+𝐛_2+𝐛_4+\alpha \pm \beta `$. Prior to rotating the anomaly into a single $`U(1)_\mathrm{A}`$, six of the FNY model’s twelve $`U(1)`$ symmetries are anomalous: Tr$`U_1=24`$, Tr$`U_2=30`$, Tr$`U_3=18`$, Tr$`U_5=6`$, Tr$`U_6=6`$ and Tr$`U_8=12`$. Thus, the total anomaly can be rotated into a single $`U(1)_\mathrm{A}`$ defined by $$U_A4U_15U_2+3U_3+U_5+U_6+2U_8.$$ (5.1) The five orthogonal linear combinations, $`U_1^{^{}}`$ $`=`$ $`2U_1U_2+U_3\text{ ; }U_2^{^{}}=U_1+5U_2+7U_3;`$ $`U_3^{^{}}`$ $`=`$ $`U_5U_6\text{ ; }U_4^{^{}}=U_5+U_6U_8;`$ (5.2) $`U_5^{^{}}`$ $`=`$ $`12U_1+15U_29U_3+25U_5+25U_6+50U_8,`$ are all traceless. A particular flat solution in the FNY model is given by the set of fields $$\{\mathrm{\Phi }_{12},\mathrm{\Phi }_{23},\overline{\mathrm{\Phi }}_{56},\mathrm{\Phi }_4,\mathrm{\Phi }_4^{},\overline{\mathrm{\Phi }}_4,\overline{\mathrm{\Phi }}_4^{},H_{31},H_{38},H_{23},V_{40},H_{28},V_{37}\}.$$ (5.3) As discussed in ref. with this set of VEVs all of the exotic states beyond the MSSM receive heavy mass from cubic or quintic order terms in the superpotential. Detailed investigation of the fermion mass texture which is generated by the $`F`$– and $`D`$–flat solutions has been performed in ref. . The analysis was performed for flat directions which utilize only non–Abelian singlet VEVs. The solution in Eq. (5.3) also contains non–Abelian fields, and was shown to be flat to all orders in ref. . Quick examination of the non–renormalizable terms suggests that the fermion mass textures generated by this flat direction are similar to those found in ref. . We then have that the light Higgs representations consist of $`\overline{h}_1`$ and a combination of $`h_1`$ and $`h_3`$. One then finds that the leading mass terms are $`Q_1u_1^c\overline{h}_1`$ and $`Q_3d_3^ch_3`$. These mass textures are therefore not phenomenologically viable as the left–handed component of the top and bottom quarks live in different multiplets. A plausible solution is to find a flat direction for which $`h_3`$ is not part of the surviving light Higgs combination. In which case a mass term for the bottom quark can appear, for example, from the quartic term $`Q_1d_1^cH_{41}H_{21}^s`$. For the purpose of our discussion here we make the assumption that the sector $`𝐛_1`$ produces the heavy generation states and $`𝐛_{2,3}`$ produce the two light generations. A more detailed study of the phenomenology of the non–Abelian flat directions will be reported in ref. . We next turn to discuss whether any combination, and which, of the $`U(1)`$ symmetries of the FNY model remains unbroken by the choice of VEVs in Eq. (5.3). Subsequently, we will examine the phenomenological characteristics of the unbroken combinations. The first observation is that the family universal $`U(1)_Z^{}`$, which is embedded in $`SO(10)`$ is broken by this choice of VEVs. Similarly, the family universal $`U(1)`$ combination which is embedded in $`E_6`$ is broken at the string scale. As we discussed above, our general expectation is that in fact these particular $`U(1)`$ symmetries cannot remain unbroken to low energies. The set of VEVs in Eq. (5.3) leaves two $`U(1)`$ combination unbroken at the string scale. The first is given by the combination $$U(1)=3U(1)_7+U_h,$$ (5.4) while the second unbroken combination is given by $$U(1)_Z^{}=5U(1)_3^{}+U(1)_73U_h.$$ (5.5) The $`U(1)`$ generators appearing in the first combination are from the Cartan sub–algebra of the hidden $`E_8`$. Therefore, the three Standard Model generations from the sectors $`𝐛_1`$, $`𝐛_2`$ and $`𝐛_3`$ are not charged with respect to this $`U(1)`$ combination and it is consequently not of interest from the point of view of low energy experiments. In the second unbroken combination $`U(1)_3^{}`$ appears and consequently the Standard Model states are charged under this unbroken $`U(1)_Z^{}`$ symmetry. ## 6 Phenomenological characteristics As we illustrated in the previous section, the $`F`$– and $`D`$–flat solution Eq. (5.3) leaves the $`U(1)_Z^{}`$ combination, Eq. (5.5), unbroken at the string scale. Several issues are important to consider in regard to the possible low energy phenomenological implications. Furthermore, many of the issues which are crucial for fully extracting the phenomenological consequences, like the fermion identification, are still not under complete control. Consequently, prior to embarking on a detailed phenomenological analysis, we have to try to isolate those characteristics which are independent of the details about which we are ignorant at this stage. If such an extraction is possible then the discussion becomes more substantial. This is the price we have to pay for trying to extract phenomenological consequences from a theory whose natural scale is vastly separated from the experiments’ natural scale. Similarly, to this level we have found a $`U(1)`$ combination which remains unbroken at the string scale. It is quite plausible that supersymmetry breaking requires the existence of an intermediate energy scale. Of course, one can devise various scenarios, like radiative breaking, by which the $`U(1)_Z^{}`$ will be broken just at the right scale, namely near the electroweak scale. The $`U(1)_Z^{}`$ breaking can be generated due to the VEV of one of the remaining light Standard Model singlets, which are charged under $`U(1)_Z^{}`$, for example $`\mathrm{\Phi }_{56}^{}`$ and $`\overline{\mathrm{\Phi }}_{56}^{}`$. But at this stage we regard the possibility that the $`U(1)_Z^{}`$ remains unbroken down to low energies as an assumption and extract the phenomenological implications from there. We see from Table 1 that the charges of the three generations and the Higgs multiplets under the $`U(1)_Z^{}`$ are completely specified. Then, up to the caveat stated above, the phenomenological implications are completely fixed. Several observations are interesting to note. First from Eq. (5.5) we see that indeed the unbroken $`U(1)_Z^{}`$ is not of $`E_6`$ or $`SO(10)`$ origin. Moreover, the unbroken $`U(1)`$ combination does not arise from the $`U(1)`$ generators of the observable $`E_8`$, but rather from $`U(1)`$ symmetries which arise from the compactified Narain lattice. Thus, the unbroken $`U(1)`$ symmetries that we may expect to arise from string vacua are not of the GUT type. Furthermore, as the fermion charges are related to the particular type of compactification, $`U(1)_Z^{}`$ experimental data may contain information on the underlying compactified manifold. Examining then the $`U(1)_3^{}`$ charges in Table 1 we see that mass mixing of the $`Z^{}`$–gauge boson with the Standard Model $`Z`$ would not arise if the light electroweak doublets are composed only of doublets from the Neveu–Schwarz sector. This in fact would be the general case if the unbroken $`U(1)`$ is solely a combination of the Cartan generators arising from the Narain lattice, and possibly hidden sector generators. That is, if it does not contain $`U(1)`$ currents from the observable $`E_8`$. In this case, as is seen from Table 1, all Neveu–Schwarz electroweak doublets are neutral with respect to $`U(1)_Z^{}`$. $`ZZ^{}`$ mass mixing could arise if the light electroweak Higgs doublets contain a state which arises from the twisted sectors. In the case of the FNY model those are $`H_{34}`$ and $`H_{41}`$ in Table 1. In general, the states of this type, arising from twisted sectors, are charged with respect to the $`U(1)`$ currents which arise from the Narain lattice. However, here it is found that also $`H_{34}`$ and $`H_{41}`$ are neutral under the particular unbroken $`U(1)`$ combination given in Eq. (5.5). Therefore, here all the electroweak Higgs doublets are neutral under $`U(1)_Z^{}`$ and $`ZZ^{}`$ mass mixing cannot arise. Possible kinetic mixing, arising from one–loop oblique corrections to the gauge boson propagator, is also highly suppressed. This follows from our assumption that the sector $`𝐛_1`$ produces the heavy generation. The heavy generation states are therefore neutral under $`U(1)_Z^{}`$, and do not contribute to the one–loop corrections. For the two light generations, arising from the sectors $`𝐛_2`$ and $`𝐛_3`$, we see from Table 1 that the charges of the states are equal in magnitude and opposite in sign, and would therefore cancel. For Standard Model fermions the kinetic mixing can therefore only be of the order of $`(\mathrm{ln}(m_1^2m_2^2)/M_Z^{}^2)(\mathrm{ln}(m_c^2m_u^2)/M_Z^{}^2)`$, which is highly suppressed even for $`M_Z^{}500`$ GeV. Gauginos, Higgsinos and the light Higgs cannot contribute to kinetic mixing because they are all neutral under this particular $`U(1)_Z^{}`$. We now turn to the supersymmetric scalar sector. Since under our assumption the sector $`𝐛_1`$ produces the heavy generation, which is neutral under $`U(1)_Z^{}`$, only the two light generation can contribute. However, up to light fermion mass corrections, and assuming universality, the two light scalar generations are degenerate in mass. Nonuniversality, could arise due to the $`U(1)_Z^{}`$ $`D`$–term contribution. However, $`D_Z^{}`$ vanishes if the VEVs of the two fields which break $`U(1)_Z^{}`$, say $`\mathrm{\Phi }_{56}^{}`$ and $`\overline{\mathrm{\Phi }}_{56}^{}`$, are equal in magnitude. Therefore, under this assumption, the scalar contribution to the scalar masses is also negligible, and kinetic mixing is highly suppressed for this particular $`Z^{}`$ combination. We now give a rough estimate of the phenomenological constraints on $`M_Z^{}`$. For this purpose we have to normalize $`U(1)_Z^{}`$ so that it has the correct normalization to produce the correct conformal dimension, $`\overline{h}=1`$, for the massless states. ¿From Eq. (5.5) we deduce that the normalization factor is $`N=1/\sqrt{78}`$. Estimating the beta function coefficients from the charges given in Table 1, we obtain $`b_Z^{}2.4`$, where we have taken the spectrum to consist of three MSSM generations, excluding the three right–handed neutrinos. Taking $`\alpha _{\mathrm{GUT}}^125`$, and extrapolating from $`M_{\mathrm{GUT}}10^{17}`$ GeV to $`M_Z`$, we obtain $`\alpha _Z^{}^1(M_Z)40`$. As seen from Table 1 the charges of the two light generations, while equal in magnitude are opposite in charge, and consequently not universal. Therefore, the strongest constraint is from Flavor Changing Neutral Currents, which arises here from fermion mixing. To estimate this constraint we use $`\mathrm{\Gamma }(K_L^0\mu ^+\mu ^{})10^8\mathrm{\Gamma }(K^+\mu ^+\nu _\mu ).`$ (6.1) Estimating the tree diagrams we obtain $`Q_Z^{}^4\alpha _Z^{}^2\mathrm{cos}^2\theta _C\mathrm{sin}^2\theta _C/M_Z^{}^4=10^8\alpha _2^2\mathrm{sin}^2\theta _C/(4M_W^4).`$ (6.2) With the appropriately normalized charges for $`Q_Z^{}`$, we obtain $`M_Z^{}25M_W2\mathrm{TeV}.`$ We remark that the additional suppression of the $`Z^{}`$ interaction is obtained because of the $`U(1)_Z^{}`$ normalization factor that we calculated above. This reflects the fact that the $`U(1)_Z^{}`$ combination contains Cartan generators of the hidden $`E_8`$ under which the Standard Model states are not charged. The consequence is that there is roughly an order of magnitude suppression of the $`Q_Z^{}`$ charges of the Standard Model states. A more stringent constraint arises by considering the mixing in the $`K_0\overline{K}_0`$ system parametrized by the mass difference $`\mathrm{\Delta }M_K=3.5\times 10^{12}\mathrm{MeV}`$ . Treating the $`Z^{}`$ as a contact interaction we have that $`\mathrm{\Delta }M_KG_2M_Kf_K^2`$, where $`f_K1.2m_\pi `$ is the kaon decay constant, $`M_K0.5`$ GeV is the kaon mass, and $`G_2=(Q_Z^{}^2/M_Z^{}^2)4\pi \alpha _Z^{}(\mathrm{cos}\theta _C\mathrm{sin}\theta _C)^2`$ is the $`Z^{}`$ contact interaction term. We then find that $`G_210^7G_F`$, where $`G_F`$ is the Fermi constant. From this we obtain $`M_Z^{}>30`$ TeV. Considering the corresponding mass difference in the $`B`$–meson system, $`\mathrm{\Delta }M_B3.1210^4\mathrm{eV}`$ , imposes only $`m_Z^{}500\mathrm{G}\mathrm{e}\mathrm{V}`$, and is therefore less restrictive, where we have used the Standard Model value for $`V_{td}`$. We do not estimate here constraints arising from FCNC in the lepton sector as the leptonic mixing parameters are not known. It is therefore expected that a $`Z^{}`$ with non–universal charges for the two light generations is constrained to be above the reach of the LHC. Nevertheless, as we discuss below, a $`Z^{}`$ with universal couplings for the first two light generations and with different couplings to the heavy generation may also arise from the free fermionic models and may in fact be a signature of the $`Z_2\times Z_2`$ orbifold which underlies the free fermionic models. We note that if a $`Z^{}`$ gauge boson of the type that we discussed above is in the region accessible to future hadron colliders, it will yield spectacular signatures. Namely, in the case of the particular $`U(1)_Z^{}`$ combination that we examined here, it will result in enhancement in the production of the two light generations, whereas a parallel enhancement in the production of the heavy generation will not be observed. Similarly, for this particular $`U(1)_Z^{}`$ combination, production of Higgs doublets and gauginos in the $`Z^{}`$ channel will not be observed. While it is not our aim to argue that the particular $`U(1)_Z^{}`$ combination examined here is necessarily phenomenologically viable, what we see is that in a given string model, and for a given flat direction, the phenomenological consequences and possible production and decay channels are completely specified and yield distinctive signatures. ## 7 Discussion We emphasize that it is not our intent to argue here that the particular $`U(1)_Z^{}`$ combination that we examined is necessarily “the” phenomenologically viable combination that may be seen in future collider experiments. What we have shown is that in a specific string model the $`U(1)_Z^{}`$ combinations which remains unbroken for specific flat directions are given. Consequently, the charges of the Standard Model fermions are specified and, in the case that the $`U(1)_Z^{}`$ symmetry remains unbroken down to low energies, the phenomenological implications are determined. The $`U(1)_Z^{}`$ that we examined here provides an illustrative example. However, we believe that more general lessons can be extracted. The first is that we anticipate that the $`U(1)`$ combinations which remain unbroken after analysis of the flat directions are not of the type which appear in $`SO(10)`$ or $`E_6`$ grand unifying theories. Therefore, it is anticipated that if a $`U(1)`$ combination remains unbroken down to low energies, it contains $`U(1)`$ factors which are external to the GUT gauge group. The second important lesson arises by examining the various $`U(1)`$ charges given in Table 1. We see that a common feature is precisely the flavor non–universality of the different $`U(1)`$ combinations. Thus, we see, for example that for $`U_1^{}`$, $`U_4^{}`$, and $`U_4`$ the charges of the two light generations are universal and differ from the charges of the heavy generation. Flat directions which preserve one of these $`U(1)`$’s as a component of an unbroken $`U(1)`$ symmetry, may therefore yield a $`Z^{}`$ gauge boson which is less severely restricted by FCNC constraints. Nevertheless, the distinctive collider signatures of a $`Z^{}`$ arising from any of those $`U(1)`$ symmetries will be a non–universality in the production of the different generations. Thus, for example, $`U_4^{}`$ would predict enhancement in the production of the two light families, without a corresponding enhancement in production of the heavy family, whereas $`U_4`$ would predict exactly the opposite. A $`Z^{}`$ with this characteristic may in fact be a consequence of the $`Z_2\times Z_2`$ orbifold with standard embedding, which underlies the free fermionic formulation for the following reason. Take, for example, the case in which the anomalous $`U(1)`$ is a combination which is embedded in $`E_6`$ and is given by $`U_A=U_1+U_2+U_3`$ in the notation of Section 3. The two anomaly free orthogonal combinations can be taken as $`U_1^{}=U_1U_2`$ and $`U_2^{}=U_1+U_22U_3`$. The states of each generation from each sector $`b_j`$ have charge $`+1/2`$ under $`U(1)_j`$ and are neutral with respect to the other two. Consequently, $`U_1^{}`$ produces charges which are equal in magnitude and opposite in sign for two generation, whereas one generation is neutral under it. This yields the same type of $`Z^{}`$ that we examined here and is strongly constrained by FCNC. On the other hand $`U_2^{}`$ is universal with respect to two families and produces different charges for the third family. This situation may, in fact, be a unique consequence of the $`Z_2\times Z_2`$ orbifold twisting, due to its cyclic permutation symmetry. From Table 1 we see that, in fact, this type of charge assignment is also frequently preserved in the three generation models. What we argue is that if a $`Z^{}`$ with universal couplings for the two light generations and different couplings for the heavy generation is observed in future experiments, it may be a key piece of evidence for the $`Z_2\times Z_2`$ orbifold compactification. In the case of a $`Z^{}`$ with charges equal in magnitude but opposite in sign for the first two generations, we may expect it to be outside the reach of forthcoming hadron colliders. However, if it is not too far above their reach, we may expect novel FCNC phenomena, and potentially new sources for CP violation. We note that an additional $`Z^{}`$ of the type that we discussed here has also been advocated as playing a role in suppressing proton decay in supersymmetric extensions of the Standard Model . We also remark that very recently it has been suggested that there exists evidence for a $`Z^{}`$ with these characteristics in electroweak precision data . All in all, nature may eventually prove to be kind for her patient and obedient servants. ## 8 Acknowledgments AF would like to heartily thank Misha Voloshin for useful discussions on the matters of this paper and physics in general. This work is supported in part by DOE Grants No. DE–FG–0294ER40823 (AEF, TtV) and DE–FG–0395ER40917 (GC,DVN). ## Appendix A Quantum Number of FNY Massless Fields \| state \| $`U_E`$ \| $`(C,L)_Y`$ \| $`U_A`$ \| $`U_C`$ \| $`U_L`$ \| $`U_1^{^{}}`$ \| $`U_2^{^{}}`$ \| $`U_3^{^{}}`$ \| $`U_4^{^{}}`$ \| $`U_5^{^{}}`$ \| $`U_4`$ \| $`(3,2,2^{^{}})_H`$ \| $`U_7`$ \| $`U_H`$ \| $`U_9`$ \| \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| --- \| \| $`Q_1`$ \| $`\frac{2,1}{3}`$ \| $`(3,2)_{\frac{1}{6}}`$ \| 8 \| 2 \| 0 \| -4 \| 2 \| 0 \| 0 \| -24 \| 2 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`Q_2`$ \| $`\frac{2,1}{3}`$ \| $`(3,2)_{\frac{1}{6}}`$ \| 12 \| 2 \| 0 \| 2 \| -10 \| 2 \| 2 \| 20 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`Q_3`$ \| $`\frac{2,1}{3}`$ \| $`(3,2)_{\frac{1}{6}}`$ \| 8 \| 2 \| 0 \| 2 \| 14 \| -2 \| 2 \| 32 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`d_1^c`$ \| $`\frac{1}{3}`$ \| $`(\overline{3},1)_{\frac{1}{3}}`$ \| 8 \| -2 \| 4 \| -4 \| 2 \| 0 \| 0 \| -24 \| -2 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`d_2^c`$ \| $`\frac{1}{3}`$ \| $`(\overline{3},1)_{\frac{1}{3}}`$ \| 8 \| -2 \| 4 \| 2 \| -10 \| -2 \| -2 \| -80 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`d_3^c`$ \| $`\frac{1}{3}`$ \| $`(\overline{3},1)_{\frac{1}{3}}`$ \| 4 \| -2 \| 4 \| 2 \| 14 \| 2 \| -2 \| -68 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`u_1^c`$ \| $`\frac{2}{3}`$ \| $`(\overline{3},1)_{\frac{2}{3}}`$ \| 8 \| -2 \| -4 \| -4 \| 2 \| 0 \| 0 \| -24 \| -2 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`u_2^c`$ \| $`\frac{2}{3}`$ \| $`(\overline{3},1)_{\frac{2}{3}}`$ \| 12 \| -2 \| -4 \| 2 \| -10 \| 2 \| 2 \| 20 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`u_3^c`$ \| $`\frac{2}{3}`$ \| $`(\overline{3},1)_{\frac{2}{3}}`$ \| 8 \| -2 \| -4 \| 2 \| 14 \| -2 \| 2 \| 32 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`e_1^c`$ \| 1 \| $`(1,1)_1`$ \| 8 \| 6 \| 4 \| -4 \| 2 \| 0 \| 0 \| -24 \| -2 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`e_2^c`$ \| 1 \| $`(1,1)_1`$ \| 12 \| 6 \| 4 \| 2 \| -10 \| 2 \| 2 \| 20 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`e_3^c`$ \| 1 \| $`(1,1)_1`$ \| 8 \| 6 \| 4 \| 2 \| 14 \| -2 \| 2 \| 32 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`N_1^c`$ \| 0 \| $`(1,1)_0`$ \| 8 \| 6 \| -4 \| -4 \| 2 \| 0 \| 0 \| -24 \| -2 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`N_2^c`$ \| 0 \| $`(1,1)_0`$ \| 8 \| 6 \| -4 \| 2 \| -10 \| -2 \| -2 \| -80 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`N_3^c`$ \| 0 \| $`(1,1)_0`$ \| 4 \| 6 \| -4 \| 2 \| 14 \| 2 \| -2 \| -68 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`L_1`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| 8 \| -6 \| 0 \| -4 \| 2 \| 0 \| 0 \| -24 \| 2 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`L_2`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| 8 \| -6 \| 0 \| 2 \| -10 \| -2 \| -2 \| -80 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`L_3`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| 4 \| -6 \| 0 \| 2 \| 14 \| 2 \| -2 \| -68 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`h_1`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| 16 \| 0 \| -4 \| -8 \| 4 \| 0 \| 0 \| -48 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`h_2`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| -20 \| 0 \| -4 \| -4 \| 20 \| 0 \| 0 \| 60 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`h_3`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| -12 \| 0 \| -4 \| -4 \| -28 \| 0 \| 0 \| 36 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`\overline{h}_1`$ \| 1, 0 \| $`(1,2)_{\frac{1}{2}}`$ \| -16 \| 0 \| 4 \| 8 \| -4 \| 0 \| 0 \| 48 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`\overline{h}_2`$ \| 1, 0 \| $`(1,2)_{\frac{1}{2}}`$ \| 20 \| 0 \| 4 \| 4 \| -20 \| 0 \| 0 \| -60 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`\overline{h}_3`$ \| 1, 0 \| $`(1,2)_{\frac{1}{2}}`$ \| 12 \| 0 \| 4 \| 4 \| 28 \| 0 \| 0 \| -36 \| 0 \| (1,1,1) \| 0 \| 0 \| 0 \| \| $`H_{34}`$ \| 1,0 \| $`(1,2)_{\frac{1}{2}}`$ \| 8 \| 3 \| 2 \| -2 \| -11 \| 2 \| -4 \| 32 \| 0 \| (1,1,1) \| -1 \| 3 \| 0 \| \| $`H_{41}`$ \| 0,-1 \| $`(1,2)_{\frac{1}{2}}`$ \| 0 \| -3 \| -2 \| 2 \| -13 \| -2 \| -4 \| 56 \| 0 \| (1,1,1) \| 1 \| -3 \| 0 \| Table 1: Gauge Charges of FNY three generation and Higgs sectors. The names of the states appear in the first column, with the states’ various charges appearing in the other columns. The entries under $`(C,L)_Y`$ denote Standard Model charges, while the entries under $`(3,2,2^{})`$ denote hidden sector $`SU(3)_H\times SU(2)_H\times SU(2)_H^{^{}}`$ charges. All $`U(1)`$ charges are multiplied by a factor of 4 relative to the definition in Eqs. (5.1) and (5.2).
warning/0002/gr-qc0002060.html	ar5iv	text	# References I. Introduction It is well known that Einstein’s general relativity can be obtained from several distinct Lagrangian formulations. One suitable formulation is the teleparallel equivalent of general relativity (TEGR), which is defined in terms of tetrad fields $`e_\mu ^a`$ ($`a,\mu `$ are SO(3,1) and space-time indices, respectively) and actually represents an alternative geometrical framework for Einstein’s equations. The Lagrangian density for the tetrad field in the TEGR is given by a sum of quadratic terms in the torsion tensor $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a`$, which is related to the anti-symmetric part of the connection $`\mathrm{\Gamma }_{\mu \nu }^\lambda =e^{a\lambda }_\mu e_{a\nu }`$. The curvature tensor constructed out of the latter vanishes identically. This connection defines a space with teleparallelism, or absolute parallelism. In a space-time with an underlying tetrad field two vectors at distant points may be called parallel if they have identical components with respect to the local tetrads at the points considered. Thus consider a vector field $`V^\mu (x)`$. At the point $`x^\lambda `$ its tetrad components are given by $`V^a(x)=e_\mu ^a(x)V^\mu (x)`$. For the tetrad components $`V^a(x+dx)`$ it is easy to show that $`V^a(x+dx)=V^a(x)+DV^a(x)`$, where $`DV^a(x)=e_\mu ^a(_\lambda V^\mu )dx^\lambda `$. The covariant derivative $``$ is constructed out of the connection $`\mathrm{\Gamma }_{\mu \nu }^\lambda =e^{a\lambda }_\mu e_{a\nu }`$. Therefore such connection defines a condition for absolute parallelism in space-time. The tetrad fields are thus required to transform under the global SO(3,1) group. The Lagrangian density for the TEGR is based on the relation $$eR(e)=e(\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a)+\mathrm{\hspace{0.17em}2}_\mu (eT^\mu ),$$ $`(1)`$ which can be verified by substituting the Levi-Civita connection $`{}_{}{}^{0}\omega _{\mu ab}^{}`$ into the scalar curvature $`R(e)`$ on the left hand side of (1) by means of the relation $`{}_{}{}^{0}\omega _{\mu ab}^{}=K_{\mu ab}`$, where $`K_{\mu ab}`$ is the contorsion tensor: $`K_{\mu ab}=\frac{1}{2}e_a^\lambda e_b^\nu (T_{\lambda \mu \nu }+T_{\nu \lambda \mu }T_{\mu \nu \lambda }`$). In empty space-time the Lagrangian density for the TEGR is given by $$L(e)=ke(\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a),$$ $`(2)`$ where $`k=\frac{1}{16\pi G}`$, $`e=det(e_\mu ^a)`$ and $`T_a=T_{ba}^b`$. As usual, tetrad fields convert space-time into SO(3,1) indices and vice-versa. Let $`\frac{\delta L}{\delta e_\mu ^a}`$ denote the field equation satisfied by $`e_\mu ^a`$. It can be verified by explicit calculations that the latter are equivalent to Einstein’s equations in tetrad form: $$\frac{\delta L}{\delta e^{a\mu }}=\frac{1}{2}e\{R_{a\mu }(e)\frac{1}{2}e_{a\mu }R(e)\}.$$ $`(3)`$ This theory has been considered long ago by Møller, although not precisely in the form presented above. More recently it has been reconsidered as a gauge theory of the translation group. Interesting developments have been achieved in the context of Ashtekar variables. It has been shown that in the teleparallel geometry the (complex) Hamiltonian becomes quadratic in the new field momenta, and that Einstein’s equations become formally of the Yang-Mills type. An important feature of the above formulation is that the tetrad fields $`e_\mu ^a`$ transform under the global SO(3,1) group. In fact the TEGR was previously considered in ref. with a local SO(3,1) symmetry. In order to make clear this point let us recall that relation (1) can be written in terms of $`e_\mu ^a`$ and an arbitrary spin connection $`\omega _{\mu ab}`$. As discussed in , $`e_\mu ^a`$ and $`\omega _{\mu ab}`$ are not related to each other via field equations. This arbitrary connection can be identically written as $`\omega _{\mu ab}=^0\omega _{\mu ab}+K_{\mu ab}`$, where $`K_{\mu ab}`$ is the same as above but now the torsion tensor is given by $`T_{\mu \nu }^a=_\mu e_\nu ^a_\nu e_\mu ^a+\omega _\mu ^a{}_{b}{}^{}e_{\nu }^{b}\omega _\nu ^a{}_{b}{}^{}e_{\mu }^{b}`$. Substituting $`\omega _{\mu ab}`$ into the scalar curvature $`R(e,\omega )=e^{a\mu }e^{b\nu }R_{ab\mu \nu }(\omega )`$ we obtain the identity $$eR(e,\omega )=eR(e)+e(\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a)2_\mu (eT^\mu ).$$ The teleparallel space is determined by the vanishing of $`R(e,\omega )`$. It was shown in ref. that the Hamiltonian formulation of the TEGR with a local SO(3,1) symmetry cannot be made consistent since the constraint algebra does not “close”, and therefore the dynamical evolution of the field quantities is not well defined. A well established Hamiltonian formulation can only be achieved if the SO(3,1) is turned into a global symmetry group. The requirement of the vanishing of the curvature tensor $`R_{b\mu \nu }^a(\omega )`$ that appears on the left hand side of the identity above has the ultimate effect of discarding the connection $`\omega _{\mu ab}`$. A global SO(3,1) symmetry leads to a theory with well defined initial value problem. Therefore we can dispense with the local symmetry of the theory together with the constraint of vanishing curvature, both considered previously in . Kopczyński (ref., section 4) has argued that even at the Lagrangian level one can require $`\omega _{\mu ab}=0`$ and obtain the same field equations. As a matter of fact the present results confirm this point of view. Thus the Hamiltonian will display a global SO(3) symmetry. We remark that the local SO(3) symmetry group has recently played a special role in connection with the Hamiltonian formulation of gravity theories, as developed in . In the latter it is carried out a Poincaré invariant foliation of the space-time in which the SO(3) group is taken as the classification subgroup of the Poincaré group (rather than the Lorentz group). A Hamiltonian formalism based on a nonlinear realization of the Poincaré group is constructed and applied to the Einstein-Cartan theory. The major motivation for considering the TEGR resides in the fact that it is possible to make definite statements about the energy and momentum of the gravitational field. In the 3+1 formulation of the TEGR we find that the Hamiltonian and vector constraints contain each one a divergence in the form of a scalar and vector densities, respectively, that can be identified as the energy and momentum densities of the gravitational field. Therefore the Hamiltonian and vector constraints are considered as energy-momentum equations. This identification has proven to be consistent, and has shown that the TEGR provides a natural setting for investigations of the gravitational energy. Several relevant applications have been presented in the literature. Among the latter we point out investigations on the gravitational energy of rotating black holes (the evaluation of the irreducible mass of a Kerr black hole) and of Bondi’s radiating metric. In this paper we carry out the Hamiltonian formulation of an arbitrary teleparallel theory, quadratic in the torsion tensor just like in (2). The Lagrangian density to be considered describes a three-parameter family of teleparallel theories. We want to investigate the existence of theories that satisfy the only criterium of having a well defined Hamiltonian formulation, which amounts to having a well posed initial value problem. For this purpose we adopt the field quantity definitions of Hayashi and Shirafuji. We also make use of their analysis of the Newtonian limit of these theories and of the restrictions implied by this requirement. The investigation will be carried out along the lines of ref. . As in the latter, we will impose the time gauge condition to the tetrad field (this condition is important in order to establish a comparison with ). A consistent implementation of the Legendre transform reduces the three-parameter to a one-parameter family of theories. The latter constitutes a well defined theory with only first class constraints. The free parameter is fixed by requiring the gravitational field to exhibit the Newtonian limit. The resulting theory is just the TEGR. In section II we establish the definitions and present the Lagrangian formulation of the teleparallel theory. The Hamiltonian formulation is established in section III. The relevant details of the Legendre transform will be presented in this section. In the last section we present our final comments. Notation: spacetime indices $`\mu ,\nu ,\mathrm{}`$ and SO(3,1) indices $`a,b,\mathrm{}`$ run from 0 to 3. In the 3+1 decomposition latin indices from the middle of the alphabet indicate space indices according to $`\mu =0,i,a=(0),(i)`$. The flat spacetime metric is fixed by $`\eta _{(0)(0)}=1`$. II. The Lagrangian formulation of an arbitrary teleparallel theory We begin by presenting the four basic postulates that the Lagrangian density for the gravitational field in empty space-time, in the teleparallel geometry, must satisfy. It must be (i) invariant under general coordinate transformations, (ii) invariant under global SO(3,1) transformations, (iii) invariant under parity transformations and (iv) quadratic in the torsion tensor. The most general Lagrangian density can be written as $$L_0=ke(c_1t^{abc}t_{abc}+c_2v^av_a+c_3a^ba_b),$$ $`(4)`$ where $`c_1,c_2,c_3`$ are constants and $$t_{abc}=\frac{1}{2}(T_{abc}+T_{bac})+\frac{1}{6}(\eta _{ac}v_b+\eta _{bc}v_a)\frac{1}{3}\eta _{ab}v_c,$$ $`(5.1)`$ $$v_a=T_{ba}^b=T_a,$$ $`(5.2)`$ $$a_a=\frac{1}{6}\epsilon _{abcd}T^{bcd},$$ $`(5.3)`$ $$T_{abc}=e_b^\mu e_c^\nu T_{a\mu \nu }.$$ Definitions (5) correspond to the irreducible components of the torsion tensor. As we mentioned earlier, we are departing from Hayashi and Shirafuji’s notation. Our analysis will make contact both with Ref. and with Ref. . We are considering an extended teleparallel theory in the sense of Müller-Hoissen and Nitsch. In order to carry out the Hamiltonian formulation in the next section we need to rewrite the three terms of $`L_0`$ in order to make explict the appearance of the torsion tensor. Therefore we rewrite $`L_0`$ as $$L_0=ke(c_1X^{abc}T_{abc}+c_2Y^{abc}T_{abc}+c_3Z^{abc}T_{abc}),$$ $`(6)`$ with the following definitions: $$X^{abc}=\frac{1}{2}T^{abc}+\frac{1}{4}T^{bac}\frac{1}{4}T^{cab}+\frac{1}{4}(\eta ^{ac}v^b\eta ^{ab}v^c),$$ $`(7.1)`$ $$Y^{abc}=\frac{1}{2}(\eta ^{ab}v^c\eta ^{ac}v^b),$$ $`(7.2)`$ $$Z^{abc}=\frac{1}{18}(T^{abc}+T^{bca}+T^{cab}).$$ $`(7.3)`$ The definitions above satisfy $`X^{abc}=X^{acb}`$, $`Y^{abc}=Y^{acb}`$ and $`Z^{abc}=Z^{acb}`$. $`X^{abc},Y^{abc}`$ and $`Z^{abc}`$ have altogether the same number of independent components of $`T^{abc}`$. It is not difficult to verify that $`X^{abc}+X^{bca}+X^{cab}0`$. Let us define the field quantity $`\mathrm{\Sigma }^{abc}`$ by $$\mathrm{\Sigma }^{abc}=c_1X^{abc}+c_2Y^{abc}+c_3Z^{abc},$$ $`(8)`$ which allow us to further rewrite $`L_0`$ according to the notation of Ref. : $$L_0=ke\mathrm{\Sigma }^{abc}T_{abc}.$$ $`(9)`$ We note that if the constants $`c_i`$ satisfy $$c_1=\frac{2}{3},c_2=\frac{2}{3},c_3=\frac{3}{2},$$ $`(10)`$ then $`\mathrm{\Sigma }^{abc}`$ reduces to the corresponding quantity of the TEGR : $$\mathrm{\Sigma }_{_{_{TEGR}}}^{abc}=\frac{1}{4}(T^{abc}+T^{bac}T^{cab})+\frac{1}{2}(\eta ^{ac}v^b\eta ^{ab}v^c),$$ $`(11)`$ for which we have $$\mathrm{\Sigma }_{_{_{TEGR}}}^{abc}T_{abc}=\frac{1}{4}T^{abc}T_{abc}+\frac{1}{2}T^{abc}T_{bac}T^aT_a.$$ $`(12)`$ In order to carry out the 3+1 decomposition of the theory we need a first order differential Lagrangian density. It will be achieved through the introduction of an auxiliary field quantity $`\mathrm{\Delta }_{abc}`$, according to the procedure developed in . Thus we consider the Lagrangian density $$L(e,\mathrm{\Delta })=ke(c_1\mathrm{\Theta }^{abc}+c_2\mathrm{\Omega }^{abc}+c_3\mathrm{\Gamma }^{abc})(\mathrm{\Delta }_{abc}2T_{abc}),$$ $`(13)`$ where $`\mathrm{\Theta }^{abc}`$, $`\mathrm{\Omega }^{abc}`$ and $`\mathrm{\Gamma }^{abc}`$ are defined in similarity with $`X^{abc}`$, $`Y^{abc}`$ and $`Z^{abc}`$, respectively: $$\mathrm{\Theta }^{abc}=\frac{1}{2}\mathrm{\Delta }^{abc}+\frac{1}{4}\mathrm{\Delta }^{bac}\frac{1}{4}\mathrm{\Delta }^{cab}+\frac{1}{4}(\eta ^{ac}\mathrm{\Delta }^b\eta ^{ab}\mathrm{\Delta }^c),$$ $`(14.1)`$ $$\mathrm{\Omega }^{abc}=\frac{1}{2}(\eta ^{ab}\mathrm{\Delta }^c\eta ^{ac}\mathrm{\Delta }^b),$$ $`(14.2)`$ $$\mathrm{\Gamma }^{abc}=\frac{1}{18}(\mathrm{\Delta }^{abc}+\mathrm{\Delta }^{bca}+\mathrm{\Delta }^{cab}).$$ $`(14.3)`$ The three quantities above are anti-symmetric in the last two indices. The field equations are most easily obtained by making use of the three identities satisfied by these expressions: $$X^{abc}\mathrm{\Delta }_{abc}=\mathrm{\Theta }^{abc}T_{abc},$$ $`(15.1)`$ $$Y^{abc}\mathrm{\Delta }_{abc}=\mathrm{\Omega }^{abc}T_{abc},$$ $`(15.2)`$ $$Z^{abc}\mathrm{\Delta }_{abc}=\mathrm{\Gamma }^{abc}T_{abc}.$$ $`(15.3)`$ These identities turn out to be useful in the variation of the action integral. We note in addition that since $`\mathrm{\Theta }^{abc}\mathrm{\Delta }_{abc}`$ is quadratic in $`\mathrm{\Delta }_{abc}`$ it follows that $$\delta (\mathrm{\Theta }^{abc}\mathrm{\Delta }_{abc})=\mathrm{\hspace{0.33em}2}\mathrm{\Theta }^{abc}\delta (\mathrm{\Delta }_{abc}),$$ $`(16)`$ and similarly for $`\mathrm{\Omega }^{abc}`$ and $`\mathrm{\Gamma }^{abc}`$ (this result can be verified by explicit calculations). Because of (16) the variation of $`kec_1\mathrm{\Theta }^{abc}(\mathrm{\Delta }_{abc}2T_{abc})`$ with respect to $`\mathrm{\Delta }_{abc}`$ is given by $$\delta \{kec_1\mathrm{\Theta }^{abc}(\mathrm{\Delta }_{abc}2T_{abc})\}=2kec_1(\mathrm{\Theta }^{abc}X^{abc})\delta (\mathrm{\Delta }_{abc}),$$ $`(17)`$ and likewise for the other terms in $`L`$. Therefore the field equations arising from (13) with respect to variations of $`\mathrm{\Delta }_{abc}`$ are given by $$c_1(\mathrm{\Theta }^{abc}X^{abc})+c_2(\mathrm{\Omega }^{abc}Y^{abc})+c_3(\mathrm{\Gamma }^{abc}Z^{abc})=\mathrm{\hspace{0.33em}0}.$$ $`(18)`$ The only solution of (18) for arbitrary constants $`c_i`$ is given by $$\mathrm{\Delta }_{abc}=T_{abc}=e_b^\mu e_c^\nu T_{a\mu \nu }.$$ $`(19)`$ Note that (18) represents 24 equations for 24 unknown quantities $`\mathrm{\Delta }_{abc}`$. If the constants $`c_i`$ satisfy (10) then the field equations obtained with respect to variations of $`e^{a\mu }`$ are, in view of (3), precisely equivalent to Einstein’s equations. It should me mentioned that teleparallel theories also arise as effective theories in the context of Poincaré gauge theories of gravity by means of a modified double duality ansatz (see, for instance, Baekler et. al.), however in the limit of vanishing curvature of the Riemann-Cartan manifold. III. The Hamiltonian formulation Although in this section we still maintain the notation given at the end of section I, we will make a change of notation regarding the tetrad field $`e_\mu ^a`$. The space-time tetrad field considered in the last section will be denoted here as $`{}_{}{}^{4}e_{\mu }^{a}`$, to emphasize that it is the tetrad field of the four-dimensional space-time. In a 3+1 decomposition the space-time tetrad field does not coincide with the tetrad field restricted (projected) to the three-dimensional spacelike hypersurface. We adopt the standard 3+1 decomposition for the tetrad field: $${}_{}{}^{4}e_{k}^{a}=e_k^a,$$ $${}_{}{}^{4}e_{}^{ai}=e^{ai}+\frac{N^i}{N}\eta ^a,$$ $$e^{ai}=\overline{g}^{ik}e_k^a,\eta ^a=N^4e^{a0},$$ $${}_{}{}^{4}e_{0}^{a}=N^ie_i^a+N\eta ^a,$$ $${}_{}{}^{4}e=Ne=N\sqrt{e_i^ae_{aj}},$$ $`(20)`$ where $`g_{ij}=e_i^ae_{aj}`$ and $`\overline{g}^{ij}g_{jk}=\delta _k^i`$. The vector $`\eta ^a`$ satisfies $$\eta _ae_k^a=0,\eta _a\eta ^a=1.$$ It follows that $$e^{bk}e_{bj}=\delta _j^k,$$ $$e_i^ae^{bi}=\eta ^{ab}+\eta ^a\eta ^b.$$ The components $`e^{ai}`$ and $`e_k^a`$ are now restricted to the spacelike hypersurface. The Hamiltonian formulation will be established by rewriting the Lagrangian density (13) in the form $`L=p\dot{q}H`$. There is no time derivative of $`{}_{}{}^{4}e_{a0}^{}`$, and therefore we will enforce the corresponding momentum $`P^{a0}`$ to vanish from the outset. In analogy with (8) let us define the quantity $`\mathrm{\Lambda }^{abc}`$, $$\mathrm{\Lambda }^{abc}=c_1\mathrm{\Theta }^{abc}+c_2\mathrm{\Omega }^{abc}+c_3\mathrm{\Gamma }^{abc},$$ $`(21)`$ in terms of which we define $`P^{ai}`$, the momentum canonically conjugated to $`e_{ai}`$: $$P^{ai}=\mathrm{\hspace{0.33em}4}k^4e\mathrm{\Lambda }^{a0i}=\mathrm{\hspace{0.33em}4}kee_b^i\eta _c\mathrm{\Lambda }^{abc}.$$ $`(22)`$ In a first step the Lagrangian density is written as $$L=P^{ai}\dot{e}_{ai}+^4e_{a0}_iP^{ai}+2Nke\mathrm{\Lambda }^{aij}T_{aij}+N^kP^{ai}T_{aik}$$ $$Nke\mathrm{\Lambda }^{abc}\mathrm{\Delta }_{abc}_i(P^{ai}{}_{}{}^{4}e_{a0}^{}),$$ $`(23)`$ where $`\mathrm{\Lambda }^{aij}=e_b^ie_c^j\mathrm{\Lambda }^{abc}`$. The task of writing $`L`$ in terms of $`e_{ai}`$, $`P^{ai}`$ and Lagrange multipliers is not trivial. The troublesome term is $`Nke\mathrm{\Lambda }^{abc}\mathrm{\Delta }_{abc}`$. We will make use of the field equations (19) and identify $$\mathrm{\Delta }_{a\mu \nu }=T_{a\mu \nu }$$ in $`L`$. The Legendre transform would be straightforward if, in view of (22), $`\mathrm{\Lambda }^{aij}`$ would depend only on $`e_{(i)j}`$ and its spatial derivatives, which at this point is not the case. The Hamiltonian density cannot depend on the components $`\mathrm{\Delta }_{a0j}=T_{a0j}`$, associated to the velocities $`\dot{e}_{ai}`$. Therefore these components will have to be eliminated in the Legendre transform. We find it convenient to establish a decomposition for $`\mathrm{\Lambda }^{abc}`$ in order to distinguish the components that contribute to the canonical momentum $`P^{ai}`$. It is given by $$\mathrm{\Lambda }^{abc}=\frac{1}{4ke}(\eta ^be_i^cP^{ai}\eta ^ce_i^bP^{ai})+e_i^be_j^c\mathrm{\Lambda }^{aij}.$$ $`(24)`$ The quantity $`\mathrm{\Lambda }^{aij}`$ in the expression above contains “velocity” terms $`\mathrm{\Delta }_{a0j}`$ that cannot be inverted and written in terms of $`P^{ai}`$. However these terms will not be present in final expression of $`L`$. This feature will be achieved in view of the Schwinger’s time gauge condition $$\eta ^a=\delta _{(0)}^a,$$ $`(25)`$ that implies $`{}_{}{}^{4}e_{(k)}^{0}=e_i^{(0)}=0`$. (25) is assumed to hold from now on, i.e., it is assumed to hold before varying the action. As a consequence $`\dot{e}_{(0)i}=0`$. Taking into account definitions (21) and (22) we find by explict calculations that $$P^{(0)k}=2ke(c_1+c_2)T_{(0)}^{(0)}{}_{}{}^{k}+ke(c_12c_2)T^k.$$ $`(26)`$ We will soon return to this expression. The time gauge condition actually reduces the configuration space of the theory, and also reduces the symmetry group from the SO(3,1) to the global SO(3) group. As a consequence the teleparallel geometry is restricted to the three-dimensional spacelike hypersurface. In the following we will rewrite the various components of $`L`$ in terms of canonical quantities. First, it is not difficult to verify that $${}_{}{}^{4}e_{a0}^{}_iP^{ai}=N^ke_{(j)k}_iP^{(j)i}N_iP^{(0)i}.$$ $`(27)`$ We also have $$_i(^4e_{a0}P^{ai})=_i(NP^{(0)i})_i(N_kP^{ki}).$$ $`(28)`$ Next we consider $`Nke\mathrm{\Lambda }^{abc}\mathrm{\Delta }_{abc}`$. By making use of (20), (24) and (25) it can be rewritten as $$Nke\mathrm{\Lambda }^{abc}\mathrm{\Delta }_{abc}=N(\frac{c_1}{4}+\frac{c_3}{18})^1\{\frac{1}{16ke}(P^{ij}P_{ji}P^{(0)i}P_i^{(0)})$$ $$+\frac{1}{2}e_{(m)i}P_j^{(m)}\mathrm{\Lambda }^{(0)ij}kee_i^{(m)}e_{(n)}^k\mathrm{\Lambda }^{(n)ij}\mathrm{\Lambda }_{(m)kj}\}$$ $$+N\left(\frac{c_1}{2c_2}1\right)\{\frac{1}{48ke}(P^2P^{(0)i}P_i^{(0)})+\frac{1}{6}P_k^{(0)}e_{(m)j}\mathrm{\Lambda }^{(m)jk}\frac{1}{3}kee_{(m)i}\mathrm{\Lambda }^{(m)ik}e_{(n)j}\mathrm{\Lambda }_k^{(n)j}\}$$ $$+N\left(\frac{c_3}{9}\frac{c_1}{4}\right)\{\frac{1}{4}\mathrm{\Delta }_{ij(0)}P^{ij}ke\mathrm{\Delta }_{i(0)j}\mathrm{\Lambda }^{(0)ij}ke\mathrm{\Delta }_{ikj}\mathrm{\Lambda }^{kij}\frac{1}{4}\mathrm{\Delta }_{(0)(0)i}P^{(0)i}\frac{1}{4}\mathrm{\Delta }_{(0)ij}P^{ij}\}.$$ $`(29)`$ Spatial indices are raised and lowered with the help of $`e_{(i)j}`$ and $`e^{(k)m}`$. We note that the first term on the second line of the expression above actually reads (except for the lapse function and for the multiplicative term) $$\frac{1}{2}P_{[ij]}\mathrm{\Lambda }^{(0)ij},$$ where $`[..]`$ denotes antisymmetrization. The expression of $`\mathrm{\Lambda }^{(0)ij}`$ contains “velocity” terms $`\mathrm{\Delta }_{a0j}`$. We know, however, that in tetrad type theories of gravity the anti-symmetric part of the momentum is contrained to vanish. Let us evaluate the expression of $`P_{[ij]}`$ directly from its definition (22), taking into account expressions (19) and (21) for $`\mathrm{\Lambda }^{abc}`$. It is given by $$P_{[ij]}+ke\left(c_1+\frac{2}{9}c_3\right)T_{(0)ij}+ke\left(c_1\frac{4}{9}c_3\right)T_{[i\|(0)\|j]}=0.$$ $`(30)`$ In the time gauge the term $`T_{(0)ij}=_ie_{(0)j}_je_{(0)i}`$ vanishes. We will return to this expression latter together with expression (26). The third term to be considered in $`L`$ is $`N^kP^{ai}T_{aik}`$. In view of the time gauge condition it reads $$N^kP^{ai}T_{aik}=N^kP^{(i)j}T_{(i)jk}.$$ $`(32)`$ Lastly, we work out the remaining term $`2Nke\mathrm{\Lambda }^{aij}T_{aij}`$. The time gauge condition simplifies this expression in two respects. First, the term for which $`a=(0)`$ vanishes. Thus $`2Nke\mathrm{\Lambda }^{aij}T_{aij}=2Nke\mathrm{\Lambda }^{kij}T_{kij}`$. Second, because of (25) it can be shown by explicit calculations that $`\mathrm{\Lambda }^{kij}T_{kij}`$ does not contain terms of the type $`\mathrm{\Delta }_{a0j}=T_{a0j}`$. Therefore $`\mathrm{\Lambda }^{kij}T_{kij}`$ is totally projected in the spacelike hypersurface. We are now in a position of bringing back expressions (26)-(30) to the Lagrangian density (23). Before carrying out the substitution we can establish the conditions under which the Lagrangian density will be exempt of the terms $`T_{a0j}`$. From expression (26) we observe that we must demand $$c_1+c_2=0.$$ $`(33)`$ Next we see that the last line of (29), which contains several $`\mathrm{\Delta }_{a0j}`$ type terms, is discarded if we require $$c_1\frac{4}{9}c_3=0.$$ $`(34)`$ We observe from (30) that (34) ensures that $`P_{[ij]}`$ vanishes in the time gauge, $$P_{[ij]}=0,$$ $`(35)`$ which in turn makes (29) exempt of the term $`\frac{1}{2}P_{[ij]}\mathrm{\Lambda }^{(0)ij}`$. Thus $`P_{[ij]}`$ will enter the Hamiltonian density $`H=p\dot{q}L`$ multiplied by a Lagrange multiplier. We can finally provide the ultimate expression of $`L`$ by collecting terms that multiply the lapse and shift functions. We choose to write it in terms of the constant $`c_1`$. Note that $`\mathrm{\Lambda }^{kij}`$ can now be substituted by $`\mathrm{\Sigma }^{kij}`$, which is a function of $`e_{(i)j}`$ only. The final expression reads $$L=P^{(i)j}\dot{e}_{(i)j}+NC+N^kC_k+\lambda ^{ij}P_{[ij]}_i(3c_1keT^i+N_kP^{ki}),$$ $`(36)`$ where $`\{\lambda ^{ij}\}`$ are Lagrange multipliers. The Lagrangian density above is invariant under the global SO(3) group. The Hamiltonian and vector constraints are given by $$C=\frac{1}{6c_1ke}\left(P^{ij}P_{ji}\frac{1}{2}P^2\right)+ke\mathrm{\Sigma }^{kij}T_{kij}_i(3c_1keT^i),$$ $`(37)`$ and $$C_k=e_{(j)k}_iP^{(j)i}+P^{(j)i}T_{(j)ik},$$ $`(38)`$ where $`T^i=\overline{g}^{ik}T_k=\overline{g}^{ik}e^{(m)j}T_{(m)jk}`$; $`\mathrm{\Sigma }^{kij}`$ is defined by (8) together with conditions (33) and (34). These latter conditions reduce the theory defined by (4) to a one-parameter theory. IV. Discussion We observed that the Legendre transform has reduced the three-parameter to a one-parameter family of teleparallel theories. A consistent implementation of the Legendre transform is a necessary condition for the Hamiltonian formulation, but not sufficient. The complete canonical formulation demands further crucial investigations. It is also necessary to verify whether the constraints constitute a first class set, namely, whether the algebra of constraints “closes”. In the TEGR the Hamiltonian formulation and the constraint algebra have been obtained in . The Hamiltonian constraint of the latter is very similar to (37), except for the presence of $`c_1`$ in the three terms of $`C`$. By making $`c_1=\frac{2}{3}`$ expression (37) becomes precisely the Hamiltonian constraint of . Constraint (38) is the same as in . Let us write $`\mathrm{\Sigma }^{kij}`$ that appears in (37) in terms of $`c_1`$, using definition (8) and conditions (33) and (34). It reads $$\mathrm{\Sigma }^{kij}=\frac{3c_1}{2}\left(\frac{2}{3}X^{kij}\frac{2}{3}Y^{kij}+\frac{3}{2}Z^{kij}\right)=\frac{3c_1}{2}\mathrm{\Sigma }_{_{_{TEGR}}}^{kij}.$$ $`(39)`$ where $`\mathrm{\Sigma }_{_{_{TEGR}}}^{kij}`$ is restricted to the three-dimensional spacelike hypersurface and is defined in similarity with (11). By means of conditions (33) and (34) Lagrangian density (9) becomes $$L_0=\frac{3c_1}{2}ke\mathrm{\Sigma }_{_{_{TEGR}}}^{abc}T_{abc}.$$ $`(40)`$ We observe then that by defining $$k^{}=\frac{3c_1}{2}k,$$ we can rewrite (37) according to $$C=\frac{1}{4k^{}e}\left(P^{ij}P_{ji}\frac{1}{2}P^2\right)+k^{}e\mathrm{\Sigma }_{_{_{TEGR}}}^{kij}T_{kij}_i(2k^{}eT^i).$$ $`(41)`$ Except for $`k^{}`$ expression above is exactly the Hamiltonian constraint of ref. . The constant $`k^{}`$ does not affect the evaluation of Poisson brackets between (38), (41) and $`P^{[ij]}`$. Thus we conclude that the constraint algebra determined by (41) and (38) is exactly the same of ref. . Therefore (37) (or (41)) and (38) are first class constraints. As a consequence field quantities have a well defined time evolution. The fixing of $`k^{}`$ is related to the Newtonian behaviour of the gravitational field. Hayashi and Shirafuji have analyzed the Lagrangian field equations derived from (4). In particular they have investigated the conditions under which a correct Newtonian approximation is obtained by studying solutions of the field equations that yield static and isotropic gravitational fields. Without imposing any a priori restriction on the parameters $`c_i`$ they concluded that the Newtonian limit is verified for a class of solutions provided $$c_2=\frac{(c_1\frac{2}{3})}{(1\frac{9}{8}c_1)}\frac{2}{3}.$$ No condition fixes $`c_3`$. By imposing the mandatory condition (33), $`c_1+c_2=0`$, in the expression above it follows that $`c_2=\frac{2}{3}`$ and $`c_1=\frac{2}{3}`$. Hence we finally arrive at the TEGR. Lenzen and Baekler et. al. have shown the emergence of free functions in exact torsion solutions in Poincaré gauge theories (PGT) of gravity. Therefore it is worth examining this question here. The Hamiltonian formulation developed above provides a suitable framework for such analysis. The emergence of free functions is related to the selection of appropriate triads (tetrads) for the space (space-time). The counting of degrees of freedom here is the same as in the usual ADM (metrical) formulation, except that there are 9 triad components in (36) rather than 6 metric functions as in the ADM formulation (the imposition of $`P_{[ij]}=0`$ together with the field equation $`\dot{e}_{(i)j}(x)=\{e_{(i)j}(x),H\}`$ leads to an expression for $`\lambda _{ij}`$ in terms of $`e_{(i)j}`$). Therefore there are 3 extra (undetermined) triad components. However we have discussed in that these 3 components may be fixed in the context of isolated material systems by the asymptotic behaviour in the limit $`r\mathrm{}`$, $$e_{(i)j}\eta _{ij}+\frac{1}{2}h_{ij}(\frac{1}{r}),$$ $`(42)`$ where $`h_{ij}=h_{ji}`$ is the first space dependent term in the asymptotic expansion of the metric tensor $`g_{ij}`$. It is not difficult to verify that this condition fixes uniquely a triad to a three-dimensional metric tensor. In fact this condition has already been suggested by Møller for the same purpose. It turns out that in the TEGR condition (42) is essential in order obtain the ADM energy out of the scalar density $`_i(eT^i)`$ (with appropriate multiplicative constants) in the Hamiltonian constraint. A further condition on the triads is also essential in the TEGR, mainly in respect to the definition of gravitational energy: we must have $`T_{(i)jk}=0`$ if we make the physical parameters of the metric tensor (such as mass, angular momentum, etc) vanish. Triads that lead to a vanishing torsion tensor in any coordinate system are called reference space triads (all applications of the TEGR have taken into account the reference space triads). The two conditions above associate uniquely a metric tensor to triads (tetrads) components, make the latter exempt of free functions and lead to a well defined notion of gravitational energy. Baekler and Mielke consider the Hamiltonian formulation of the most general class of PGT theories, constructed out of tetrad fields $`e_{a\mu }`$ and connections $`\omega _{\mu ab}`$. They claim that a first class algebra is achieved irrespective of any prior gauge condition on $`\omega _{0ab}`$, and for a theory with arbitrary multiplicative constants for the squared torsion and curvature terms (the $`a_i`$ constants of the PGT theories are related to $`c_i`$ according to $`c_1=\frac{2}{3}a_1`$, $`c_2=\frac{1}{3}a_2`$ and $`c_3=3a_3`$). In view of this result the time evolution of field quantities would be, in principle, well defined (Hecht et. al. also consider the initial value problem for some PGT theories; they argue that there is freedom in the choice of the PGT parameters in the Lagrangian density such that the theory acquires a mathematically well defined initial value problem, but no Hamiltonian analysis is carried out). Certainly the analysis of is not in agreeement with that of ref., where the fixation of $`\omega _{0ab}`$ is mandatory. However it was shown by Kopczyński and Müller-Hoissen and Nitsch that the TEGR defined in terms of tetrad fields and connections $`\omega _{\mu ab}`$, supplemented by the condition of vanishing curvature (as developed in ) faces difficulties with respect to the Cauchy problem. They have shown that in general six components of the torsion tensor are not determined from evolution of the initial data. On the other hand in the context of the constraints of the theory become a first class set provided we fix the six quantities $`\omega _{0ab}=0`$ before varying the action (this point is also discussed in ). Although we have no proof we believe that the two properties above (the failure of the Cauchy problem and the fixation of $`\omega _{0ab}`$) are related to each other. For this reason we dispense with the constraint of vanishing curvature and consider the theory defined by (2). We note finally that a similar analysis has been developed by Blagojević and Nikolić, who considered PGT theories of the type $`R+R^2+T^2`$. Although the latter has been worked out in the framework of Riemannian geometry, with both tetrad and connection fields, it would be possible, in principle, to establish a comparison of their work with our analysis. However, the constraints of Ref. are only formally indicated. They are not explicitly expressed in terms of canonical variables as in (37) and (38), and therefore an objective comparison cannot be made. Moreover the constraint algebra in the particular case $`c_1=\frac{2}{3},c_2=\frac{2}{3}`$, $`c_3=\frac{3}{2}`$, as obtained in ref., has not been established in this earlier investigation. A recent investigation on the Hamiltonian formulation of PGT theories has been carried out along the lines of . Instead of actually performing an ordinary Legendre transform the authors make use of the if constraint formalism of . For the restricted sector of torsion squared terms they obtain a three parameter Hamiltonian density with second class constraints. However their Hamiltonian analysis does not single out the “viable” conditions $`a_1+2a_3=0`$ and $`2a_1+a_2=0`$ (in the notation of ), which they take into account in order to carry out their analysis. These conditions are enforced by hand and correspond precisely to conditions (33) and (34).
warning/0002/cond-mat0002163.html	ar5iv	text	# Convergence to the critical attractor of dissipative maps: Log-periodic oscillations, fractality and nonextensivity ## I introduction Nonlinear low-dimensional dissipative maps can describe a great variety of systems with few degrees of freedom. The underlying nonlinearity can induce the system to exhibit a complex behavior with quite structured paths in the phase space. The sensitivity to initial conditions is a relevant aspect associated to the structure of the dynamical attractor. In general, the sensitivity is measured as the effect of any uncertainty on the system’s variables. For systems exhibiting periodic or chaotic orbits, the effect of any uncertainty on initial conditions depicts an exponential temporal evolution with $`\xi (t)lim_{\mathrm{\Delta }x(0)0}\mathrm{\Delta }x(t)/\mathrm{\Delta }x(0)e^{\lambda t}`$, where $`\lambda `$ is the Lyapunov exponent, $`\mathrm{\Delta }x(0)`$ and $`\mathrm{\Delta }x(t)`$ are the uncertainties at times $`0`$ and $`t`$. When the Lyapunov exponent $`\lambda <0`$, $`\xi (t)`$ characterizes the rate of contraction towards periodic orbits. On the other hand, for $`\lambda >0`$, it characterizes the rate of divergence of chaotic orbits. At bifurcation and critical points (i.e., onset to chaos) the Lyapunov exponent $`\lambda `$ vanishes. Recently, it was shown that this feature is related to a power-law sensitivity to initial conditions on the form $$\xi (t)=[1+(1q)\lambda _qt]^{1/(1q)},$$ (1) with $`\lambda _q`$ defining a characteristic time scale after which the power-law behavior sets up. A quantitative way to measure the sensitivity to initial conditions is to follow, from a particular partition of the phase space, the temporal evolution of the number of cells $`W(t)`$ occupied by an ensemble of identical copies of the system. For periodic and chaotic orbits $`W(t)=W(0)e^{\lambda t}`$. In the particular case of equiprobability, the well-known Pesin equality reads $`K=\lambda `$ if $`\lambda 0`$ with $`K`$ being the Kolmogorov-Sinai entropy defined as the variation per unit time of the standard Boltzmann-Gibbs entropy $`S=p_i\mathrm{ln}p_i`$. This equality provides a link between the sensitivity to initial conditions and the dynamic evolution of the relevant entropy. At bifurcation and critical points and for an ensemble of initial conditions concentrated in a single cell, i.e. $`W(0)=1`$, it has been shown that $$W(t)=[1+(1q)K_qt]^{1/(1q)},$$ (2) with $`K_q`$ being the generalized Kolmogorov-Sinai entropy defined as the rate of variation of the non-extensive Tsallis entropy $`S_q=(1p_i^q)/(q1)`$. The Pesin equality can be generalized as $`K_q=\lambda _q`$ if $`\lambda _q0`$. Tsallis entropies have been successfully applied to recent studies of a series of non-extensive systems and provided a theoretical background to the understanding of some of their unusual physical properties. The expansion towards the critical attractor of an ensemble of initial conditions concentrated around the inflexion point of the map can be characterized by a proper $`S_q`$ evolving at a constant rate. Scaling arguments have shown that the appropriate entropic index $`q`$ is related to the multifractal structure of the critical dynamical attractor by $$\frac{1}{1q}=\frac{1}{\alpha _{min}}\frac{1}{\alpha _{max}}$$ (3) where $`\alpha _{min}`$ and $`\alpha _{max}`$ are the extremal singularity strengths of the multifractal spectrum of the critical attractor. The above scaling relation has been shown to hold for the families of generalized logistic and circle maps. However, the temporal evolution of critical dynamical systems can be strongly dependent on the particular initial ensemble. Although some scaling laws can be found for an ensemble of initial conditions concentrated around the map inflexion point, these are usually not universal with respect to a general ensemble. In this work, we are going to numerically investigate the critical temporal evolution of the volume of the phase space occupied by an ensemble of initial conditions spread over the entire phase space. This ensemble is expected to contract towards the critical attractor. Using a family of one-dimensional generalized logistic maps having $`d_f<1`$, we will perform a detailed study of the parametric dependence of $`W(t)`$ on the fractal dimension of the critical attractor. Due to the discrete scale invariance of the critical attractor, the convergence displays log-periodic oscillations. We are also going to explore the dependence of the amplitude of these oscillations with respect to the attractor’s fractal dimension. Further, the behavior of $`W(t)`$ will be investigated for the one-dimensional critical circle map having $`d_f=1`$. For this map, the temporal evolution is expected to display distinct trends since the critical attractor is dense. ## II The convergence to the critical attractor of generalized logistic maps Logistic-like maps are the simplest one-dimensional nonlinear dynamical systems which allow a close investigation of a series of critical exponents related to the onset of chaotic orbits. This family reads $$x_{t+1}=1a\|x_t\|^z;(z>1;0<a<2;t=0,1,2,\mathrm{};x_t[1,1]).$$ (4) Here $`z`$ is the inflexion of the map in the neighborhood of the extremal point $`\overline{x}=0`$. These maps are well known to have topological properties not dependent of $`z`$. However, the metrical properties, such as Feigenbaum exponents and the multifractal spectrum of the critical attractor do depend on $`z`$. In particular the fractal dimension of the critical attractor $`d_f(z)<1`$ and therefore it does not fill a finite fraction of the phase space. For a set of initial conditions spread in the vicinity of the inflexion point, it was found that the volume in phase space occupied by the ensemble grows following a rich pattern with the upper bounds $`W_{max}(t)`$ governed by a power-law $`W_{max}(t)t^{1/(1q)}`$, where $`q`$ is the entropic index characterizing the relevant Tsallis entropy that grows at a constant rate. It has been shown that the dynamic exponent $`1/(1q)`$ is directly related to geometric scaling exponents related to the extremal sets of the dynamic attractor. Due to the presence of long-range spatial and temporal correlations at criticality, one expects the critical exponent governing the temporal evolution to be sensitive to the particular initial ensemble. Indeed, the multifractal spectrum characterizing the critical dynamical attractor indicates that an infinite set of exponents are needed to fully characterize the scaling behavior. In particular, an ensemble consisting of a set of identical systems whose initial conditions is spread over the entire phase-space is a common one when studying non-linear as well as thermodynamical systems. Here, we will follow the dynamic evolution, in phase space, of an ensemble of initial conditions uniformly distributed over the phase-space and explore its relation with the generalized fractal dimensions of the critical attractor. In practice, a partition of the phase space on $`N_{box}`$ cells of equal size is performed and a set of $`N_c`$ identical copies of the system is followed whose initial conditions are uniformly spread over the phase-space. The ratio $`r=N_c/N_{box}`$ is a control parameter giving the degree of sampling of the phase-space. Within the non-extensive Tsallis statistics, there is a proper entropy $`S_q`$ evolving at a constant rate such that $$K_q=\underset{N_{box}\mathrm{}}{lim}[S_q(t)S_q(0)]/t$$ (5) goes to a constant value as $`t\mathrm{}`$. Notice that $`K_q<0`$ for the process of convergence towards the critical attractor. Assuming that all cells of the partition are occupied with equal probability, the entropy $`S_q(t)`$ can be written as $$S_q(t)=\frac{1_{i=1}^{W(t)}p_i^q}{q1}=\frac{W(t)^{1q}1}{1q}$$ (6) The last two equations imply that the number of occupied cells evolves in time as $$W(t)=[W(0)^{1q}+(1q)K_qt]^{1/(1q)}$$ (7) with the exponent $`\mu =1/(q1)>0`$ governing the asymptotic power-law decay. In figure 1, we show our results for $`W(t)/N_{box}`$ in the standard logistic map with inflexion $`z=2`$ and from distinct partitions of the phase space with sampling ratio $`r=0.1`$ . We observe that, after a short transient period when $`W(t)`$ is nearly constant, a power-law contraction of the volume occupied by the ensemble sets up. $`W(t)`$ saturates at a finite fraction corresponding to the phase space volume occupied by the critical attractor on a given finite partition. The saturation is postponed when a finer partition is used once the fraction occupied by the critical attractor vanishes in the limit $`N_{box}\mathrm{}`$. In figure 2, we show $`W(t)/N_{box}`$ for a given fine partition of the phase-space and distinct sampling ratios $`r`$. We notice that the crossover regime to the power-law scaling is quite short for large values of $`r`$ so that a clear power-law scaling regime sets up even at early times. This feature is consistent with eq.(7) which states that the crossover time $`\tau `$ scales as $`\tau 1/W(0)^{q1}`$. Further, the scaling regime exhibits log-periodic oscillations once the multifractal nature of the critical attractor is closely probed by such dense ensemble. A general form for $`W(t)`$ reflecting the discrete scale invariance of the attractor can be written as $$W(t)=t^\mu P\left(\frac{\mathrm{ln}t}{\mathrm{ln}\lambda }\right)$$ (8) where $`P`$ is a function of period unity and $`\lambda `$ is the characteristic scaling factor between the periods of two consecutive oscillations. These log-periodic oscillations have been observed in a large number of systems exhibiting discrete scale invariance. In general the amplitude of these oscillations ranges form $`10^4`$ up to $`10^1`$. Keeping only the first term in a Fourier series of $`P(\mathrm{ln}t/\mathrm{ln}\lambda )`$, one can write $`W(t)`$ in the form $$W(t)=c_0t^\mu \left[1+2\frac{c_1}{c_0}\mathrm{cos}\left(2\pi \frac{\mathrm{ln}t}{\mathrm{ln}\lambda }+\varphi \right)\right].$$ (9) Log-periodic modulations correcting a pure power-law have been found in several systems, as, for example, diffusion-limited-aggregation, crack growth, earthquakes and financial markets. It has also been observed in thermodynamic systems with fractal-like energy spectrum. The factors controlling the log-periodic relative amplitude $`2c_1/c_0`$ are not well known for most of the systems where it has been observed. In the present study, we can closely investigate the factors which may control these amplitudes by measuring it as a function of the map inflexion $`z`$ for a fixed partition and sampling ratio (see figure 3). We found that these oscillations have amplitudes decaying exponentially with $`z`$ as shown in figure 4. It is interesting to point out that the fractal dimension of the attractor is a monotonically decreasing function of $`z`$. Therefore, the above trend indicates a possible correlation between the amplitude of the log-periodic oscillations and the fractal dimension of the dynamical attractor. We also measured the critical exponent $`\mu `$ as a function of the map inflexion $`z`$. Our results are summarized in the Table. It is a decreasing function of $`z`$ as can be seen in figure 5. The volume occupied by the ensemble depicts a fast contraction for $`z1`$ where the fractal dimension is small. On the other side, a very slow contraction is observed for large values of $`z`$, pointing towards a saturation or at most to a logarithmic decrease of $`W(t)`$ in the limit of dense attractors. We would like to point out here that the exponent governing the expansion of the volume occupied by an ensemble of initial conditions concentrated around the inflexion point exhibits a reversed trend. Although scaling arguments have shown that this exponent can be written in terms of scaling exponents characterizing the extremal sets in the attractor, we could not devise a simple scaling relation between $`\mu `$ and the multifractal singularity spectrum. However, we observed that, when plotted against the fractal dimension of the attractor as shown in figure 6, the dynamic exponent $`\mu `$ is very well fitted by $`\mu (1d_f)^2`$, which indicates $`d_f`$ as the relevant geometric exponent coupled to the dynamics of the uniform ensemble. We would like to mention here that the same dynamic exponents were obtained for the generalized periodic maps which belong to the same universality class of logistic-like maps. ## III The convergence to the critical attractor of the circle map The results from the previous section indicate that a slow convergence to the critical attractor shall be expected for dense critical attractors. However, it is not clear in what fashion this convergence will take place when the dynamical attractor fills the phase space with a multifractal probability density as occurs for the one-dimensional critical circle map $$\theta _{t+1}=\theta _t+\mathrm{\Omega }\frac{1}{2\pi }\mathrm{sin}(2\pi \theta _t)mod(1)$$ (10) where $`0\theta _t<1`$ is a point on a circle. The circle map describes dynamical systems possessing a natural frequency $`\omega _1`$ which are driven by an external force of frequency $`\omega _2`$ ($`\mathrm{\Omega }=\omega _1/\omega _2`$ is the bare winding number) and belongs to the same universality class of the forced Rayleigh-Bénard convection. For $`\mathrm{\Omega }=0.606661\mathrm{}`$ the circle map has a cubic inflexion ($`z=3`$) in the vicinity of the point $`\overline{\theta }=0`$. Starting from a given point on the circle, it generates a quasi-periodic orbit which fills the phase-space and the dynamical attractor is a multifractal with fractal dimension $`d_f=1`$. In figure 7 we show our results for the temporal evolution of the phase-space volume occupied by an ensemble of initial conditions uniformly spread over the circle. $`W(t)`$ exhibits a rich pattern which resembles the one observed for the sensitivity function associated to the expansion of the phase-space from initial conditions concentrated around the inflexion point. However, $`W(t)`$ does not present any power-law regime. Instead, the lower bounds display a slow logarithmic decrease with time, saturating at a finite volume fraction. The saturation is a feature related to the finite partition used in the numerical calculation. This minimum decreases logarithmically with the number of cells in the phase-space as shown in figure 8. We also observed the same behavior for generalized circle maps with an arbitrary inflexion $`z`$. The critical attractors within this family have all $`d_f=1`$ although they exhibit a $`z`$-dependent multifractal singularity spectra. The $`z`$-independent scenario for $`W(t)`$ corroborates the conjecture that $`d_f`$ is the relevant geometric exponent coupled to the dynamics of the uniform ensemble. ## IV Summary and Conclusions In this work, we studied the temporal evolution in phase space of an ensemble of identical copies of one-dimensional nonlinear dissipative maps. We found that the phase-space volume occupied by an initially uniform ensemble displays a power-law decay with log-periodic oscillations whenever the dynamical attractor has a fractal dimension $`d_f<1`$. The amplitude of the oscillations was found to depict a monotonic parametric dependence on $`d_f`$. For dense multifractal attractors, $`W(t)`$ presents only a slow logarithmic contraction of its lower bounds followed by a rich pattern. The critical exponent characterizing the contraction of the uniform ensemble was found to have no direct relation to the one governing the expansion from a set of initial conditions concentrated around the inflexion point. In particular, no power-law was found for the contraction in the standard and generalized circle maps, in contrast to the $`z`$-dependent power-law expansion. These results indicate that the relevant Tsallis entropy evolving at a constant rate (modulated by log-periodic oscillations) is characterized by an entropic index $`q`$ that depends on the initial ensemble. It would be valuable to investigate the possible existence of classes of ensembles with a common dynamics in phase-space and, therefore, characterized by the same entropic index $`q`$. The non-universality of $`q`$ with respect to the initial ensemble is related to the multifractal character of the dynamical attractor. However, as for the ensemble concentrated at the vicinity of the inflexion point, the exponent governing the dynamics of the uniform ensemble is coupled to a geometric scaling exponent, in particular to the proper fractal dimension of the attractor. This result comes in favor of the concept that the degree of nonextensivity of the entropy measure evolving at a constant rate is related to the fractal nature of the dynamical attractor. ## V acknowledgments UT acknowledges the partial support of BAYG-C program of TUBITAK (Turkish agency) as well as CNPq and PRONEX (Brazilian agencies). This work was partially supported by CNPq and CAPES (Brazilian research agencies). MLL would like to thank the hospitality of the Physics Department at Universidade Federal de Pernambuco during the Summer School 2000 where this work was partially developed. ## FIGURE AND TABLE CAPTIONS Figure 1 - The volume occupied by the ensemble $`W(t)`$ as a function of time in the standard logistic map ($`z=2`$) and with sampling ratio $`r=0.1`$. From top to bottom $`N_{box}=2000,8000,32000,128000`$. Figure 2 - The volume occupied by the ensemble $`W(t)`$ as a function of time in the standard logistic map ($`z=2`$) and for a partition containing $`N_{box}=128000`$ cells. Notice the emergence of log-periodic oscillations for large sampling ratios. Figure 3 - The periodic function $`W(t)/(c_0t^\mu )`$ versus time within the scaling regime and for $`r=10`$. Data from map inflexions $`z=1.1,1.25,1.5,2.0`$ are shown. The amplitude of the oscillations decreases monotonically as $`z`$ increases, but the characteristic scaling factor between the periods of two consecutive oscillations is roughly $`z`$-independent. Figure 4 - The amplitude of the log-periodic oscillations $`2c_1/c_0`$ as a function of the map inflexion $`z`$ for sampling ratio $`r=10`$. The monotonic decrease of the oscillations indicates a close relation between these and the fractal dimension of the underlying dynamical attractor. Figure 5 - The dynamic exponent $`\mu `$ governing the contraction of the occupied phase space volume \[$`W(t)t^\mu `$\] as a function of the map inflexion $`z`$. Figure 6 \- The parametric dependence of the dynamic exponent $`\mu `$ with the fractal dimension $`d_f`$ of the critical attractor. It is very well fitted to the form $`\mu (1d_f)^2`$, indicating that $`d_f`$ seems to be the relevant geometric exponent coupled to the dynamics of the uniform ensemble. Figure 7 - The volume occupied by the ensemble $`W(t)`$ as a function of time in the standard critical circle map. The lower bounds display a slow logarithmic decay with time saturating at a finite volume fraction due to the finite partition of the phase space. Figure 8 - The asymptotic lower bounds for the occupied volume in the phase space versus the number of cells $`N_{box}`$. The logarithmic decay agrees with the prediction that $`\mu (d_f1)0`$. The same behavior was observed for the family of generalized circle maps and corroborates the conjecture that $`d_f`$ is the relevant geometric exponent coupled to the dynamics of the uniform ensemble. Table \- Numerical values, within the $`z`$-generalized family of logistic maps, of: i) the dynamic exponent $`\mu `$ governing the contraction towards the critical attractor of the uniform ensemble; ii) the entropic index $`q`$ of the proper Tsallis entropy decreasing at a constant rate; iii) the fractal dimension $`d_f`$ of the critical attractor. These values also hold for the generalized periodic maps. The last line represents our results for the $`z`$-generalized circle maps. Table \| $`z`$ \| $`\mu =1/(1q)`$ \| $`q`$ \| $`d_f`$ \| \| --- \| --- \| --- \| --- \| \| $`1.10`$ \| $`1.62\pm 0.02`$ \| $`1.62\pm 0.01`$ \| $`0.32\pm 0.02`$ \| \| $`1.25`$ \| $`1.23\pm 0.01`$ \| $`1.81\pm 0.01`$ \| $`0.40\pm 0.01`$ \| \| $`1.5`$ \| $`0.95\pm 0.01`$ \| $`2.05\pm 0.01`$ \| $`0.47\pm 0.01`$ \| \| $`1.75`$ \| $`0.80\pm 0.01`$ \| $`2.25\pm 0.015`$ \| $`0.51\pm 0.01`$ \| \| $`2.0`$ \| $`0.71\pm 0.01`$ \| $`2.41\pm 0.02`$ \| $`0.54\pm 0.01`$ \| \| $`2.5`$ \| $`0.59\pm 0.01`$ \| $`2.70\pm 0.02`$ \| $`0.58\pm 0.01`$ \| \| $`3.0`$ \| $`0.515\pm 0.005`$ \| $`2.94\pm 0.02`$ \| $`0.60\pm 0.01`$ \| \| $`5.0`$ \| $`0.395\pm 0.005`$ \| $`3.53\pm 0.03`$ \| $`0.66\pm 0.01`$ \| \| $`z`$-circular \| \| \| \| \| maps \| $`0.0`$ \| $`\mathrm{}`$ \| $`1.0`$ \|
warning/0002/hep-ph0002040.html	ar5iv	text	# Time-reversal violating rotation of polarization plane of light in gas placed in the electric field ## 1 Introduction Violation of time reversal invariance was discovered more then 30 years ago in $`K`$ meson decay . Up to now it remains one of the great unsolved problems in elementary particle physics. Many experiments were devoted to the search of any other manifestation of time reversal noninvariance. Among them for example are measurements of electrical dipole moment (EDM) of neutrons , atoms and molecules . No EDM was found but these experiments impose strong restrictions on theory. In particular search of EDM in heavy atoms set a tight limits on parameters of electron - nucleon $`PT`$ violating interactions and value of electron EDM. At present more precise schemes of experiment are actively discussed. One of them is the observation of the light polarization plane rotation caused by pseudo-Zeeman splitting of magnetic sub-levels of atom with nonzero EDM in electric field. This effect arises due to interaction $`W=\stackrel{}{d}_a\stackrel{}{E}`$ of atomic EDM with external electric field . We should note here that discussions of these experiments take into account only static EDM of atom. According to atom has another $`PT`$ noninvariant characteristic that describes its response to the external electric field. It is a $`P`$ and $`T`$ \- odd polarizability $`\beta _E^{PT}`$ that arises in electric field due to interference of $`PT`$ \- odd and Stark - induced transition amplitudes. As was shown in $`PT`$ \- odd polarizability also leads to the rotation of photon polarization plane and circular dichroism of atomic gas in external electric field. This contribution should exist even in a hypothetical case when atomic EDM in a ground and excited states are occasionally equal to zero and pseudo-Zeeman splitting of atomic levels is absent. Unlike rotation of polarization plane due to atomic EDM that manifest Macaluso-Corbino dependence of angle on light frequency, rotation caused by $`\beta _E^{PT}`$ is a kinematic analog of Faraday rotation in a Van-Vleck paramagnetic. Moreover $`PT`$ noninvariant polarizability $`\beta _E^{PT}`$ cause magnetization of atom in electric field . This magnetization in turn induce magnetic field $`H_{ind}`$. The energy of interaction of magnetic moment of atom with this field is $`W_H=\stackrel{}{\mu }_a\stackrel{}{H}_{ind}(E)`$. Therefore the total splitting of atomic levels that cause rotation of polarization plane of light is equal to $`W_H=\stackrel{}{d}_a\stackrel{}{E}\stackrel{}{\mu }_a\stackrel{}{H}_{ind}(E)`$. This splitting appears even if atomic EDM $`d_a`$ is equal to zero. In this paper the mechanisms of rotation of polarization plane of light due to both $`PT`$ noninvariant electron-nucleon interactions and nonzero electron EDM are considered. Estimates of expected angle of polarization plane rotation near highly forbidden magnetic dipole transition $`6S_{1/2}7S_{1/2}`$ in cesium ($`\lambda =539`$ nm) are performed. Possible experimental schemes to observe this rotation are discussed. ## 2 PT - odd mixing We start with the simplest case. Let us place an atom in a ground state $`s_{1/2}`$ to the electric field. If we take into account admixture of the nearest $`p_{1/2}`$ state due to $`P`$ and $`T`$ odd interactions and interaction with the electric field, then the wave function of atom takes the form. $$\|\stackrel{~}{s}_{1/2}>=\frac{1}{\sqrt{4\pi }}(R_0(r)R_1(r)(\stackrel{}{\sigma }\stackrel{}{n})\eta R_1(r)(\stackrel{}{\sigma }\stackrel{}{n})(\stackrel{}{\sigma }\stackrel{}{E})\delta )\|\chi _{1/2}$$ (1) Here $`\stackrel{}{\sigma }`$ \- are the Pauli matrices, $`\stackrel{}{n}=\stackrel{}{r}/r`$ is the unit vector along the direction of $`\stackrel{}{r}`$, $`R_0`$ and $`R_1`$ are radial parts of $`s_{1/2}`$ and $`p_{1/2}`$ wave functions respectively, $`\|\chi _{1/2}`$ is the spin part of wave function, $`\eta `$ and $`\delta `$ are the mixing coefficients due to $`P`$ and $`T`$ noninvariant interactions and electric field respectively. Let us consider orientation of electron spin in atom. In order to find the spatial distribution of spin direction we can calculate the matrix element of electron spin operator in respect to the spin part of atomic wave function. Only terms proportional to the product of electrical field strength $`\stackrel{}{E}`$ and $`PT`$ odd mixing coefficient $`\eta `$ are important for our consideration because only they cause $`PT`$ \- odd rotation of polarization plane of light. The change of spin direction due to these terms is $`\mathrm{\Delta }\stackrel{}{s}(\stackrel{}{r})`$ $`=`$ $`{\displaystyle \frac{\eta \delta }{8\pi }}R_1^2\chi _{1/2}\|(\stackrel{}{\sigma }\stackrel{}{n})\stackrel{}{\sigma }(\stackrel{}{\sigma }\stackrel{}{n})(\stackrel{}{\sigma }\stackrel{}{E})+(\stackrel{}{\sigma }\stackrel{}{E})(\stackrel{}{\sigma }\stackrel{}{n})\stackrel{}{\sigma }(\stackrel{}{\sigma }\stackrel{}{n})\|\chi _{1/2}`$ (2) $`=`$ $`{\displaystyle \frac{\eta \delta R_1^2}{8\pi }}\left(4\stackrel{}{n}(\stackrel{}{n}\stackrel{}{E})2\stackrel{}{E}\right)`$ The vector field $`4\stackrel{}{n}(\stackrel{}{n}\stackrel{}{E})2\stackrel{}{E}`$ is shown in Fig. 1. Since $`\mathrm{\Delta }\stackrel{}{s}`$ does not depend on initial direction of atomic spin, this spin structure appears even in non-polarized atom. Let us note that the spin vector averaged over spatial variables differs from zero and is directed along the vector of electric field strength $`\stackrel{}{E}`$. The photons with directions of angular moment parallel and antiparallel to the electric field will interact with such spin structure in a different ways causing rotation of polarization plane of light. Polarization of atoms produce magnetic moment of gas . Thus we have another interesting $`PT`$ \- odd effect. If we place gas in electric field, small magnetic field will appear. Magnetic field in turn will interact with magnetic moment of atom giving another contribution to the rotation of the polarization plane of light . According to if light propagates along the direction of electric field, then the amplitude of elastic coherent forward scattering of light by nonpolarized atoms has the form $$f_{ik}(0)=f_{ik}^{ev}+\frac{\omega ^2}{c^2}(i\beta _s^Pϵ_{ikl}n_{\gamma l}+i\beta _E^{PT}ϵ_{ikl}n_{El})$$ (3) Here $`f_{ik}^{ev}`$ is $`P`$ and $`T`$ invariant part of scattering amplitude, $`\beta _s^P`$ is the $`P`$ \- odd but $`T`$ \- even scalar atomic polarizability, $`\beta _E^{PT}`$ is $`P`$ and $`T`$ \- odd scalar polarizability of atom, $`\stackrel{}{n}_\gamma =\stackrel{}{k}/k`$ is the unit vector along the direction of photon propagation, $`\stackrel{}{n}_E=\stackrel{}{E}/E`$ is the unit vector along the direction of electric field, $`ϵ_{ijk}`$ is the third rank antisymmetric tensor. The refraction index of gas can be written as $$n=1+\frac{2\pi N}{k^2}f(0)$$ (4) where $`N`$ is the number of atoms per $`cm^3`$, $`k`$ is the photon wave vector. Using (3) expression (4) can be rewritten as follows $$n_\pm =1+\frac{2\pi N}{k^2}(f^{ev}(0)\frac{\omega ^2}{c^2}[\beta _s^P+\beta _E^{PT}(\stackrel{}{n}_E\stackrel{}{n}_\gamma )])$$ (5) Indices $`+`$ and $``$ stands for left and right circular polarization of incident light respectively. Angle of rotation of polarization plane has the form $$\varphi =\frac{1}{2}kRe(n_+n_{})=\frac{2\pi N\omega }{c}(\beta _s^P+\beta _E^{PT}(\stackrel{}{n}_E\stackrel{}{n}_\gamma ))$$ (6) Term proportional to $`\beta _s^P`$ describes well known phenomenon of $`P`$ \- odd but $`T`$ \- even rotation of polarization plane of light. Term proportional to $`\beta _E^{PT}`$ describes $`P`$ and $`T`$ noninvariant rotation of polarization plane of light about the direction of electric field. $`P`$ and $`T`$ \- odd rotation of polarization plane of light change sign with the reversal of electric field direction in contrast to $`P`$ odd but $`T`$ even rotation. This allow us to distinguish $`PT`$ \- odd rotation from the other possible mechanisms of rotation of polarization plane. Refraction index of gas has both real and imaginary parts. Since imaginary part of refraction index for left and right circularly polarized photons are also different due to $`P`$ and $`T`$ odd interactions, the admixture of circular polarization to the linearly polarized light traveling in gas (circular dichroism) appears. Let us consider $`PT`$ noninvariant polarizability $`\beta _E^{PT}`$. According to the tensor of dynamical polarizability of an atom (molecule) has the form $$\alpha _{ik}^n=\underset{m}{}\{\frac{\stackrel{~}{g}_n\|d_i\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_k\|\stackrel{~}{g}_n}{E_{em}E_{gn}\mathrm{}\omega }+\frac{\stackrel{~}{g}_n\|d_k\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_i\|\stackrel{~}{g}_n}{E_{em}E_{gn}+\mathrm{}\omega }\}$$ (7) where $`\|\stackrel{~}{g}_n>`$ and $`\|\stackrel{~}{e}_m>`$ are wave functions of atom in ground and excited states perturbed by electric field and $`PT`$ \- noninvariant interactions, $`d`$ is the operator of dipole transition, $`\omega `$ is the frequency of incident light, $`E_{em}`$ and $`E_{gn}`$ are the energies of atom in states $`\|\stackrel{~}{g}_n>`$ and $`\|\stackrel{~}{e}_m>`$ respectively. In general case atoms are distributed to the sub - levels of ground state $`g_n`$ with the probability $`P(n)`$. Therefore, $`\alpha _{ik}^n`$ should be averaged over all states $`n`$. As a result, the polarizability can be written $$\alpha _{ik}=\underset{n}{}P(n)\alpha _{ik}^n$$ (8) The tensor $`\alpha _{ik}`$ can be decomposed into the irreducible parts as $$\alpha _{ik}=\alpha _0\delta _{ik}+\alpha _{ik}^s+\alpha _{ik}^a.$$ (9) Here $`\alpha _0=\frac{1}{3}_i\alpha _{ii}`$ is the scalar, $`\alpha _{ik}^s=\frac{1}{2}(\alpha _{ik}+\alpha _{ki})\alpha _0\delta _{ik}`$ is the symmetric tensor of rank two, $`\alpha _{ik}^a=\frac{1}{2}(\alpha _{ik}\alpha _{ki})`$ is the antisymmetric tensor of rank two, $$\alpha _{ik}^a=\frac{\omega }{\mathrm{}}\underset{n}{}P(n)\underset{m}{}\{\frac{\stackrel{~}{g}_n\|d_i\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_k\|\stackrel{~}{g}_n\stackrel{~}{g}_n\|d_k\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_i\|\stackrel{~}{g}_n}{\omega _{em,gn}^2\omega ^2}\}$$ (10) where $`\omega _{em,gn}=(E_{em}E_{gn})/\mathrm{}`$. Let atoms (molecules) be nonpolarized. The antisymmetric part of polarizability (10) is equal to zero in the absence of T- and P- odd interactions . It should be reminded that for the P-odd and T-even interactions the antisymmetric part of polarizability differs from zero only for both the electric and magnetic dipole transitions . We can evaluate the antisymmetric part $`\alpha _{ik}^a`$ of the tensor $`\alpha _{ik}`$ of dynamical polarizability of atom (molecule), and, as a result, obtain the expression for $`\beta _E^{PT}`$ in the following way. According to (5) the magnitude of the $`PT`$-odd effect is determined by the polarizability $`\beta _E^{PT}`$ or by the amplitude $`f_\pm (0)`$ of elastic coherent scattering of a photon by an atom (molecule). If $`\stackrel{}{n}_E\stackrel{}{n}_\gamma `$ the amplitude $`f_\pm (0)`$ in the dipole approximation can be written as $`f_\pm =\omega ^2\alpha _{ik}e_i^{(\pm )}e_k^{(\pm )}/c^2=\omega ^2\beta _E^{PT}/c^2`$. As a result, in order to obtain the amplitude $`f_\pm `$, the polarizability (7) for photon polarization states $`\stackrel{}{e}=\stackrel{}{e}_\pm `$ should be found. Using (7) we can present the polarizability $`\beta _E^{PT}`$ as follows: $$\beta _E^{PT}=\frac{\omega }{\mathrm{}}\underset{n}{}P(n)\underset{m}{}\{\frac{\stackrel{~}{g}_n\|d_{}\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_+\|\stackrel{~}{g}_n\stackrel{~}{g}_n\|d_+\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_{}\|\stackrel{~}{g}_n}{\omega _{em,gn}^2\omega ^2}\}$$ (11) For further analysis the more detailed expressions for atom (molecule) wave functions are necessary. The constants of $`PT`$ noninvariant interactions are very small. Therefore we can use the perturbation theory. Let $`\|\overline{g}`$ and $`\|\overline{e}`$ be the wave function of ground and excited states of atom (molecule) interacting with an electric field $`\stackrel{}{E}`$ in the absence of $`PT`$ \- odd interactions. Switch on $`PT`$ noninvariant interaction $`(H_T0)`$. According to the perturbation theory the wave functions $`\|\stackrel{~}{g}`$ and $`\|\stackrel{~}{e}`$ can be written in this case as $`\|\stackrel{~}{g}>`$ $`=`$ $`\|\overline{g}>+{\displaystyle \underset{n}{}}\|n>{\displaystyle \frac{n\|H_T\|\overline{g}}{E_gE_n}}`$ $`\|\stackrel{~}{e}>`$ $`=`$ $`\|\overline{e}>+{\displaystyle \underset{n}{}}\|n>{\displaystyle \frac{n\|H_T\|\overline{e}}{E_eE_n}}`$ (12) where $`H_T`$ is Hamiltonian of $`T`$ noninvariant interactions. It should be mentioned that both numerator and denominator of (11) contain $`H_T`$. Suppose $`H_T`$ to be small one can represent the total polarizability $`\beta _E^{PT}`$ as the sum of two terms $$\beta _E^T=\beta _{mix}^T+\beta _{split}^T.$$ (13) Here $$\beta _{mix}^{PT}=\frac{\omega }{\mathrm{}}\underset{n}{}P(n)\underset{m}{}\{\frac{\stackrel{~}{g}_n\|d_{}\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_+\|\stackrel{~}{g}_n\stackrel{~}{g}_n\|d_+\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_{}\|\stackrel{~}{g}_n}{\omega _{\overline{e}m,\overline{g}n}^2\omega ^2}\}$$ (14) where $`\omega _{\overline{e}m,\overline{g}n}`$ does not include the $`PT`$ noninvariant shift of atomic levels, and $$\beta _{split}^{PT}=\frac{\omega }{\mathrm{}}\underset{n}{}P(n)\underset{m}{}\{\frac{\overline{g}_n\|d_{}\|\overline{e}_m\overline{e}_m\|d_+\|\overline{g}_n\overline{g}_n\|d_+\|\overline{e}_m\overline{e}_m\|d_{}\|\overline{g}_n}{\omega _{em,gn}^2\omega ^2}\}$$ (15) $$\omega _{em,gn}=(E_{em}(\stackrel{}{E})E_{gn}(\stackrel{}{E}))/\mathrm{}$$ It should be reminded that energy levels $`E_{e,m}(\stackrel{}{E})`$ and $`E_{g,n}(\stackrel{}{E})`$ contain shifts caused by interaction of electric dipole moment of atom with electric field $`\stackrel{}{E}`$ and magnetic moment of atom with T-odd induced magnetic field $`\stackrel{}{H}_{ind}(\stackrel{}{E})`$. Below we will consider nonpolarized atoms and small detuning of radiation frequency from atomic transition. Therefore (14) and (15) can be written as follows. $$\beta _{mix}^{PT}=\frac{1}{2\mathrm{}(2j_g+1)}\underset{n,m}{}\{\frac{\stackrel{~}{g}_n\|d_{}\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_+\|\stackrel{~}{g}_n\stackrel{~}{g}_n\|d_+\|\stackrel{~}{e}_m\stackrel{~}{e}_m\|d_{}\|\stackrel{~}{g}_n}{\omega _{\overline{e}m,\overline{g}n}\omega }\}$$ (16) $$\beta _{split}^{PT}=\frac{1}{2\mathrm{}(2j_g+1)}\underset{m,n}{}\{\frac{\overline{g}_n\|d_{}\|\overline{e}_m\overline{e}_m\|d_+\|\overline{g}_n\overline{g}_n\|d_+\|\overline{e}_m\overline{e}_m\|d_{}\|\overline{g}_n}{\omega _{em,gn}\omega }\}$$ (17) In this section we will study only rotation of polarization plane associated with $`\beta _{mix}`$. The rotation associated with $`\beta _{split}`$ will be considered in the next section. Due to Doppler shift resonance frequency of transition for a single atom depends on velocity of atom in a gas. In order to obtain expressions for absorption length and angle of polarization plane rotation we should average (14) over Maxwell distribution of atomic velocity. After standard calculations (see e. q. ) expressions take the form. $`\varphi ={\displaystyle \frac{2\pi N\omega }{c}}\mathrm{Re}\beta _{mix}^{PT}_v`$ $`=`$ $`\pi Nl{\displaystyle \frac{\omega }{\mathrm{\Delta }_D\mathrm{}c}}g(u,v)[\|\overline{A^+}\|^2\|\overline{A^{}}\|^2]`$ $`L^1=2k\mathrm{Im}n_\pm _v`$ $`=`$ $`4\pi N{\displaystyle \frac{\omega }{\mathrm{\Delta }_D\mathrm{}c}}f(u,v)\|\overline{A^\pm }\|^2`$ (18) where $`_v`$ denotes the averaging over atomic velocity, $`\|\overline{A^+}\|^2`$ and $`\|\overline{A^{}}\|^2`$ are the squares of transition amplitudes for left and right circularly polarized photons averaged over atomic polarization. $$\|\overline{A^\pm }\|^2=\frac{1}{(2j_g+1)}\underset{m_g}{}\stackrel{~}{g}\|d^\pm \|\stackrel{~}{e}\stackrel{~}{e}\|d^{}\|\stackrel{~}{g}$$ (19) $`\mathrm{\Delta }_D=\omega _0\sqrt{2kT/Mc^2}`$ is Doppler linewidth, $`f(u,v)`$ and $`g(u,v)`$ are equal to $$\begin{array}{c}g(u,v)\hfill \\ f(u,v)\hfill \end{array}\}=\begin{array}{c}\mathrm{Im}\hfill \\ \mathrm{Re}\hfill \end{array}\}\sqrt{\pi }e^{w^2}\left(1\mathrm{\Phi }(iw)\right)$$ (20) here $`w=u+iv`$, $`\mathrm{\Phi }(z)=\frac{2}{\sqrt{\pi }}_0^z𝑑te^{t^2}`$, $`u=(\omega \omega _0)/\mathrm{\Delta }_D`$, $`v=\mathrm{\Gamma }/2\mathrm{\Delta }_D`$, $`\mathrm{\Gamma }`$ is the recoil linewidth, Let us assume that electric field is small enough and we can use first order perturbation theory. Perturbed states $`\|\stackrel{~}{g}>`$ and $`\|\stackrel{~}{e}>`$ in this case have the form $`\|\stackrel{~}{g}>`$ $`=`$ $`\|g>+{\displaystyle \underset{n}{}}\|n>{\displaystyle \frac{n\|H_T\|g}{E_gE_n}}+{\displaystyle \underset{m}{}}\|m>{\displaystyle \frac{m\|\stackrel{}{d}\stackrel{}{E}_z\|g}{E_gE_m}}`$ $`\|\stackrel{~}{e}>`$ $`=`$ $`\|e>+{\displaystyle \underset{n}{}}\|n>{\displaystyle \frac{n\|H_T\|e}{E_eE_n}}+{\displaystyle \underset{n}{}}\|m>{\displaystyle \frac{m\|\stackrel{}{d}\stackrel{}{E}_z\|e}{E_eE_m}}`$ (21) Here $`H_T`$ is Hamiltonian of $`PT`$ noninvariant interactions, $`\|g>`$ and $`\|e>`$ are unperturbed ground and excited states of atom. Only terms proportional to products of $`H_T`$ and $`\stackrel{}{d}\stackrel{}{E}`$ leads to phenomenon of interest. Using (21) we can write. $$\|\overline{A^+}\|^2\|\overline{A^{}}\|^2=\frac{2}{2j_g+1}\mathrm{Re}\underset{m_g}{}(g\|d_+^{PT}\|ee\|d_{}^{St}\|gg\|d_{}^{PT}\|ee\|d_+^{St}\|g)$$ (22) where $`d_i^{PT}`$ is the admixture of $`E1`$ amplitude due to $`PT`$ \- odd interactions $$g\|d_i^{PT}\|e=\underset{m}{}\frac{g\|H_T\|mm\|d_i\|e}{E_mE_g}+\frac{g\|d_i\|mm\|H_T\|e}{E_mE_e}$$ (23) and $`d_i^{St}=\mathrm{\Lambda }_{ik}E_k`$ is the Stark - induced amplitude. $$g\|d_i^{St}\|e=E_k\mathrm{\Lambda }_{ik}=E_k\underset{n}{}\frac{g\|d_k\|nn\|d_i\|e}{E_nE_g}+\frac{g\|d_i\|nn\|d_k\|e}{E_nE_e}$$ (24) Here $`\mathrm{\Lambda }_{ik}`$ is the tensor of transition atomic polarizability. According to Stark - induced $`E1`$ amplitude has the form $$e\|d_ϵ^{St}\|g=\underset{q,q^{}}{}\mathrm{\Lambda }_{q,q^{}}E_qϵ_q^{}=\underset{K,Q}{}(1)^Q\mathrm{\Lambda }_Q^K(Eϵ)_Q^K,$$ (25) where $`E_q`$ is external electric field strength, $`ϵ_q^{}`$ is the strength of electric field in a laser wave, $`\mathrm{\Lambda }_Q^K`$ and $`(Eϵ)_Q^K`$ are the components of irreducible spherical tensors. Using Wigner - Ekhard theorem we can represent $`\mathrm{\Lambda }_Q^K`$ as follows. $$\mathrm{\Lambda }_Q^K=(1)^{j_em_e}\left(\begin{array}{ccc}j_e& K& j_g\\ m_e& Q& m_g\end{array}\right)\mathrm{\Lambda }^K$$ Reduced matrix elements $`\mathrm{\Lambda }^K`$ ($`K=0,1,2`$) are proportional to scalar, vector and tensor transition polarizability respectively. Due to orthogonality of $`3j`$-symbols only terms proportional to the vector part remains in (22) after summation over magnetic sub-levels. In order to obtain the angle of polarization plane rotation and absorption length for atoms with the nuclear spin, we should take into account hyperfine structure of atomic levels. After necessary transformations (see e. q. ) equation (18) can be rewritten using reduced matrix elements of corresponding transitions. $`\varphi `$ $`=`$ $`4\pi N_Fl{\displaystyle \frac{\omega }{\mathrm{}c\mathrm{\Delta }_D}}g(u,v){\displaystyle \frac{1}{3(2F_g+1)}}K^2Re(gd^{PT}eed^{St}g)`$ $`L^1`$ $`=`$ $`4\pi N_F{\displaystyle \frac{\omega }{\mathrm{}c\mathrm{\Delta }_D}}f(u,v){\displaystyle \frac{1}{3(2F_g+1)}}K^2\|g\|\|A\|\|e\|^2`$ (26) Here $`F_g`$, $`F_e`$ are total angular moments of atom in ground and exited states, $`j_g`$, $`j_e`$ are the total angular moments of electrons in these states, $`i`$ is the nuclear spin, $$N_F=N\frac{2F_g+1}{(2i+1)(2j_g+1)}$$ is the density of atoms with total moment $`F_g`$, $$K^2=(2F_g+1)(2F_e+1)\left\{\begin{array}{ccc}i& j_g& F_g\\ 1& F_e& J_e\end{array}\right\}$$ and $`gd^{St}e`$ is proportional to the $`\mathrm{\Lambda }^1`$. We assume that electric field is parallel to the direction of light propagation. ## 3 Rotation of polarization plane due to atomic EDM Presence of EDM in ground or excited state of atom also cause rotation of polarization plane of light. We can derive the expression for the angle of polarization plane rotation performing the calculations similar to those described in section 2, but using $`\beta _{split}`$ instead of $`\beta _{mix}`$. But in this case the calculations can be greatly simplified if we note that the mechanism of $`PT`$ noninvariant rotation caused by atomic EDM is analogous to the Faraday rotation of polarization plane in a week magnetic field. Indeed according to application of weak magnetic field to the atomic gas affect the refractive index in two ways: through the changes in the energies of the magnetic sub-levels and through the mixing of hyperfine states. If we consider only terms proportional to the magnetic field strength $`H`$ and neglect the higher order terms then the levels shift becomes $$\mathrm{\Delta }E_i=Hi\|\mu _z\|i$$ The magnetic field $`H`$ mixes states of the same $`F_z`$ but different $`F`$, so the state $`\|j`$ becomes $$\|\overline{j}=\|j\underset{kj}{}H_z\frac{\|kk\|\mu _z\|j}{E_kE_j}$$ If atom has an EDM then applied electric field affect the refraction index in the same way (see (15)). It shifts the atomic levels $$\mathrm{\Delta }E_i=Ei\|d_z\|i$$ and mixes the hyperfine states of atom with the same $`F_z`$ but different $`F`$ $$\|\overline{j}=\|j\underset{kj}{}E_z\frac{\|kk\|d_z\|j}{E_kE_j}$$ Therefore after substitutions $`EH`$, $`\mu _gd_g`$, $`\mu _ed_e`$ where $`d_e`$, $`d_g`$ are EDM of atom in ground and excited states, $`\mu _i`$ is the magnetic moment of state $`i`$, we can use in calculations the expression of for rotation of polarization plane of light in a weak magnetic field. If we take into account only dipole transitions (it is possible for example for $`6s_{1/2}7s_{1/2}`$ transition in cesium), then the angle of polarization plane rotation has the form $$\varphi =\frac{2\pi Nl}{(2i+1)(2j_g+1)}\frac{\omega }{\mathrm{\Delta }_D\mathrm{}c}\frac{E_z}{\mathrm{}\mathrm{\Delta }_D}\|A\|^2\left(\frac{g(u,v)}{u}\delta _1+2g(u,v)\gamma _1\right).$$ (27) where $`A`$ is the reduced matrix element of transition amplitude. The expressions for parameters $`\gamma _1`$ and $`\delta _1`$ are given below $`\gamma _1`$ $`=`$ $`{\displaystyle \frac{(2F_g+1)(2F_e+1)}{\sqrt{6}}}(1)^i\left\{\begin{array}{ccc}i& j_g& F_g\\ 1& F_e& j_e\end{array}\right\}[d_e(1)^{j_e+F_g}\sqrt{{\displaystyle \frac{(j_e+1)(2j_e+1)}{j_e}}}`$ $`({\displaystyle \frac{\mathrm{\Delta }_D}{\mathrm{\Delta }_{hf}(F_e,F_e1)}}(2F_e1)\left\{\begin{array}{ccc}i& j_g& F_g\\ 1& F_e1& j_e\end{array}\right\}\left\{\begin{array}{ccc}i& j_e& F_e\\ 1& F_e1& j_e\end{array}\right\}\left\{\begin{array}{ccc}F_g& 1& F_e\\ 1& F_e1& 1\end{array}\right\}`$ $`+{\displaystyle \frac{\mathrm{\Delta }_D}{\mathrm{\Delta }_{hf}(F_e,F_e+1)}}(2F_e+3)\left\{\begin{array}{ccc}i& j_g& F_g\\ 1& F_e+1& j_e\end{array}\right\}\left\{\begin{array}{ccc}i& j_e& F_e\\ 1& F_e+1& j_e\end{array}\right\}\left\{\begin{array}{ccc}F_g& 1& F_e\\ 1& F_e+1& 1\end{array}\right\})`$ $`(j_ej_g,F_eF_g,d_ed_g,)]`$ and $`\delta _1`$ $`=`$ $`{\displaystyle \frac{(2F_g+1)(2F_e+1)}{\sqrt{6}}}(1)^i\left\{\begin{array}{ccc}i& j_g& F_g\\ 1& F_e& j_e\end{array}\right\}^2[d_e(1)^{j_e+F_g}\sqrt{{\displaystyle \frac{(j_e+1)(2j_e+1)}{j_e}}}`$ $`(2F_e+1)\left\{\begin{array}{ccc}i& j_e& F_e\\ 1& F_e& j_e\end{array}\right\}\left\{\begin{array}{ccc}F_g& F_e& 1\\ 1& 1& F_e\end{array}\right\}+(j_ej_g,F_eF_g,d_ed_g,)]`$ Here $`\mathrm{\Delta }_{hf}`$ is the hyperfine level splitting, $`\mathrm{\Delta }_D`$ is the Doppler linewidth, $`j_g`$ and $`j_e`$ are the angular moments of electrons in atom, $`i`$ is the nuclear spin, $`F_g`$ and $`F_e`$ are total angular moments of atom in a ground and excited states respectively. First term in (27) arise from level splitting in electric field. It describes effect similar to Macaluso - Corbino rotation of polarization plane in magnetic field. Second term appears due to mixing of hyperfine levels with the different total moment $`F`$ but the same $`F_z`$ in electric field. It describes the $`T`$ noninvariant analog of polarization plane rotation due to Van-Vleck mechanism. ## 4 Estimates Let us compare the magnitude of $`PT`$ \- odd polarization plane rotation for different transitions. If spin of atomic nucleus is zero than the angle of rotation of polarization plane per absorption length due to $`PT`$ \- odd level mixing according to (18) has the form $$\varphi (L_{abs})=\frac{g(u,v)}{4f(u,v)}\frac{\|\overline{A^+}\|^2\|\overline{A^{}}\|^2}{\|\overline{A^\pm }\|^2}$$ (28) If detuning $`\mathrm{\Delta }\mathrm{\Delta }_D`$ then $`gf1`$. As in the case of $`P`$ odd but $`T`$ even interaction , we will discuss the rotation of polarization plane of light near magnetic dipole transitions. We can estimate the angle of polarization plane rotation per absorption length as follows. $$\varphi (L_{abs})\frac{(d^2H_T/\mathrm{\Delta }E)(dE_z/\mathrm{\Delta }E)}{\mu ^2}\frac{H_T}{\alpha ^2\mathrm{\Delta }E}\frac{dE_z}{\mathrm{\Delta }E}$$ (29) Here $`dea_0`$, $`\mu \alpha d`$ are the values of $`E1`$ and $`M1`$ transition amplitudes, $`H_T`$ is the matrix element of $`PT`$ noninvariant interaction, $`\mathrm{\Delta }ERy`$ is the typical space between energies of opposite parity states, $`a_0`$ is the Bohr radius, $`\alpha =1/137`$ is the fine structure constant. Near strongly forbidden magnetic dipole transitions (e. q. $`6s_{1/2}7s_{1/2}`$ in $`Cs`$) Stark-induced $`E1`$ amplitude is several orders of magnitude greater than $`M1`$ one. Absorption of light here depends primarily on Stark-induced amplitude. We can write for angle of polarization plane rotation per absorption length the following expression $$\varphi (L_{abs})\frac{(d^2H_T/\mathrm{\Delta }E)(dE_z/\mathrm{\Delta }E)}{(d^2E_z/\mathrm{\Delta }E)^2}\frac{H_T}{\mathrm{\Delta }E}\frac{\mathrm{\Delta }E}{dE_z}$$ (30) Usually $`dE_z/\mathrm{\Delta }E`$ is less than $`10^310^4`$. Therefore angle of rotation per absorption length is higher for strongly forbidden $`M1`$ transition. Angle of rotation per absorption length near allowed $`E1`$ transition is essentially lower than (29) and (30). $$\varphi (L_{abs})\frac{(d^2H_T/\mathrm{\Delta }E)(dE_z/\mathrm{\Delta }E)}{d^2}\frac{H_T}{\mathrm{\Delta }E}\frac{dE_z}{\mathrm{\Delta }E}$$ (31) It should be noted here that the angle of rotation of polarization plane per unit length has the same order of magnitude in all three cases. The difference in angle of rotation per absorption length is caused by different absorption of light near the corresponding transition. It is interesting to compare these estimates with the rotation of polarization plane caused by nonzero EDM of atom. In the absence of hyperfine structure the angle of rotation per absorption length caused by splitting of magnetic sub-levels in electric field can be estimated using (18) and (27) as follows $$\varphi (L_{abs})=\frac{1}{2(2j_g+1)}\frac{E_z\delta }{f\mathrm{}\mathrm{\Delta }_D}\frac{g}{u}\frac{d_{at}E_z}{\mathrm{}\mathrm{\Delta }_D}\frac{dE}{\mathrm{}\mathrm{\Delta }_D}\frac{H_T}{\mathrm{\Delta }E},$$ (32) where $`d_{at}dH_T/\mathrm{\Delta }E`$ is the EDM of atom, $`\mathrm{\Delta }_D(10^510^6)\mathrm{\Delta }E`$ is the Doppler linewidth. The value $`\varphi (L_{abs})`$ here does not depend on magnitude of transition amplitude $`A`$ and has the same order of magnitude for all kinds of transition considered above. ## 5 P and T - odd interactions in atom Several mechanisms can cause the $`PT`$ noninvariant interactions in atom. According to they include $`PT`$ \- odd weak interactions of electron and nucleon, the interaction of electric dipole moment of electron with the electric field inside the atom, interaction of electrons with electric dipole and magnetic quadrupole moments of nucleus and $`PT`$ odd electron - electron interaction. Below we will consider two kinds of $`PT`$ \- odd interactions that according to give the dominant contribution in our case. This is $`PT`$ \- odd electron nucleon interaction and interaction of electron EDM with the electric field inside atom. According to Hamiltonian of $`T`$ \- violating interaction between electron and hadron has the form $$H_T=C_s\frac{G}{\sqrt{2}}(\overline{e}i\gamma _5e)(\overline{n}n)+C_t\frac{G}{\sqrt{2}}(\overline{e}i\gamma _5\sigma _{\mu \nu }e)(\overline{n}\sigma ^{\mu \nu }n)$$ (33) where $`G=1.05510^5m_p^2`$ is Fermi constant, $`e`$ and $`n`$ are electron and hadron field operators respectively, $`C_s`$ and $`C_t`$ are dimensionless constants that characterize the strengths of $`T`$ \- violating interactions relative to usual $`T`$ \- conserving weak interaction. The first term in (33) describes scalar hadronic current coupling to pseudoscalar electronic current, and the second one describes tensor hadron current coupling to the pseudotensor electron current. Matrix elements for this $`T`$ \- odd Hamiltonians according to is equal to $$s_{1/2}H_{\mathrm{Todd}}p_{1/2}=\frac{Gm_e^2\alpha ^2Z^2R}{2\sqrt{2}\pi }\frac{\mathrm{𝐑𝐲}}{\sqrt{\nu _s\nu _p}^3}2\gamma C_sA$$ (34) where $`m_e`$ is the electron mass, $`\mathrm{Ry}=13.6\mathrm{eV}`$ is Rydberg energy constant, $`\nu _i`$ is the effective principal quantum number of state $`i`$, $`A`$ is the atomic number, $`R`$ is the relativistic factor ($`R=2.8`$ for cesium), $`\gamma =\sqrt{(j+1/2)^2Z^2\alpha ^2}`$ and $`j`$ is the total angular moment of atom. We neglect here tensor part of interaction for simplicity. Hamiltonian of interaction of electron EDM and electric field inside atom that mix opposite parity atomic states has the form $$H_d=\underset{k}{}(\gamma _{0k}1)\stackrel{}{\mathrm{\Sigma }}_k\stackrel{}{E}_k$$ (35) where $`E_k`$ is the electric field strength acting upon electron $`k`$. When summation in (35) is performed over one valence electron only and electric field strength near the nucleus approximately equals to $`\stackrel{}{E}=Z\alpha \stackrel{}{r}/r^3`$, matrix element of operator $`H_d`$ can be written as follows $$j,l=j+1/2\|\|H_d\|\|j,l^{}=j1/2=\frac{4(Z\alpha )^3}{\gamma (4\gamma ^21)(\nu _l\nu _l^{})^{3/2}a_0^2}$$ (36) where $`l`$ and $`l^{}`$ are the orbital angular moments, $`a_0`$ is the Bohr radius. ## 6 Estimates for $`6s_{1/2}7s_{1/2}`$ transition in cesium Let us estimate the $`PT`$ \- odd rotation of polarization plane for highly forbidden $`M1`$ transition $`6s_{1/2}7s_{1/2}`$ in cesium. The schemes of atomic levels of cesium is shown in Fig. 2. ### 6.1 Rotation of polarization plane of light due electron nucleon interactions In order to obtain the angle of polarization plane rotation for the highly forbidden $`M1`$ transition $`6s_{1/2}7s_{1/2}`$ due to $`P`$ and $`T`$ noninvariant interactions between electron and nucleus we can use the well known results for $`P`$ \- odd but $`T`$ \- even weak interactions. Matrix elements for $`T`$ \- even Hamiltonian $`H_w`$ has the form . $`s_{1/2}H_wp_{1/2}`$ $`=`$ $`i{\displaystyle \frac{Gm_e^2\alpha ^2Z^2R}{2\sqrt{2}\pi }}{\displaystyle \frac{\mathrm{𝐑𝐲}}{\sqrt{\nu _s\nu _p}^3}}Q_w`$ (37) $`Q_w=N+Z(14\mathrm{sin}\theta _W)`$ is the week nuclear charge, $`N`$ and $`Z`$ are the number of neutrons and protons in nucleus respectively, $`\mathrm{sin}^2\theta _W`$ is the Weinberg angle. Comparing this expressions with (34) and using the wave functions of $`6s_{1/2}`$ and $`7s_{1/2}`$ cesium states perturbed by $`P`$ \- odd but $`T`$ \- even interactions $`\|\stackrel{~}{6s_{1/2}}>`$ $`=`$ $`\|6s_{1/2}>+i10^{11}\left({\displaystyle \frac{Q_w}{N}}\right)(1.17\|6p_{1/2}>+0.34\|7p_{1/2}>)`$ $`\|\stackrel{~}{7s_{1/2}}>`$ $`=`$ $`\|7s_{1/2}>+i10^{11}\left({\displaystyle \frac{Q_w}{N}}\right)(0.87\|6p_{1/2}>1.33\|7p_{1/2}>)`$ (38) one can obtain atomic wave functions perturbed by $`P`$ and $`T`$ noninvariant interactions $`\|\stackrel{~}{6s_{1/2}}>`$ $`=`$ $`\|6s_{1/2}>+10^{11}\left(2\gamma {\displaystyle \frac{A}{N}}C_s\right)(1.17\|6p_{1/2}>+0.34\|7p_{1/2}>)`$ $`\|\stackrel{~}{7s_{1/2}}>`$ $`=`$ $`\|7s_{1/2}>+10^{11}\left(2\gamma {\displaystyle \frac{A}{N}}C_s\right)(0.87\|6p_{1/2}>1.33\|7p_{1/2}>).`$ (39) After simple calculations reduced matrix element of $`PT`$ odd $`E1`$ transition can be written as follows $`6s_{1/2}d^{PT}7s_{1/2}`$ $`=`$ $`\sqrt{3}10^{11}\left(2\gamma {\displaystyle \frac{A}{N}}C_s\right)(0.876s_{1/2}\|\|d\|\|6p_{1/2}`$ $`1.176p_{1/2}\|\|d\|\|7s_{1/2}0.347p_{1/2}\|\|d\|\|7s_{1/2})`$ Using values of radial integrals $`\rho (6s_{1/2},6p_{1/2})=5.535`$, $`\rho (7s_{1/2},6p_{1/2})=5.45`$, $`\rho (7s_{1/2},7p_{1/2})=12.30`$ we obtain $$6s_{1/2}d^{PT}7s_{1/2}=1.2710^{10}\|e\|a_0C_s$$ (40) The matrix element of Stark - induced $`6s_{1/2}7s_{1/2}`$ transition in cesium has traditionally been written in form $$6s_{1/2},m^{}\|d_i^{St}\|7s_{1/2},m=\alpha E_i\delta _{mm^{}}+i\beta ϵ_{ijk}E_jm^{}\|\sigma _k\|m$$ where $`m`$ and $`m^{}`$ are magnetic quantum numbers of ground and excited states of cesium, $`E_i`$ is the electric field strength, $`\sigma _k`$ is the Pauli matrix, $`\alpha `$ and $`\beta `$ are the scalar and vector transition polarizability (aee also (25)). The value of $`6s_{1/2}d^{St}7s_{1/2}`$ introduced in (26) can be expressed for cesium via the vector transition polarizability $`\beta `$ as follows. $`6s_{1/2}d^{St}7s_{1/2}=\sqrt{6}\beta E`$. Value of $`\beta `$ is well known from theoretical calculations as well as from experiment . According to it is equal to $`\beta =27.0a_0^3`$. Therefore $$7s_{1/2}d^{St}6s_{1/2}=1.2810^8\|e\|a_0E(V/cm)$$ (41) When temperature is $`T=750K`$ pressure of $`Cs`$ vapor is $`p=10`$ kPa , concentration of atoms is $`N=10^{18}cm^3`$ and Doppler linewidth is $`\mathrm{\Delta }_D/\omega _0=10^6`$. For transition between hyperfine levels $`F_g=4F_e=4`$, where coefficient $`K^2`$ is maximal ($`K^2=15/8`$) when detuning $`\mathrm{\Delta }\mathrm{\Delta }_D`$, $`v=\mathrm{\Gamma }/2\mathrm{\Delta }_D0.1`$ and $`f1`$, $`g0.7`$, absorption length in longitudinal electric field $`E`$ is equal to $`L_{abs}=710^{10}/E^2(V/cm)`$ (length is measured in centimeters) and angle of $`PT`$ noninvariant rotation of polarization plane is $$\varphi =1.010^{13}C_slE.$$ If $`E=10^4\mathrm{V}/\mathrm{cm}`$ then $`L_{abs}=7\mathrm{m}`$. The best signal to noise ratio is achieved when $`l=2L_{abs}`$ . In this case $`\|\varphi \|=1.310^6C_s`$. The lowest limit to the parameters of electron-nucleon interaction $`C_s<410^7`$ was set in . Corresponding limit to the angle of rotation of polarization plane is $`\|\varphi \|<0.510^{12}\mathrm{rad}.`$ #### 6.1.1 Cesium EDM Using wave functions (39) we can obtain EDM in $`6s_{1/2}`$ and $`7s_{1/2}`$ states of $`Cs`$ $$d_{6s_{1/2}}=1.3510^{10}C_s\|e\|a_0$$ $$d_{7s_{1/2}}=4.3910^{10}C_s\|e\|a_0.$$ Under the same condition of experiment as before expression (27) yields the angle of polarization plane rotation due to level splitting in electric field $$\|\varphi _1\|=1.410^{24}lC_sE_z^3(\mathrm{V}/\mathrm{cm})<810^{16}$$ and the angle of rotation due to hyperfine levels mixing $$\|\varphi _2\|=2.110^{24}lC_sE_z^3(\mathrm{V}/\mathrm{cm})<1.210^{15}.$$ (We assume here that for detuning $`\mathrm{\Delta }\mathrm{\Delta }_D`$ functions $`g(u,v)0.7`$, $`g(u,v)/u1.1`$). These angles are two orders of magnitude lower than one arising from interference of Stark-induced and $`PT`$ noninvariant transition amplitudes. ### 6.2 Rotation of polarization plane of light due to electron EDM If the $`PT`$ noninvariant interaction in atom is caused by interaction of electron EDM with the electric field of nucleus then wave functions of $`6s`$ and $`7s`$ states of cesium take the form $$\|\stackrel{~}{n}s_{1/2}>=\|ns_{1/2}>+\underset{m}{}\|mp_{1/2}>\frac{mp_{1/2}\|H_d\|ns_{1/2}}{E_gE_n}$$ (42) where matrix element of operator $`H_d`$ is given in (36). We should take into account admixture of states $`6p_{1/2}`$ and $`7p_{1/2}`$ to $`6s_{1/2}`$ and $`7s_{1/2}`$ states of cesium. After necessary calculations perturbed wave functions of atom can be written as follows $`\|\stackrel{~}{6}s_{1/2}>`$ $`=`$ $`\|6s_{1/2}>(35\|6p_{1/2}>+10.5\|7p_{1/2}>)d_e/(ea_0)`$ $`\|\stackrel{~}{7}s_{1/2}>`$ $`=`$ $`\|7s_{1/2}>+(27.7\|6p_{1/2}>36.2\|7p_{1/2}>)d_e/(ea_0)`$ (43) Using Eq. (43) we can obtain reduced matrix element of electric dipole transitions between $`\stackrel{~}{6}s_{1/2}`$ and $`\stackrel{~}{7}s_{1/2}`$ states $`6s_{1/2}d^{PT}7s_{1/2}`$ $`=`$ $`\sqrt{2/3}d_e(35\rho (6p_{1/2},7s_{1/2})+10.5\rho (7p_{1/2},7s_{1/2})`$ (44) $`+27.7\rho (6p_{1/2},6s_{1/2}))=72d_e`$ and electric dipole moment of cesium in ground state $`d_{6s_{1/2}}`$ and excited state $`d_{7s_{1/2}}`$ $`d_{6s_{1/2}}`$ $`=`$ $`131d_e`$ $`d_{7s_{1/2}}`$ $`=`$ $`400d_e,`$ (45) where $`d_e`$ is the electron EDM. As we mention before two effects induce $`T`$ \- noninvariant rotation of polarization plane in electric field. First of them is the interference of $`PT`$ \- odd and Stark induced transition amplitudes and second is the interaction of atomic EDM with electric field. After substitution of (44) to equation (26) one can obtain the angle of rotation due to interference of amplitudes $`\|\varphi \|<0.610^{12}`$ under the same experimental conditions as before. The rotation due to atomic EDM is a sum of two contributions. Using first term in Eq. (27) and Eq. (45) one can obtain for the rotation induced by splitting of magnetic sub-levels in electric field the angle $`\|\varphi _1\|<1.310^{15}`$. The mixing of hyperfine components (second term in equation (27)) gives the contribution $`\|\varphi _2\|<210^{15}`$. For estimates we use experimental limit on electron EDM from $`\|d_e\|<410^{27}\|e\|\mathrm{cm}`$. Limits on angles quoted above are close to those obtained for $`PT`$ \- odd electron nucleon interactions. ## 7 Discussion of experiment The simplest experimental scheme to observe the pseudo - Faraday rotation of polarization plane of light in electric field consist of cell with atomic gas placed in the electric field and sensitive polarimeter. In the case of large absorption length one can place this cell in resonator or delay line optical cavity to reduce the size of experimental setup (see e. q. ). Several schemes was proposed to increase the sensitivity of measurements. One of them is based on the nonlinear magneto - optic effect (NMOE) . Due to ultra - narrow width of dispersion like shaped Faraday rotation caused by NMOE the sensitivity of this experiments to the PT noninvariant rotation of polarization plane (the change of rotation angle with the change of applied electric field) is several orders of magnitude higher than in the conventional scheme. The authors of hopes to achieve the sensitivity to the cesium EDM $`d_{Cs}<10^{26}\|e\|`$ cm. The corresponding limit to the electron EDM is $`d_{Cs}<10^{28}\|e\|`$ cm. Even higher sensitivity can probably provide the method of measurements of polarization plane rotation proposed in . This method is based on observation of evolution of polarization of light in a a cell with atomic vapor and amplifying media with inverse population of atomic levels placed in the resonator. According to the compensation of absorption of light in the cell allows to increase the observed angle of polarization plane rotation. ## 8 Conclusion In the present article we have considered phenomenon of rotation of polarization plane of light in gas placed in electric field. Calculations of angle of polarization plane rotation were performed for $`6S_{1/2}7S_{1/2}`$ transition in atomic cesium. Two mechanisms that cause this effect are considered. They are interference of Stark-induced and $`PT`$ noninvariant transition amplitude and atomic EDM. Both of them can be induced by $`PT`$ noninvariant interaction between electrons and atomic nucleus and by interaction of electron EDM with electric field inside atom. For the highly forbidden $`M1`$ transition $`6S_{1/2}7S_{1/2}`$ in cesium we can expect the angle of polarization plane rotation per absorption length due to $`PT`$ \- odd atomic polarizability $`\beta _{mix}^{PT}`$ $`\varphi <10^{12}`$. This angle is three orders of magnitude greater than one caused by atomic EDM for this transition. Angle of polarization plane rotation can be significantly greater for other atoms, for example rare-earth elements, where additional amplification goes from close levels of opposite parity. The interesting example is the transition $`6s^2{}_{}{}^{1}S_{0}^{}6s5d{}_{}{}^{3}D_{1}^{}`$, ($`\lambda =408\mathrm{n}\mathrm{m}`$) in $`Yb`$ , where the value of $`PT`$ odd angle can be two orders of magnitude higher then for cesium due to larger $`PT`$ noninvariant amplitude. Therefore we can hope that experimental measurement of described phenomenon can achieve sensitivity in measurements of parameters of $`PT`$ noninvariant interactions between electron and nucleus and electron EDM, comparable with the current atomic EDM experiments.
warning/0002/hep-th0002219.html	ar5iv	text	# 1 Introduction. ## 1 Introduction. The non-abelian affine Toda (NAAT) equations are integrable (multi-component) generalizations of the sine-Gordon equation for a bosonic field that takes values in a non-abelian Lie group, in contrast with the usual Toda field theories where the field takes values in the (abelian) Cartan subgroup of a Lie group. In (see also ), it was found the subset of NAAT equations that can be written as the classical equations of motion of an action with a positive-definite kinetic term, a real potential term bounded from below and, moreover, with a mass-gap in order to make possible an $`S`$-matrix description. The resulting theories were named Homogeneous sine-Gordon (HSG) and Symmetric Space sine-Gordon (SSSG) theories, which are associated with the different compact Lie groups and compact symmetric spaces, respectively. Actually, the HSG and SSSG theories of are particular examples of the deformed coset models constructed by Park in and the Symmetric Space sine-Gordon models constructed by Bakas, Park and Shin in , respectively, where the specific form of the potential makes them exhibit a mass gap. There are some general features of these theories that have to be emphasized. The first one is that all of them have soliton solutions. Taking into account that both the HSG and SSSG theories are described by actions with sensible properties, this is in contrast with the usual affine Toda field theories where the condition of having soliton solutions (imaginary coupling constant) always leads to an ill defined action which makes the quantum theory problematic . The second is that they are defined by an action of the form $$S[h]=\frac{1}{\beta ^2}\left\{S_{WZNW}[h]d^2xV(h)\right\},$$ (1.1) where $`h=h(x,t)`$ is a field that takes values in a non-abelian compact group $`G`$, $`V(h)`$ is some potential function on the group manifold, and $`S_{WZNW}`$ is either the WZNW action for the group $`G`$ or the gauged WZNW action for a coset of the form $`G/H`$, where $`H`$ is an abelian subgroup of $`G`$ to be specified. Therefore, if the quantum theory is to be well defined then the coupling constant has to be quantized: $`\beta ^2=1/k`$, for some positive integer $`k`$ (see (3.9) for a more precise form of this quantization rule). Such a quantization does not occur in the sine-Gordon theory or the usual affine Toda theories because the field takes values in an abelian group in those cases. An important consequence of this is that, in the quantum theory, the $`\beta ^2`$ will not be a continuous coupling constant. However, the quantum theory will have other continuous coupling constants that appear in the potential and, in particular, determine the mass spectrum. A third feature which follows from (1.1) is that these theories are naturally described as perturbations of conformal field theory (CFT) coset models. Therefore, they will provide a Lagrangian formulation for some already known integrable perturbations of CFT’s and, furthermore, they will also lead us to discovering new ones. Finally, these theories are expected to exhibit realistic properties of quantum particles not captured by other integrable field theories; for instance, the presence of unstable particles. The main quantum properties of the HSG theories are quite well understood. They are integrable perturbations of the $`G`$-parafermion theories, which are coset CFT’s of the form $`G_k/H`$, where $`G`$ is a compact simple Lie group of rank $`r_g`$, $`HU(1)^{r_g}`$ is a maximal torus, and the ‘level’ $`k`$ is a positive integer. The perturbation is given by a spinless primary field of conformal dimension $`\mathrm{\Delta }=\overline{\mathrm{\Delta }}=h_g^{}/(k+h_g^{})`$, with $`h_g^{}`$ the dual Coxeter number of $`G`$. The simplest HSG theory is associated to $`G=SU(2)`$, whose equation of motion is the complex sine-Gordon equation . This theory corresponds to the perturbation of the usual $`\text{}_k`$-parafermions by the first thermal operator , whose exact factorizable scattering matrix is the minimal one associated to $`a_{k1}`$ . The quantum integrability of the HSG theories for arbitrary $`G`$ was established in , its soliton spectrum was obtained in , and a proposal for the exact factorizable $`S`$-matrices of the theories related to simply laced Lie groups $`G`$ has been made in . The main feature of those $`S`$-matrices is that they possess resonance poles which can be associated directly to the presence of unstable particles in the spectrum via the classical Lagrangian. Actually, to the best of our knowledge, these are the only known integrable quantum field theories that describe unstable particles. The $`S`$-matrices of have been probed in using the thermodynamic Bethe ansatz. In particular, this analysis confirms the expected value of the central charge of the unperturbed CFT for any simply laced $`G`$, and supports the interpretation of the resonance poles as a trace of the existence of unstable particles in the theory. These scattering matrices have been recently generalized in a Lie algebraic sense by Fring and Korf in . In contrast, the quantum properties of the generality of SSSG theories are not known, and the purpose of this paper is to partially fill this gap. Namely, we will find the class of perturbed CFT’s corresponding to the SSSG theories, investigate their quantum integrability, and discuss the general features of their spectrum of solitons. The SSSG theories are related to a compact symmetric space $`G/G_0`$, with a $`G_0`$-valued field, and describe perturbations of either the WZNW CFT corresponding to $`G_0`$ or a coset CFT of the form $`G_0/H`$, where $`HU(1)^p`$ is a torus of $`G_0`$, not necessarily maximal. The equations of motion of this kind of theories for more general choices of the normal subgroup $`H`$ were originally considered in the context of the, so called, reduced two-dimensional $`\sigma `$-models , although their Lagrangian formulation was not known until much later . The results of show that they fit quite naturally into the class of non-abelian affine Toda theories and, what is more important, that the condition of having a mass gap requires that $`H`$ is either trivial or abelian. The simplest SSSG theories are the ubiquitous sine-Gordon field theory, which corresponds to $`G/G_0=SU(2)/SO(2)`$, and the complex sine-Gordon theory, which is related this time to $`Sp(2)/U(2)`$ (recall that it is also the HSG theory associated to $`SU(2)`$). Actually, these two theories serve as paradigms of what can be expected in more complex situations. Another theories already discussed in the literature that belong to the class of SSSG theories are the integrable perturbations of the $`SU(2)_k`$ WZNW model and its $`\stackrel{~}{so(2)}`$ reduction constructed by Brazhnikov . Both of them are related to the symmetric space $`SU(3)/SO(3)`$ and, moreover, the second is identified with the perturbation of the usual $`\text{}_k`$-parafermions by the second thermal operator. The classification of the SSSG theories as perturbed CFT’s is achieved through the calculation of the central charge of the unperturbed CFT and the conformal dimension of the perturbation. Since the unperturbed CFT is always a coset CFT of the form $`G_0/H`$, the calculation of its central charge is straightforward. In contrast, the calculation of the conformal dimension of the perturbation requires the knowledge of the structure of the symmetric space. The symmetric space $`G/G_0`$ is associated with a Lie algebra decomposition $`g=g_0g_1`$ that satisfies the commutation relations $$[g_0,g_0]g_0,[g_0,g_1]g_1,[g_1,g_1]g_0,$$ (1.2) where $`g`$ and $`g_0`$ are the Lie algebras of $`G`$ and $`G_0`$, respectively. Then, the conformal properties of the perturbation depend on the structure of the representation of $`g_0`$ provided by $`[g_0,g_1]g_1`$. First of all, if the perturbation is to be given by a single primary field, then this representation has to be irreducible. This amounts to restrict the choice of $`G/G_0`$ to the, so called, ‘irreducible symmetric spaces’ , which have been completely classified by Cartan and are labelled by type I and type II. In Section 2, we summarize the main features of the SSSG theories. In Section 3, the conformal dimension of the perturbation corresponding to all the SSSG models related to type I symmetric spaces is calculated by making use of the relationship between the classification of type I symmetric spaces and the classification of the finite order automorphisms of complex Lie algebras. It is worth noticing that this analysis only depends on the structure of the representation of $`g_0`$ given by $`[g_0,g_1]g_1`$. Therefore, our results apply to any SSSG related to a type I symmetric space irrespectively of the choice of the normal subgroup $`H`$ that determines the coset $`G_0/H`$ and specifies the underlying CFT. For example, they provide the conformal dimension of the perturbation in the SSSG models constructed by Bakas, Park and Shin in , which include the generalized sine-Gordon models related to the NAAT equations based on $`sl(2)`$ embeddings constructed by Hollowood, Miramontes and Park in . The classical integrability of all these theories is a consequence of the relationship between their equations of motion and the NAAT equations, which ensure the existence of an infinite number of classically conserved quantities. Concerning quantum integrability, it can be established by invoking a well known result due to Parke, which affirms that the existence of two higher spin conserved quantities of different spin in a two-dimensional quantum field theory is enough to ensure that there is no particle production in scattering processes and that the $`S`$-matrix is factorizable . This way, the quantum integrability of the HSG theories was demonstrated in by checking that the classically conserved quantities of spin $`\pm 2`$ remain conserved in the quantum theory after an appropriate renormalization. The majority of SSSG theories also have classically conserved quantities of spin $`\pm 2`$, and we expect that the analysis in can be generalized to those cases without much effort. However, there is a class of SSSG where the simplest higher spin classically conserved quantities are of spin $`\pm 3`$: the theories related to symmetric spaces of maximal rank, i.e., those which satisfy $`\mathrm{rank}(G/G_0)=\mathrm{rank}(G)`$, where the rank of the symmetric space is the dimension of the maximal abelian subspaces contained in $`g_1`$. In Section 4, we explicitly check that at least two of those classically conserved quantities of spin $`\pm 3`$ remain conserved in the quantum theory for all the theories related to a symmetric space of maximal rank where $`H`$ is trivial. We name these theories ‘Split models’, and they are the only SSSG theories which correspond to massive quantum integrable perturbations of WZNW theories. It is worth noticing that (marginal) perturbations of WZNW models have been recently considered in the context of $`\mathrm{AdS}_3`$ black holes . The results of Section 4, together with those of , support the conjecture that all the SSSG are quantum integrable. This implies that they should admit a factorizable $`S`$-matrix and the next stage of analysis consists in establishing its form. Since we expect that it should be possible to infer the form of the exact $`S`$-matrix through the semiclassical quantization of the solitons, in Section 5 we investigate the general features of the spectrum of solitons. In view of the different kinds of SSSG theories, we have restricted the analysis to the Split models, which illustrate the main properties to be expected. Like the sine-Gordon and complex sine-Gordon theories, the fundamental particles of the theory can be identified with some of the classical soliton solutions. Moreover, the solitons of the Split models are topological and do not carry Noether charges, again like the solitons of the sine-Gordon theory. The origin of the topological charge is both the existence of different vacua and the fact that $`G_0`$ can be non simply connected. This is in contrast with the solitons of the HSG theories which are not topological but carry a $`U(1)^{r_g}`$ abelian Noether charge . In general, the solitons of a generic SSSG theory corresponding to a perturbation of a coset CFT of the form $`G/H`$ are expected to carry both topological charges and abelian Noether charges related to a global symmetry of the classical action specified by $`H`$. In this sense, they are analogous to the dyons in four-dimensional non-abelian gauge theories . Another relevant feature of the solitons of the Split models, which is expected to be shared with other SSSG theories, is that their mass spectrum suggests that some of them might describe unstable particles in the quantum theory, which is analogous to what happens in the HSG theories . Actually, the unstability of the heaviest solitons in the spectrum has been checked by Brazhnikov in the SSSG theories related to $`SU(3)/SO(3)`$. Our conclusions are presented in Section 6, and we have collected the explicit expressions for the classically conserved densities of the Split models together with some useful algebraic notation in the Appendix. ## 2 The Symmetric Space sine-Gordon theories. The different SSSG models of are characterized by four algebraic data: $`\{g,\sigma ,\mathrm{\Lambda }_\pm \}`$. $`g`$ is a compact semisimple finite Lie algebra, and $`\sigma `$ is an involutive ($`\sigma ^2=1`$) automorphism of $`g`$ that induces the decomposition $$g=g_0g_1,$$ (2.1) with $`g_0`$ the set of fixed points of $`\sigma `$ and $`g_1=\{ug\sigma (u)=u\}`$. The subspaces $`g_0`$ and $`g_1`$ satisfy the commutation relations (1.2), which exhibits that $`G/G_0`$ is a compact symmetric space, where $`G`$ and $`G_0`$ are the Lie groups corresponding to the compact Lie algebras $`g`$ and $`g_0`$, respectively. Actually, the different choices for $`\{g,\sigma \}`$ are in one-to-one relation with the different compact symmetric spaces $`G/G_0`$, which prompted the name chosen in for this class of models. Finally, $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$, are two semisimple elements in $`g_1`$ whose choice is only restricted by the condition that $$\mathrm{Ker}(\mathrm{ad}_{\mathrm{\Lambda }_+})g_0=\mathrm{Ker}(\mathrm{ad}_\mathrm{\Lambda }_{})g_0=g_0^0$$ (2.2) is abelian, which is required to ensure the existence of a mass gap . They play the role of continuous coupling constants. The SSSG model associated to $`\{g,\sigma ,\mathrm{\Lambda }_\pm \}`$ or, equivalently, to $`\mathrm{\Lambda }_\pm `$ and the symmetric space $`G/G_0`$ is specified by the action $$S_{\mathrm{SSSG}}=\frac{1}{\beta ^2}\left\{S_{\mathrm{WZNW}}[h]+\frac{m^2}{\pi }d^2x\mathrm{\Lambda }_+,h^{}\mathrm{\Lambda }_{}h\right\},$$ (2.3) where $`h`$ is a bosonic field taking values in $`G_0`$, $`S_{\mathrm{WZNW}}`$ is the gauged WZNW action corresponding to the coset $`G_0/H`$, with $`H`$ the abelian Lie group corresponding to $`g_0^0`$, and $`m`$ is the only mass scale of the theory . The SSSG action (2.3) is invariant with respect to the abelian gauge transformations $$he^\alpha he^{\tau (\alpha )}$$ (2.4) for any $`\alpha =\alpha (x,t)g_0^0`$. Since the potential $`V(h)=m^2\mathrm{\Lambda }_+,h^{}\mathrm{\Lambda }_{}h/\pi `$ introduced in (2.3) has $`H\times H`$ left-right symmetry, this gauge symmetry is essential to make the SSSG theories exhibit a mass gap. The precise form of the group of gauge transformations is specified by $`\tau `$, which is an automorphism of $`g_0^0`$ that will not play any role in the following sections; we refer the reader to for further details about the conditions to be satisfied by $`\tau `$. The gauge symmetry leaves a residual global $`H=U(1)^p`$-symmetry, with $`p=\mathrm{dim}(g_0^0)`$, which, in particular, makes the classical solitonic solutions carry conserved abelian Noether charges, a feature that is shared with the HSG theories. Moreover, as will be discussed in Section 5, the solitonic solutions of the SSSG theories may also carry topological charges. This possibility does not exist in the case of the HSG theories where all the possible vacuum configurations get identified modulo gauge transformations. The action (2.3) can be obtained by Hamiltonian reduction of a gauged two loop WZNW model associated to the affine Kac-Moody algebra $`\overline{g}^{(r)}`$, where $`\overline{g}`$ is the complexification of $`g`$ and $`r=1,2,`$ or $`3`$ is the least positive integer for which $`\sigma ^r`$ is an inner automorphism (see the comments below (4.3)). In fact, the Hamiltonian reduction approach has been recently followed by Gomes et al. to construct another class of (classical) affine non-abelian Toda models different to the SSSG theories . The classical integrability of these theories is a consequence of the connection between their equations of motion and the non-abelian affine Toda equations . This is made explicit by considering a particular gauge fixing prescription where the classical equations of motion reduce to $`_{}(h^{}_+h)=m^2[\mathrm{\Lambda }_+,h^{}\mathrm{\Lambda }_{}h],`$ (2.5) $`P(h^{}_+h)=P(_{}hh^{})=\mathrm{\hspace{0.17em}0},`$ (2.6) where $`P`$ is a projector onto the subalgebra $`g_0^0`$, and $`x_\pm =t\pm x`$ are the light-cone variables. The first equation in (2.6) is a non-abelian affine Toda equation, and the second provide a set of constraints which come from the variation of the action (2.3) with respect to the abelian gauge connections in the LS gauge . Notice that the equations (2.6) are left invariant by the transformation $$xx,hh^{},\mathrm{\Lambda }_\pm \eta ^{\pm 1}\mathrm{\Lambda }_{},$$ (2.7) for an arbitrary real number $`\eta `$, which shows that the SSSG theories are parity invariant only if $`\mathrm{\Lambda }_+=\eta \mathrm{\Lambda }_{}`$ for some real number $`\eta `$, i.e., if $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$ are chosen to be parallel or anti-parallel . At the quantum level, the SSSG theories will be described as perturbed CFT’s of the form $$S=S_{\mathrm{CFT}}+\frac{m^2}{\pi \beta ^2}d^2x\mathrm{\Phi }(x,t).$$ (2.8) If, according to (2.2), $`g_0^0=u(1)^p`$ with $`p0`$, $`S_{\mathrm{CFT}}`$ will be the action of either the CFT associated to the coset $`G_0/U(1)^p`$ ($`p0`$) or the WZNW model corresponding to $`G_0`$ ($`p=0`$). Moreover, the perturbation is given by $`\mathrm{\Phi }=\mathrm{\Lambda }_+,h^{}\mathrm{\Lambda }_{}h`$, which will be understood as a matrix element of the WZNW field taken in the representation of $`G_0`$ provided by $`[g_0,g_1]g_1`$. As shown originally by Bakas for the complex sine-Gordon theory , and used in to define the HSG theories at the quantum level, these identifications constitute a non-perturbative definition of the SSSG theories. For a given symmetric space $`G/G_0`$, it is important to emphasise that the form of the coset $`G_0/U(1)^p`$ is fixed by the choice of $`\mathrm{\Lambda }_\pm `$. Since $`\mathrm{\Lambda }_\pm `$ are semisimple elements of $`g`$, different choices will lead to different values of $`p`$ in the range $$0\mathrm{rank}(G)\mathrm{rank}(G/G_0)p\mathrm{min}[\mathrm{rank}(G_0),\mathrm{rank}(G)\nu ],$$ (2.9) where the rank of the symmetric space, $`\mathrm{rank}(G/G_0)`$, is the dimension of the maximal abelian subspaces contained in $`g_1`$, and $`\nu =2`$ or 1 depending on whether $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$ are linearly independent or not, respectively. In particular, the lower bound is reached when $`\mathrm{\Lambda }_+,\mathrm{\Lambda }_{}g_1`$ are regular and, hence, $`\mathrm{Ker}(\mathrm{ad}_{\mathrm{\Lambda }_\pm })`$ is already a maximal abelian subspace of $`g`$. All this implies that the SSSG theories provide a rich variety of different integrable models that include, for $`p=\mathrm{rank}(G_0)`$, new massive perturbations of the theory of $`G_0`$-parafermions different than those provided by the Homogeneous sine-Gordon theories . Notice that this case happens only if the symmetric space satisfies $`\mathrm{rank}(G_0)\mathrm{rank}(G)\nu `$. Another particularly interesting class of models occurs when $`\mathrm{rank}(G)=\mathrm{rank}(G/G_0)`$ and $`p=0`$. In this case, the SSSG theory is just a massive perturbation of the WZNW model corresponding to $`G_0`$. We will concentrate on the theories where the perturbation in (2.8) is given by a single spinless primary field of the CFT which, taking into account the properties of the WZNW field , amounts to restrict ourselves to the SSSG theories associated with symmetric spaces where the representation of $`g_0`$ provided by $`[g_0,g_1]g_1`$ is irreducible. Otherwise, the perturbation will be the sum of more than one primary field. The symmetric spaces with that property are called ‘irreducible’, and have been completely classified by Cartan . There are two types of compact irreducible symmetric spaces: * Type I, where the compact Lie algebra $`g`$ is simple. * Type II, where the compact Lie algebra is of the form $`g=g_1g_2`$ with $`g_1=g_2`$ simple, and the involution $`\sigma `$ interchanges $`g_1`$ and $`g_2`$. In the following we will only consider the SSSG theories associated to the type I symmetric spaces, which admit a thorough classification. ## 3 The type I SSSG theories as perturbed CFT’s. In this section, we calculate the central charge of the unperturbed CFT and the conformal dimension of the perturbation corresponding to the SSSG theories associated with the symmetric spaces of type I. We will make use of the relationship between the Cartan classification of this type of symmetric spaces and the Kac classification of the automorphisms of finite order of complex Lie algebras, which provides a systematic and very convenient description of the involutive automorphism $`\sigma `$. This represents an important advantage with respect to previous works on integrable systems associated with symmetric spaces which generally make use of explicit parametrizations of the field $`h`$ based on some matrix representation for the symmetric space. ### 3.1 Type I symmetric spaces and finite order automorphisms. The basic result is due to Kac, who established the following correspondence between the involutions of a complex Lie algebra $`\overline{g}`$ and the involutions of its compact real form $`g`$. ###### Theorem 1 (Proposition 1.4 in , Ch. X) Let $`Aut(g)`$ denote the set of automorphisms of $`g`$, $`Inv(g)`$ the subset containing the involutions, and $`Inv(g)/Aut(g)`$ the set of conjugacy classes in $`Aut(g)`$ of the elements in $`Inv(g)`$. We define $`Inv(\overline{g})/Aut(\overline{g})`$ similarly. Each automorphism $`\sigma Inv(g)`$ extends uniquely to $`\overline{\sigma }Inv(\overline{g})`$ and if $`\sigma _1`$, $`\sigma _2`$ are conjugate within $`Aut(g)`$, then $`\overline{\sigma }_1`$, $`\overline{\sigma }_2`$ are conjugate within $`Aut(\overline{g})`$. Taking into account all this, it can be proved that the mapping: $$\tau :Inv(g)/Aut(g)Inv(\overline{g})/Aut(\overline{g}),$$ (3.1) induced by $`\sigma \overline{\sigma }`$ is a bijection. Recall now that the different compact symmetric spaces of type I associated with a compact simple Lie algebra $`g`$ are in one-to-one relation with the different involutive automorphisms of $`g`$, not distinguishing automorphisms which are conjugate by the group $`Aut(g)`$. Therefore, they are also in one-to-one relation with the involutive automorphisms of its complexification $`\overline{g}`$, modulo conjugations by $`Aut(\overline{g})`$. In order to summarize the Kac classification of the automorphisms of finite order of complex Lie algebras, it is necessary to introduce the following notation (see , Ch. 8, for more details). Let $`\overline{g}`$ denote a complex simple Lie algebra and $`\mu `$ an automorphism of $`\overline{g}`$ induced by an automorphism of its Dynkin diagram of order $`r=1,2`$ or $`3`$. $`\mu `$ induces a $`\text{}/r\text{}`$–gradation of $`\overline{g}`$, which means that $`\overline{g}`$ can be decomposed into the sum of a set of subspaces labelled by an integer $`0kr1`$ that satisfy $`\mu (u)=\mathrm{e}^{\frac{2\pi i}{r}j}uu\overline{g}_j(\mu ),`$ $`\overline{g}={\displaystyle \underset{k=1}{\overset{r1}{}}}\overline{g}_k(\mu ),[\overline{g}_j(\mu ),\overline{g}_k(\mu )]\overline{g}_{j+k\mathrm{mod}r}(\mu ).`$ (3.2) Then there is a particular set of generators of $`\overline{g}`$, $`\{E_0,E_1,\mathrm{},E_l\}`$, where $`l`$ is the rank of the invariant subalgebra $`\overline{g}_0(\mu )`$, with the following properties: * $`\{E_1,\mathrm{},E_l\}\overline{g}_0(\mu )`$ for $`(\overline{g},r)(A_{2l},2)`$, and $`\{E_0,E_1,\mathrm{},E_{l1}\}\overline{g}_0(\mu )`$ for $`(\overline{g},r)=(A_{2l},2)`$. * If $`(\overline{g},r)(A_{2l},2)`$, then $`E_0\overline{g}_1(\mu )`$ is the lowest–weight vector of the irreducible representation of $`\overline{g}_0(\mu )`$ given by $`[\overline{g}_0(\mu ),\overline{g}_1(\mu )]\overline{g}_1(\mu )`$. Conversely, when $`(\overline{g},r)=(A_{2l},2)`$ this role is played by $`E_l`$. * $`\{E_1,\mathrm{},E_l\}`$ are positive Chevalley generators for $`\overline{g}_0(\mu )`$, except for $`(\overline{g},r)=(A_{2l},2)`$ where the Chevalley generators are $`\{E_0,E_1,\mathrm{},E_{l1}\}`$. * Results a), b), and c) correspond to the case $`r>1`$. When $`r=1`$, $`l=\mathrm{rank}(\overline{g})`$, $`E_1=E_{\alpha _1},\mathrm{},E_l=E_{\alpha _l}`$, and $`E_0=E_{\mathrm{\Psi }_g}`$, where $`\{\stackrel{}{\alpha }_1,\stackrel{}{\alpha }_2,\mathrm{},\stackrel{}{\alpha }_l\}`$ is a set of simple roots of $`\overline{g}`$, and $`\stackrel{}{\mathrm{\Psi }}_g`$ is the highest root. The classification is given by the following theorem. ###### Theorem 2 (Theorem 8.6 in ) Let $`\stackrel{}{s}=(s_0,s_1,\mathrm{},s_l)`$ be a sequence of non-negative relatively prime integers, and put $$m=r\underset{i=0}{\overset{l}{}}a_is_i$$ (3.3) where $`a_0,a_1,\mathrm{},a_l`$ are the Kac labels corresponding to the Dynkin diagram of the (twisted if $`r1`$) affine Kac-Moody algebra $`\overline{g}^{(r)}`$. Then: * The relations $$\sigma _{\stackrel{}{s};r}(E_j)=e^{2\pi is_j/m}E_j,j=0,\mathrm{},l,$$ (3.4) define uniquely an $`m`$-th order automorphism $`(\stackrel{}{s};r)`$ of $`\overline{g}`$. * Up to conjugacy by an automorphism of $`\overline{g}`$, the automorphisms $`\sigma _{\stackrel{}{s};r}`$ exhaust all $`m`$-th order automorphisms of $`\overline{g}`$. * The elements $`\sigma _{\stackrel{}{s};r}`$ and $`\sigma _{\stackrel{}{s}^{};r^{}}`$ are conjugate by an automorphism of $`\overline{g}`$ if, and only if, $`r=r^{}`$ and the sequence $`\stackrel{}{s}`$ can be transformed into the sequence $`\stackrel{}{s}^{}`$ by an automorphism of the Dynkin diagram of $`\overline{g}^{(r)}`$. Taking into account Theorems 1 and 2, the classification of the symmetric spaces of type I is equivalent to working out the equation $$m=r\underset{i=0}{\overset{l}{}}a_is_i=\mathrm{\hspace{0.17em}2},$$ (3.5) which has only three possible types of solutions: $`[\mathrm{A1}]r=1,a_{i_0}=2,s_{i_0}=1,\mathrm{and}s_i=0\mathrm{for}ii_0,`$ $`[\mathrm{A2}]r=2,a_{i_0}=1,s_{i_0}=1,\mathrm{and}s_i=0\mathrm{for}ii_0,`$ $`[B]r=1,a_{i_0}=a_{i_1}=1,s_{i_0}=s_{i_1}=1,\mathrm{and}s_i=0\mathrm{for}ii_0,i_1.`$ (3.6) This classifies the type I (compact) symmetric spaces and, hence, the corresponding SSSG theories, into type A1, A2, and B . Given an automorphism $`\sigma _{\stackrel{}{s};r}`$, the subspaces $`\overline{g}_0`$ and $`\overline{g}_1`$ can be easily characterized as follows. ###### Theorem 3 (Proposition 8.6 in ) * Let $`i_1,i_2,\mathrm{},i_p`$ be all the indices for which $`s_{i_1}=\mathrm{}=s_{i_p}=0`$. Then the Lie algebra $`\overline{g}_0`$ is isomorphic to a direct sum of the (l-p)–dimensional centre and a semisimple Lie algebra whose Dynkin diagram is the subdiagram of the Dynkin diagram of $`\overline{g}^{(r)}`$ consisting of the vertices $`i_1,\mathrm{},i_p`$. * Let $`j_1,\mathrm{},j_n`$ be all the indices for which $`s_{j_1}=\mathrm{}=s_{j_n}=1`$. Then the representation of $`\overline{g}_0`$ provided by $`[\overline{g}_0,\overline{g}_1]\overline{g}_1`$ is isomorphic to a direct sum of $`n`$ irreducible modules with highest weights $`\stackrel{}{\alpha }_{j_1},\mathrm{},\stackrel{}{\alpha }_{j_n}`$. Eq. (3.6) and Theorem 3 imply that $`g_0`$ is of the form $$g_0=\{\begin{array}{cc}_{i=1}^qg^{(i)}\hfill & \text{for type A1 and A2,}\hfill \\ \multicolumn{2}{c}{}\\ _{i=1}^qg^{(i)}u(1)\hfill & \text{for type B,}\hfill \end{array}$$ (3.7) where $`q`$ can be either 1 or 2 in both cases, and $`g^{(i)}`$ is always compact and simple. Therefore, if the symmetric space is of type A1 or A2, $`S_{\mathrm{CFT}}`$ in (2.8) is the action of the CFT associated to a coset of the form $$\underset{i=1}{\overset{q}{}}G_{k_i}^{(i)}/U(1)^p,$$ (3.8) where $`G^{(i)}`$ is the compact simple Lie group corresponding to $`g^{(i)}`$, and $`p`$ is some integer in the range (2.9). Since $`g_0`$ is non-abelian, the consistency of the quantum theory requires that the coupling constant in (2.3) is quantized. The precise form of the quantization rule is $$\frac{1}{\mathrm{}\beta ^2}=\frac{\stackrel{}{\mathrm{\Psi }}_{g^{(i)}}^2}{2}k_i,$$ (3.9) where $`k_i`$ is an integer for each simple factor in (3.7), and $`\stackrel{}{\mathrm{\Psi }}_{g^{(i)}}^2`$ is the square length of the long roots of $`g^{(i)}`$ with respect to the bilinear form of $`g`$. In (3.9) we have shown explicitly the Plank constant to exhibit that, just as in the sine-Gordon theory, the semi-classical limit is the same as the weak coupling limit, and that both are recovered when $`k_i\mathrm{}`$. Since there is a unique coupling constant $`\beta ^2`$, eq. (3.9) implies that the ‘levels’ in (3.8) are related by means of $$\frac{k_i}{k_j}=\frac{\stackrel{}{\mathrm{\Psi }}_{g^{(j)}}^2}{\stackrel{}{\mathrm{\Psi }}_{g^{(i)}}^2}.$$ (3.10) Nevertheless, our calculations show that $`\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}^2=\stackrel{}{\mathrm{\Psi }}_{g^{(2)}}^2`$ for all the type I symmetric spaces, with the only exception of $`G_2/SU(2)\times SU(2)`$ where $`\stackrel{}{\mathrm{\Psi }}_{g^{(2)}}^2=3\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}^2`$ (see Section 3.2.2), and we will use the notation $`k=k_1`$ in the following. Therefore, in this case, the central charge of the unperturbed CFT is $$c_{\mathrm{CFT}}=\underset{i=1}{\overset{q}{}}\frac{k_i\mathrm{dim}(g^{(i)})}{k_i+h_i^{}}p,$$ (3.11) where $`h_i^{}`$ is the dual Coxeter number of $`g^{(i)}`$. Since the perturbation $`\mathrm{\Phi }`$ is just a matrix element of the WZNW field in the representation of $`G_0`$ provided by $`[g_0,g_1]g_1`$, which is irreducible, $`\mathrm{\Phi }`$ is a spinless primary field with conformal dimension $$\mathrm{\Delta }_\mathrm{\Phi }=\overline{\mathrm{\Delta }}_\mathrm{\Phi }=\underset{i=1}{\overset{q}{}}\frac{C_2(g^{(i)})/\stackrel{}{\mathrm{\Psi }}_{g^{(i)}}^2}{k_i+h_i^{}},$$ (3.12) where $`C_2(g^{(i)})`$ is the quadratic Casimir. Then, Theorem 3 shows that this representation is a highest weight representation and, hence, the quadratic Casimir is given by $$C_2(g^{(i)})=\stackrel{}{\mathrm{\Lambda }},\stackrel{}{\mathrm{\Lambda }}+\mathrm{\hspace{0.17em}2}\stackrel{}{\delta }^{(i)}$$ (3.13) where $`\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{\alpha }_{i_0}`$ is the highest weight, and $`\stackrel{}{\delta }^{(i)}`$ is half the sum of the positive roots of $`\overline{g}^{(i)}`$. Let us consider now the SSSG theories associated with symmetric spaces of type B. The main difference with the SSSG models of type A1 or A2 is that $`g_0`$ includes now a one-dimensional centre: the $`u(1)`$ factor in (3.7). In this case, it will be convenient to choose the Cartan subalgebra of $`g`$ such that it contains the Cartan subalgebras of $`_{i=1}^qg^{(i)}`$ in addition to the generator of the centre. Then, $`u(1)=\text{}i\stackrel{}{u}\stackrel{}{h}`$, where $`\stackrel{}{u}`$ is a vector that is orthogonal to all the roots of $`_{i=1}^qg^{(i)}`$, and the components of $`i\stackrel{}{h}`$ provide a basis for the Cartan subalgebra of $`g`$. According to 3, another important difference is that the representation of $`\overline{g}_0`$ given by $`[\overline{g}_0,\overline{g}_1]\overline{g}_1`$ is the sum of two irreducible highest weight representations with highest weights $`\stackrel{}{\mathrm{\Lambda }}_1=\stackrel{}{\alpha }_{i_0}`$ and $`\stackrel{}{\mathrm{\Lambda }}_2=\stackrel{}{\alpha }_{i_1}`$, namely, $$\overline{g}_1=L(\stackrel{}{\mathrm{\Lambda }}_1)L(\stackrel{}{\mathrm{\Lambda }}_2).$$ (3.14) Let us consider the identity $$\underset{i=0}{\overset{l}{}}a_i\stackrel{}{\alpha }_i=\mathrm{\hspace{0.25em}0}$$ (3.15) satisfied by the Kac labels of the Dynkin diagram of $`\overline{g}^{(1)}`$, where $`\{\stackrel{}{\alpha }_1,\stackrel{}{\alpha }_2,\mathrm{},\stackrel{}{\alpha }_l\}`$ is a set of simple roots of $`\overline{g}`$, and $`\stackrel{}{\mathrm{\Psi }}_g`$ is the highest root. Then, the two highest weights satisfy $$\stackrel{}{\mathrm{\Lambda }}_1+\stackrel{}{\mathrm{\Lambda }}_2=\underset{\genfrac{}{}{0pt}{}{i=0}{ii_0,i_1}}{\overset{l}{}}a_i\stackrel{}{\alpha }_i,$$ (3.16) which implies that $`\stackrel{}{\mathrm{\Lambda }}_1`$ is minus the lowest weight of the representation with highest weight $`\stackrel{}{\mathrm{\Lambda }}_2`$. This manifests that the two highest weight representations are conjugate, $`L(\mathrm{\Lambda }_2)=L^{}(\mathrm{\Lambda }_1)`$, which is consistent with the fact that the representation of $`g_0`$ given by $`[g_0,g_1]g_1`$ is irreducible. Therefore, since $`\mathrm{\Lambda }_\pm g_1`$, it can be decomposed as $$\mathrm{\Lambda }_\pm =\lambda _\pm +\lambda _\pm ^{},\mathrm{with}\lambda _\pm L(\mathrm{\Lambda }_1).$$ (3.17) Taking again into account Theorem 3, eq. (3.16) also implies that $`\stackrel{}{\mathrm{\Lambda }}_1+\stackrel{}{\mathrm{\Lambda }}_2`$ is a linear combination of the roots of $`g_0`$, which are orthogonal to $`\stackrel{}{u}`$, and, hence, $$\stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_1=\stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_2.$$ (3.18) Our calculations show that, in all the symmetric spaces of type B, $`\stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_10`$. Since $`\mathrm{\Lambda }_\pm g_1`$, this implies that $`[\stackrel{}{u}\stackrel{}{h},\mathrm{\Lambda }_\pm ]0`$ and, hence, the $`u(1)`$ is not in $`g_0^0`$ for any choice of $`\mathrm{\Lambda }_\pm `$ (see (2.2)). Consider a generic field configuration $$h=\stackrel{~}{h}\mathrm{exp}\left(i\phi \stackrel{}{u}\stackrel{}{h}\right),$$ (3.19) where $`\stackrel{~}{h}`$ is a field taking values in the compact semisimple Lie group $`_{i=1}^qG^{(i)}`$, and $`\phi =\phi (x,t)`$ is a real scalar field; for convenience we will normalize $`\stackrel{}{u}`$ such that $`\stackrel{}{u}\stackrel{}{u}=4\pi `$. Then, the action (2.3) becomes $`S_{\mathrm{SSSG}}={\displaystyle \frac{1}{\beta ^2}}\{S_{\mathrm{WZNW}}[\stackrel{~}{h}]+{\displaystyle \frac{1}{2}}{\displaystyle }d^2x_\mu \phi ^\mu \phi `$ $`+{\displaystyle \frac{m^2}{\pi }}{\displaystyle }d^2x(\mathrm{e}^{i\phi \stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_1}\lambda _+^{},\stackrel{~}{h}^{}\lambda _{}\stackrel{~}{h}+\mathrm{e}^{+i\phi \stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_1}\lambda _+,\stackrel{~}{h}^{}\lambda _{}^{}\stackrel{~}{h})\}.`$ (3.20) Therefore, if the symmetric space is of type B, $`S_{\mathrm{CFT}}`$ in (2.8) is the action of the CFT associated to a coset of the form $$\left[\underset{i=1}{\overset{q}{}}G_{k_i}^{(i)}/U(1)^p\right]\times U(1),$$ (3.21) i.e., a coset of the form (3.8) plus a massless boson, whose central charge is $$c_{\mathrm{CFT}}=\underset{i=1}{\overset{q}{}}\frac{k_i\mathrm{dim}(g^{(i)})}{k_i+h_i^{}}+\mathrm{\hspace{0.25em}1}p,$$ (3.22) with the levels $`k_i`$ defined by the quantization rule (3.9). In (3.21), we have already taken into account that the centre of $`g_0`$ is not in $`g_0^0`$ as a consequence of (3.17) and (3.18). Therefore, in this case, $`\mathrm{rank}(G_0)`$ has to be substituted for $`\mathrm{rank}(G_0)1`$ on the right-hand-side of eq. (2.9). Concerning the perturbation $`\mathrm{\Phi }`$ in (3.20), it is a primary field of conformal dimension $$\mathrm{\Delta }_\mathrm{\Phi }=\overline{\mathrm{\Delta }}_\mathrm{\Phi }=\underset{i=1}{\overset{q}{}}\frac{C_2(g^{(i)})/\stackrel{}{\mathrm{\Psi }}_{g^{(i)}}^2}{k_i+h_i^{}}+\frac{(\stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_1)^2}{k\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}^2},$$ (3.23) where the quadratic Casimir is given by (3.13) with $`\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{\mathrm{\Lambda }}_1`$ or $`\stackrel{}{\mathrm{\Lambda }}_2`$, both leading to identical results. At this stage, we would like to correct a wrong statement in concerning the fields associated to the centre of $`g_0`$. In that article, it was said that those fields can always be decoupled whilst preserving integrability. However, the SSSG of type B shows that this is not true in general. In our case, there is only one field associated to the centre of $`g_0`$, $`\phi =\phi (x,t)`$, whose classical equation of motion is $$_+_{}\phi =im^2\frac{\stackrel{}{u}\stackrel{}{\mathrm{\Lambda }_1}}{4\pi }\left(\mathrm{e}^{i\phi \stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_1}\lambda _+^{},\stackrel{~}{h}^{}\lambda _{}\stackrel{~}{h}\mathrm{e}^{+i\phi \stackrel{}{u}\stackrel{}{\mathrm{\Lambda }}_1}\lambda _+,\stackrel{~}{h}^{}\lambda _{}^{}\stackrel{~}{h}\right),$$ (3.24) which clearly shows that $`\phi `$ cannot be decoupled simply by putting $`\phi (x,t)=0`$ unless $`\lambda _+^{},\stackrel{~}{h}^{}\lambda _{}\stackrel{~}{h}`$ is real, which is equivalent to $$[\mathrm{\Lambda }_+,\stackrel{~}{h}^{}\mathrm{\Lambda }_{}\stackrel{~}{h}]\underset{i=1}{\overset{q}{}}g^{(i)}.$$ (3.25) This condition was already noticed in . There are two general features of the conformal dimensions given by (3.12) and (3.23) that is important to emphasize. The first one is that the conformal dimension of the perturbation is independent of the value of $`p`$ in eqs. (3.8) and (3.21), which is a consequence of the fact that the potential in (2.3) and, hence, $`\mathrm{\Phi }`$ are invariant with respect to the gauge transformations (2.4). The second is that $`\mathrm{\Delta }_\mathrm{\Phi }`$ decreases with $`k`$, which means that the perturbation is always relevant for $`k`$ above some minimal value characteristic of each SSSG theory. Actually, $`\mathrm{\Delta }_\mathrm{\Phi }`$ vanishes when $`k\mathrm{}`$, which shows that the theory consists of $`\mathrm{dim}(g_0)p`$ bosonic massive particles in the the semi–classical and/or weak coupling limit. ### 3.2 Explicit calculation of $`\mathrm{\Delta }_\mathrm{\Phi }`$. In the following we will illustrate the general procedure to calculate the conformal dimension of the perturbation by considering three particular cases where $`g_0`$ is either simple ($`SU(2n)/SO(2n)`$), semisimple ($`G_2/SU(2)\times SU(2)`$), or the direct sum of a simple ideal and a one-dimensional centre ($`Sp(n)/U(n)`$). In all these examples, $`\mathrm{rank}(G/G_0)=\mathrm{rank}(G_0)`$, which, according to (2.9), means that $`\mathrm{\Lambda }_\pm `$ can be chosen such that $`p=0`$ in (3.8) and (3.21). The results for all the type I symmetric spaces are presented in table 2 and 3. Other useful features of the symmetric spaces of type I have been collected in tables 46. In these tables, we have already taken into account the following isomorphisms of Lie algebras: $$su(2)so(3)sp(1),so(5)sp(2),so(4)su(2)su(2),su(4)so(6).$$ In particular, this shows that $`SO(4)`$ is not simple and, therefore, no symmetric space with $`G=SO(4)`$ appears in the tables because it would not be of type I. Moreover, the symmetric space $`SU(2)/SO(2)`$ corresponds to the well known sine-Gordon theory where the field takes values in the abelian group $`SO(2)U(1)`$, and it has not been included in the tables. When $`G_0`$ is simple, it is worthwhile noticing that $`\mathrm{\Delta }_\mathrm{\Phi }`$ admits the general expression $$\mathrm{\Delta }_\mathrm{\Phi }=\frac{\stackrel{}{\mathrm{\Psi }}_g^2}{\stackrel{}{\mathrm{\Psi }}_{g_0}^2}\frac{h_g^{}}{2(k+h_{g_0}^{})}.$$ (3.26) #### 3.2.1 Example I: $`G/G_0=SU(2n)/SO(2n)`$, $`n>2`$. In this case, $`\overline{g}=A_{2n1}`$, $`r=2`$, and $`\stackrel{}{s}=(0,\mathrm{},0,1)`$, which follows from Theorem 3, part a), and the observation that the Dynkin diagram of $`\overline{g}_0=D_n`$ is a subdiagram of the Dynkin diagram of $`\overline{g}^{(r)}=A_{2n1}^{(2)}`$, as can be seen in fig. 1. Therefore, and taking into account (3.6), this symmetric space is of type A2. Fig. 1 also shows that the roots of $`\overline{g}_0=D_n`$ are short roots in the Dynkin diagram of $`\overline{g}^{(r)}=A_{2n1}^{(2)}`$, which means that $`\stackrel{}{\mathrm{\Psi }}_{D_n}^2=\stackrel{}{\mathrm{\Psi }}_{A_{2n1}^{(2)}}^2/2`$. Theorem 3, part b), also implies that $`[\overline{g}_0,\overline{g}_1]\overline{g}_1`$ is an irreducible representation of $`\overline{g}_0`$ with highest weight $`\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{\alpha }_n`$. Then, we can use the identity (3.15) in order to write $`\stackrel{}{\mathrm{\Lambda }}`$ as a linear combination of the roots of $`\overline{g}_0=D_n`$, namely $$\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{\alpha }_n=\stackrel{}{\alpha }_0+\stackrel{}{\alpha }_1+\mathrm{\hspace{0.17em}2}\underset{i=2}{\overset{n1}{}}\stackrel{}{\alpha }_i=\stackrel{}{\beta }_n+\stackrel{}{\beta }_{n1}+\mathrm{\hspace{0.17em}2}\underset{i=1}{\overset{n2}{}}\stackrel{}{\beta }_i.$$ (3.27) Moreover, taking into account that the highest root of $`D_n`$ is $$\stackrel{}{\mathrm{\Psi }}_{D_n}=\stackrel{}{\beta }_1+\stackrel{}{\beta }_n+\stackrel{}{\beta }_{n1}+\mathrm{\hspace{0.17em}2}\underset{i=2}{\overset{n2}{}}\stackrel{}{\beta }_i,$$ (3.28) eq. (3.27) simplifies to $`\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{\mathrm{\Psi }}_{D_n}+\stackrel{}{\beta }_1`$. All this allows one to easily calculate the quadratic Casimir of this highest weight representation: $$C_2(D_n)=\stackrel{}{\mathrm{\Lambda }},\stackrel{}{\mathrm{\Lambda }}+\mathrm{\hspace{0.17em}2}\stackrel{}{\delta }_{D_n}=\mathrm{\hspace{0.17em}2}n\stackrel{}{\mathrm{\Psi }}_{D_n}^2,$$ (3.29) where we have used the standard realization of the root system of $`D_n`$ as a sublattice of the real euclidean space $`\text{}^n`$: $$\mathrm{\Pi }_{D_n}=\left\{\stackrel{}{\beta }_1=\stackrel{}{v}_1\stackrel{}{v}_2,\mathrm{},\stackrel{}{\beta }_{n1}=\stackrel{}{v}_{n1}\stackrel{}{v}_n,\stackrel{}{\beta }_n=\stackrel{}{v}_{n1}+\stackrel{}{v}_n\right\},$$ (3.30) where $$\stackrel{}{v}_i\stackrel{}{v}_j=\frac{\stackrel{}{\mathrm{\Psi }}_{D_n}^2}{2}\delta _{ij},$$ (3.31) together with $$\stackrel{}{\mathrm{\Psi }}_{D_n}=\stackrel{}{v}_1+\stackrel{}{v}_2,\mathrm{and}2\stackrel{}{\delta }_{D_n}=\mathrm{\hspace{0.25em}2}\underset{i=1}{\overset{n}{}}(ni)\stackrel{}{v}_i.$$ (3.32) Therefore, taking into account (3.11) and (3.12), we conclude that the SSSG’s associated with the symmetric space $`SU(2n)/SO(2n)`$ are integrable perturbations of either the WZNW model corresponding to $`SO(2n)`$ at level $`k`$ ($`p=0`$) or a coset CFT of the form $`SO(2n)_k/U(1)^p`$ whose central charge is $$c_{\mathrm{CFT}}=\frac{kn(2n1)}{k+2(n1)}p,$$ (3.33) where $`0pn`$, and the perturbation has conformal dimension $$\mathrm{\Delta }_\mathrm{\Phi }=\frac{2n}{k+\mathrm{\hspace{0.17em}2}(n\mathrm{\hspace{0.17em}1})}.$$ (3.34) Notice that the perturbation is relevant for $`k>2`$. #### 3.2.2 Example II: $`G/G_0=G_2/SU(2)\times SU(2)`$. In this case, $`\overline{g}=G_2`$, $`r=1`$, and $`\stackrel{}{s}=(0,1,0)`$, as can be seen in fig. 2. Therefore, this symmetric space is of type A1. $`\overline{g}_0`$ is of the form $`\overline{g}_0=\overline{g}^{(1)}\overline{g}^{(2)}`$ with $`\overline{g}^{(1)}=\overline{g}^{(2)}=A_1`$. $`\overline{g}^{(1)}`$ and $`\overline{g}^{(2)}`$ are associated with the roots $`\stackrel{}{\alpha }_2`$ and $`\stackrel{}{\alpha }_0`$ of the Dynkin diagram of the affine algebra $`G_2^{(1)}`$, respectively, whose length is different and, therefore, $`\stackrel{}{\mathrm{\Psi }}_{g^{(2)}}^2=3\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}^2=\stackrel{}{\mathrm{\Psi }}_{G_2^{(1)}}^2`$. Theorem 3 implies that $`[\overline{g}_0,\overline{g}_1]\overline{g}_1`$ gives an irreducible representation of $`\overline{g}_0`$ with highest weight $`\stackrel{}{\mathrm{\Lambda }}=\stackrel{}{\alpha }_1`$. Then, since $`\stackrel{}{\alpha }_0+2\stackrel{}{\alpha }_1+3\stackrel{}{\alpha }_2=0`$, one can write $$\stackrel{}{\mathrm{\Lambda }}=\frac{3}{2}\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}+\frac{1}{2}\stackrel{}{\mathrm{\Psi }}_{g^{(2)}},$$ (3.35) which leads to $`C_2(g^{(1)})=\stackrel{}{\mathrm{\Lambda }},\stackrel{}{\mathrm{\Lambda }}+\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}={\displaystyle \frac{15}{4}}\stackrel{}{\mathrm{\Psi }}_{g^{(1)}}^2,`$ $`C_2(g^{(2)})=\stackrel{}{\mathrm{\Lambda }},\stackrel{}{\mathrm{\Lambda }}+\stackrel{}{\mathrm{\Psi }}_{g^{(2)}}={\displaystyle \frac{3}{4}}\stackrel{}{\mathrm{\Psi }}_{g^{(2)}}^2.`$ (3.36) Therefore, we conclude that the SSSG’s associated with this symmetric space are integrable perturbations of either the WZNW model corresponding to $`SU(2)_k\times SU(2)_{3k}`$ ($`p=0`$) or a coset CFT of the form $`SU(2)_k\times SU(2)_{3k}/U(1)^p`$ whose central charge is $$c_{\mathrm{CFT}}=\frac{3k}{k+2}+\frac{9k}{3k+2}p.$$ (3.37) Notice the relationship between the levels of the two $`SU(2)`$ factors, which is a consequence of (3.10). The perturbation has conformal dimension $$\mathrm{\Delta }_\mathrm{\Phi }=\frac{3}{4(k+2)}+\frac{15}{4(3k+2)}.$$ (3.38) Moreover, since in this case $`\mathrm{rank}(G_0)=\mathrm{rank}(G)=2`$, according to (2.9) there is a different SSSG theory for each integer $`p`$ in the range $`0p2\nu `$. The perturbation is relevant for $`k>2`$. #### 3.2.3 Example III: $`G/G_0=Sp(n)/U(n)=Sp(n)/U(1)\times SU(n)`$, $`n>1`$. In this case, $`\overline{g}=C_n`$, $`r=1`$, and $`\stackrel{}{s}=(1,0\mathrm{},0,1)`$, as can be seen in fig. 3, and the symmetric space is of type B. Then, $`\overline{g}_0=\overline{g}^{(1)}u(1)`$ with $`\overline{g}^{(1)}=A_{n1}`$, and $`\stackrel{}{\mathrm{\Psi }}_{A_{n1}}^2=\stackrel{}{\mathrm{\Psi }}_{C_n^{(1)}}^2/2`$. Recall the standard realization of the root systems of $`C_n`$ and $`A_{n1}`$ as sublattices of the real euclidean space $`\text{}^n`$: $`\mathrm{\Pi }_{C_n}=\left\{\stackrel{}{\alpha }_1=\stackrel{}{v}_1\stackrel{}{v}_2,\mathrm{},\stackrel{}{\alpha }_{n1}=\stackrel{}{v}_{n1}\stackrel{}{v}_n,\stackrel{}{\alpha }_n=\mathrm{\hspace{0.17em}2}\stackrel{}{v}_n\right\},\stackrel{}{\mathrm{\Psi }}_{C_n}=\stackrel{}{\alpha }_0=\mathrm{\hspace{0.17em}2}\stackrel{}{v}_1,`$ $`\mathrm{\Pi }_{A_{n1}}=\left\{\stackrel{}{\alpha }_1=\stackrel{}{v}_1\stackrel{}{v}_2,\mathrm{},\stackrel{}{\alpha }_{n1}=\stackrel{}{v}_{n1}\stackrel{}{v}_n\right\}\mathrm{\Pi }_{C_n},\stackrel{}{\mathrm{\Psi }}_{A_{n1}}=\stackrel{}{v}_1\stackrel{}{v}_n,`$ (3.39) where $$\stackrel{}{v}_i\stackrel{}{v}_j=\frac{\stackrel{}{\mathrm{\Psi }}_{C_n}^2}{4}\delta _{ij}.$$ (3.40) Using (3.39), $`\stackrel{}{\mathrm{\Lambda }}_1=2\stackrel{}{v}_1`$, which has to be split in two components. One, in the weight lattice of $`A_{n1}`$, and another, corresponding to the $`u(1)`$ subalgebra of $`\overline{g}_0`$, orthogonal to it. The decomposition is as follows $$\stackrel{}{\mathrm{\Lambda }}_1=\stackrel{}{\lambda }_{A_{n1}}+\stackrel{}{\lambda }_{u(1)}$$ (3.41) with $$\stackrel{}{\lambda }_{A_{n1}}=\underset{i=1}{\overset{n1}{}}\frac{2(ni)}{n}(\stackrel{}{v}_i\stackrel{}{v}_{i+1}),\stackrel{}{\lambda }_{u(1)}=\frac{2}{n}\underset{i=1}{\overset{n}{}}\stackrel{}{v}_i=\frac{1}{\sqrt{\pi n}}\stackrel{}{u}.$$ (3.42) This allows one to calculate the required quadratic Casimir $$C_2(A_{n1})=\stackrel{}{\lambda }_{A_{n1}},\stackrel{}{\lambda }_{A_{n1}}+\mathrm{\hspace{0.17em}2}\stackrel{}{\delta }_{A_{n1}}=\frac{(n1)(n+2)}{n}\stackrel{}{\mathrm{\Psi }}_{A_{n1}}^2,$$ (3.43) where we have used $$2\stackrel{}{\delta }_{A_{n1}}=\underset{i=1}{\overset{[n/2]}{}}(\stackrel{}{v}_i\stackrel{}{v}_{n+1i})(n+12i),$$ (3.44) with $`[n/2]`$ the integer part of $`n/2`$. The same result can be obtained by considering the other representation with highest weight $`\stackrel{}{\mathrm{\Lambda }}_2`$. Therefore, the SSSG theories associated with the symmetric space $`Sp(n)/U(n)`$ are integral perturbations of either the WZNW model corresponding to $`SU(n)`$ at level $`k`$ plus a massless boson ($`p=0`$) or a coset CFT of the form $`[SU(n)_k/U(1)^p]\times U(1)`$, whose central charge is given by (3.22): $$c_{\mathrm{CFT}}=\frac{k(n^21)}{k+n}+\mathrm{\hspace{0.25em}1}p.$$ (3.45) The conformal dimension of the perturbation is given by (3.23): $$\mathrm{\Delta }_\mathrm{\Phi }=\frac{(n1)(n+2)}{n(k+n)}+\frac{2}{kn},$$ (3.46) and there is a different SSSG theory for $`0pn1`$. The perturbation is relevant for $`k>2`$. ## 4 Integrability of the SSSG theories. The classical integrability of the Homogeneous (HSG) and Symmetric Space (SSSG) sine-Gordon theories is a consequence of the relationship between their equations of motion and the non-abelian affine Toda equations, which admit a zero-curvature representation. This implies the existence of an infinite number of conserved quantities, whose construction by means of the Drinfel’d-Sokolov procedure will be summarized in the first part of this section. Concerning the HSG theories, their quantum integrability was established in by explicitly checking that the conserved quantities of scale dimension $`\pm 2`$ remain conserved in the quantum theory after an appropriate renormalization, and invoking a well known argument due to Parke . For the SSSG theories, we expect that something similar happens and, hence, that they are also quantum integrable. However, the variety of different types of SSSG theories makes difficult to thoroughly check this conjecture. In general, and similarly to the HSG models, most of the SSSG theories exhibit classically conserved quantities of scale dimension $`\pm 1,\pm 2,\mathrm{}`$, and it should be possible to generalize the proof in to show that the conserved quantities of scale dimension $`\pm 2`$ also give rise to quantum conserved quantities. Nevertheless, there is an exception to this pattern: the SSSG theories associated with symmetric spaces of maximal rank, i.e., with symmetric spaces $`G/G_0`$ whose rank is $`\mathrm{rank}(G/G_0)=\mathrm{rank}(G)`$, which only have conserved quantities of odd scale dimension $`\pm 1,\pm 3,\mathrm{}`$. According to (2.9), they are relevant perturbations of either the WZNW model corresponding to $`G_0`$ ($`p=0`$) or a coset CFT of the form $`G_0/U(1)^p`$ with $`0p\mathrm{rank}(G_0)`$. The type I symmetric spaces of maximal rank are the following : $`SO(2n)/SO(n)\times SO(n),SO(2n+1)/SO(n)\times SO(n+1),`$ $`SU(n)/SO(n),Sp(n)/U(n),E_6/Sp(4),E_7/SU(8),`$ $`E_8/SO(16),F_4/Sp(3)\times SU(2),G_2/SU(2)\times SU(2).`$ (4.1) Notice that there is one for each simple compact Lie group $`G`$, which is related to its (unique) maximally non compact real form, also known as ‘split form’. Besides $`SU(2)/SO(2)`$, which corresponds to the sine-Gordon theory, the simplest symmetric space of maximal rank is $`SU(3)/SO(3)`$. Then, eq. (2.9) reads $`0p1`$, which means that it gives rise to two different SSSG theories, depending on the choice of $`\mathrm{\Lambda }_\pm `$. They are just the integrable perturbations of the $`SU(2)_k`$ WZNW model $`(p=0)`$ and its $`\stackrel{~}{so(2)}`$ reduction $`(p=1)`$ constructed by Brazhnikov , which is identified with the perturbation of the usual $`\text{}_k`$-parafermions by the second thermal operator. In this section we explicitly prove the quantum integrability of the SSSG theories associated with a maximal rank symmetric space and $`p=0`$, which will be called ‘Split models’. Actually, they can be distinguished as the only ones that provide relevant perturbations of WZNW models. Namely, we will check that their simplest higher spin classically conserved densities, which have spin $`\pm 3`$ instead of $`\pm 2`$, remain conserved in the quantum theory after an appropriate renormalization. According to , this is enough to establish their quantum integrability, which, together with the results of for the HSG theories, supports the conjecture that all the SSSG theories are quantum integrable. ### 4.1 Classical integrability. The zero-curvature form of the equations of motion (2.6) is $$[_++mh\mathrm{\Lambda }_+h^{},_{}+m\mathrm{\Lambda }_{}_{}hh^{}]=0.$$ (4.2) while the constraints arise naturally in the group theoretical description of the non-abelian affine Toda equations . In order to use the generalized Drinfel’d-Sokolov construction of to get the infinite number of conserved densities of (4.2), we have to associate the zero-curvature equation with a loop algebra. In our case, the relevant loop algebra is $`(g,\sigma )=_j\text{}(g,\sigma )_j`$ with $$(g,\sigma )_{2j}=\lambda ^jg_0,\mathrm{and}(g,\sigma )_{2j+1}=\lambda ^jg_1,$$ (4.3) where $`\lambda `$ is a spectral parameter; if $`\sigma =\sigma _{(\stackrel{}{s};r)}`$, then $`(g,\sigma )`$ is related to the (twisted if $`r1`$) affine Kac-Moody algebra $`\overline{g}^{(r)}`$ without central extension. Moreover, since $$[(g,\sigma )_j,(g,\sigma )_k](g,\sigma )_{j+k},$$ (4.4) the subspaces (4.3) define an integer gradation of $`(g,\sigma )`$. The equations of motion (4.2) remain unchanged under the transformation $`\mathrm{\Lambda }_\pm \lambda ^1\mathrm{\Lambda }_\pm `$ and, hence, the zero-curvature equation can actually be associated with $`(g,\sigma )`$ and the Lax operator $`L=_{}+\mathrm{\Lambda }+q`$, where $$\mathrm{\Lambda }=m\lambda \mathrm{\Lambda }_{}(g,\sigma )_1\mathrm{and}q=_{}hh^{}(g,\sigma )_0.$$ (4.5) The generalized Drinfel’d-Sokolov construction goes as follows . First, there is some local function $`y`$ of the ‘potential’ $`q`$ of the form <sup>1</sup><sup>1</sup>1By a local function of $`q`$ we mean a $`_{}`$-differential polynomial in the components of $`q`$, i.e., a polynomial in the components of $`q`$, $`_{}q`$, $`_{}^2q`$, $`\mathrm{}`$ $$y=\underset{n>0}{}y^a(n)\lambda ^kt^aIm(ad_\mathrm{\Lambda })_{<0},$$ (4.6) that ‘abelianizes’ the Lax operator in the following sense: $$e^yLe^y=e^y(_{}+\mathrm{\Lambda }+q)e^y=_{}+\mathrm{\Lambda }+H,$$ (4.7) where $`H\mathrm{Ker}(\mathrm{ad}_\mathrm{\Lambda })_0`$ is another local function of $`q`$, and $`\{t^a\}`$ is a basis for $`g`$ whose standard realization is presented in the Appendix. Then, eqs. (4.2) and (4.7) imply $$e^y\left(_++mh(\lambda ^1\mathrm{\Lambda }_+)h^{}\right)e^y=_++\overline{H}.$$ (4.8) where $`\overline{H}`$ also takes values in $`\mathrm{Ker}(\mathrm{ad}_\mathrm{\Lambda })_0`$ and, therefore, the zero-curvature equation becomes $$_{}\overline{H}_+H=[\overline{H},H].$$ (4.9) The components of this equation along the centre of $`\mathrm{Ker}(\mathrm{ad}_\mathrm{\Lambda })`$ provide an infinite number of local conservation laws. To be precise, let $`b\mathrm{Cent}\left(\mathrm{Ker}(\mathrm{ad}_\mathrm{\Lambda })\right)_j`$ with $`j>0`$, and define $$_j^{(0)}[b]=b,H,\overline{}_j^{(0)}[b]=b,\overline{H}.$$ (4.10) The corresponding local conservation law is $$_+_j^{(0)}[b]=_{}\overline{}_j^{(0)}[b].$$ (4.11) Moreover, one can check that the conserved densities given by (4.10) have scale dimension $`j+1>1`$ with respect to the scale transformations $`x_\pm x_\pm /\rho `$, which means that they give rise to conserved quantities of scale dimension $`j>0`$. The same procedure can be repeated by changing the Lax operator $`L`$ by $$\overline{L}=_++m\lambda ^1\mathrm{\Lambda }_++h^{}_+h$$ (4.12) in eq. (4.7), and the result is the construction of another infinite number of local conserved quantities with scale dimension $`j<0`$, i.e., negative scale dimension. Notice that $`L`$ and $`\overline{L}`$ are conjugate by the transformation (2.7) and, moreover, that the dimensions of $`(g,\sigma )_j`$ and $`(g,\sigma )_j`$ are equal. Therefore, for each conserved quantity of positive scale dimension $`j`$ there will be another conserved quantity of scale dimension $`j`$, and both are conjugate by (2.7). In particular, if the theory is parity invariant both conserved quantities are parity conjugate. This allows one to restrict the analysis to the conserved quantities with positive scale dimension. Taking into account all this, the resulting number of classically conserved quantities with scale dimension $`\pm j`$ is given by the dimension of $`\mathrm{Cent}\left(\mathrm{Ker}(\mathrm{ad}_\mathrm{\Lambda })\right)_{\pm j}`$. These dimensions can be easily calculated when $`\mathrm{\Lambda }_\pm `$ are chosen to be regular, which means that $`\mathrm{Ker}(\mathrm{ad}_{\mathrm{\Lambda }_\pm })`$ is a Cartan subalgebra of $`g`$ and $`p=\mathrm{rank}(G)\mathrm{rank}(G/G_0)`$ in (2.9). Then, the Drinfel’d-Sokolov construction produces exactly $`\mathrm{rank}(G/G_0)`$ local conserved quantities for each odd scale dimension $`\pm 1,\pm 3,\mathrm{}`$, and $`\mathrm{rank}(G)\mathrm{rank}(G/G_0)`$ for each even scale dimension $`\pm 2,\pm 4,\mathrm{}`$. Consider now a SSSG theory related to a symmetric space of maximal rank with $`p`$ in the range (2.9) ($`p=0`$ corresponds to the case when $`\mathrm{\Lambda }_\pm `$ are regular). Then, there are $`\mathrm{rank}(G/G_0)p`$ local conserved quantities for each odd scale dimension, $`\pm 1,\pm 3,\mathrm{}`$, and no conserved quantities with even scale dimension. This can be easily proved by choosing the Cartan subalgebra such that it contains the abelian subalgebra $`g_0^0=u(1)^p`$ given by (2.2). These results agree with what has been anticipated at the beginning of this section: in general, the simplest higher spin conserved quantities of a SSSG theory have spin $`\pm 2`$. However, if $`\mathrm{rank}(G/G_0)=\mathrm{rank}(G)`$, the simplest conserved quantities will have scale dimension $`\pm 3`$. ### 4.2 Quantum integrability of the Split models. In the following, we will check that, after a suitable renormalization, the conserved densities of spin $`\pm 1`$ and $`\pm 3`$ remain conserved in the quantum version of the Split models. The explicit expressions for the relevant classical conserved densities are given in the Appendix. The proof will be similar to the one presented in for the HSG theories, which uses conformal perturbation theory. However, since the simplest higher spin conserved quantity is of spin $`3`$ instead of $`2`$, this case will be even more involved, and in order to avoid unnecessary complications we will restrict ourselves to the Split models with $`G_0`$ simple. We will also fix the normalization of the invariant bilinear form of $`g`$, $`,`$, such that $`\stackrel{}{\mathrm{\Psi }}_{g_0}^2=2`$. Since the quantum SSSG theories can be described as perturbed conformal field theories, the existence of quantum conserved quantities can be investigated by using the methods of . In the presence of the perturbation (2.8) any chiral field $`(z)`$, which in the unperturbed CFT satisfies $`\overline{}(z)=0`$, acquires a $`\overline{z}`$ dependence given by $$\overline{}(z,\overline{z})=km^2_z\frac{dw}{2\pi i}\mathrm{\Phi }(w,\overline{z})(z),$$ (4.13) where we have introduced the notation $`z=x_{}`$, $`\overline{z}=x_+`$, $`=_{}`$, and $`\overline{}=_+`$ reminiscent of euclidean space. This contribution actually corresponds to the lowest order in perturbation theory; however, if the condition of super-renormalizability at first order $`2\mathrm{\Delta }_\mathrm{\Phi }1`$ is satisfied, no counterterms are needed to renormalize (2.8), and the previous equation is expected to be exact . Actually, for any SSSG theory, this condition is always fulfilled for $`k`$ above some minimal value characteristic of the theory (see the comments at the end of Section 3.1). Therefore, in the perturbed CFT the chiral field $`(z)`$ will become a conserved quantity if the right-hand-side of (4.13) is a total $``$ derivative, i.e. if (4.13) can be written as $`\overline{}=\overline{}`$ where $`\overline{}`$ is another field of the original CFT. For this condition to be satisfied, the residue of the simple pole in the OPE between $`(z)`$ and $`\mathrm{\Phi }(w,\overline{z})`$ has to be a total $``$ derivative; namely, $$\mathrm{\Phi }(w,\overline{z})(z)=\underset{n>1}{}\frac{\{\mathrm{\Phi }\}_n(z,\overline{z})}{(wz)^n}+\frac{\overline{}(z,\overline{z})}{(wz)}+\mathrm{}$$ (4.14) Notice that the residues of the simple poles in the OPE’s $`\mathrm{\Phi }(w,\overline{z})(z)`$ and $`(z)\mathrm{\Phi }(w,\overline{w})`$ differ only in a total $``$ derivative and, in practice, we will always consider the latter whose expression is usually simpler. For the Split models, the unperturbed CFT is just the WZNW model corresponding to $`G_0`$ at level $`k`$. Using the conventions of the Appendix, the subset of generators of $`g`$ given by $`t^\alpha `$ for each positive root $`\stackrel{}{\alpha }`$ of $`g`$ provides a suitable basis of generators for the Lie subalgebra $`g_0`$. This means that the operator algebra of the $`G_0`$–WZNW model can be realized as a subset of the operator algebra of the WZNW model associated to $`G`$ that, in particular, includes the chiral currents $`J^\alpha (z)`$ and $`\overline{J}^\alpha (\overline{z})`$ which satisfy the OPE $$J^\alpha (w)J^\beta (z)=\frac{\mathrm{}^2k\delta ^{\alpha \beta }}{(wz)^2}+\frac{\mathrm{}f^{\alpha \beta \gamma }J^\gamma (z)}{(wz)}+\mathrm{}$$ (4.15) Moreover, it will be useful to recall the following well known identities . The first one is $$J^\alpha (z)h(w,\overline{w})=\mathrm{}\frac{t^\alpha h(w,\overline{w})}{zw}+\mathrm{},$$ (4.16) which is satisfied in an arbitrary representation of $`G_0`$ and exhibits that the WZNW field $`h`$ is a primary field. The second is the relation between the WZNW field and the chiral currents $$\mathrm{}(k+h_{g_0}^{})h=\underset{\alpha }{}(J^\alpha t^\alpha h),$$ (4.17) where $`(AB)(z)`$ is the normal ordered product of two operators $`A(z)`$ and $`B(z)`$, which will be defined by adopting the conventions of . At this point, it is also convenient to recall that the classical expressions are recovered from the quantum ones by means of (see (4.5) and (4.17)) $$J^\alpha t^\alpha =(\mathrm{}k)q,(\mathrm{}k)\frac{1}{\beta ^2},\mathrm{and}k\mathrm{},$$ (4.18) which, according to (3.9), amounts to take the classical ($`\mathrm{}0`$) or weak coupling ($`\beta ^20`$) limits. Therefore, the differential polynomials in the components of $`q`$ that appear in the classical local conserved densities have to be changed into differential polynomials in the chiral currents $`J^\alpha `$. In the rest of this section we will show explicitly the Plank constant to make the distinction between classical and quantum contributions simpler. #### 4.2.1 Quantum spin-2 conserved densities (spin-1 conserved quantities). Consider a generic normal ordered local spin-2 operator $$_2(z)=D_{\alpha \beta }(J^\alpha J^\beta )(z),$$ (4.19) where the numerical coefficients satisfy $`D_{\alpha \beta }=D_{\beta \alpha }`$. Taking into account that the perturbing operator $`\mathrm{\Phi }=\mathrm{\Lambda }_{},P`$, with $`P=(h\mathrm{\Lambda }_+h^{})`$, is just the WZNW field taken in the representation of $`G_0`$ provided by $`[g_0,g_1]g_1`$, eqs. (4.16) and (4.17) become $`J^\alpha (z)v,P(w,\overline{w})=\mathrm{}{\displaystyle \frac{[v,t^\alpha ],P(w,\overline{w})}{(zw)}}+\mathrm{}`$ (4.20) $`\mathrm{}(k+h_{g_0}^{})v,P=\left(J^\alpha [v,t^\alpha ],P\right),`$ (4.21) where $`v`$ is an arbitrary element in $`g_1`$. Then, it is easy to show that the residue of the simple pole in the OPE $`_2(z)\mathrm{\Phi }(w,\overline{w})`$ is given by $$\mathrm{Res}\left(_2(z)\mathrm{\Phi }(w,\overline{w})\right)=2\mathrm{}\left(J^\beta D_{\alpha \beta }[\mathrm{\Lambda }_{},t^\alpha ],P(w,\overline{w})\right).$$ (4.22) Taking into account (4.21), this residue is a total $``$ derivative if the tensor $`D_{\alpha \beta }`$ satisfies <sup>2</sup><sup>2</sup>2Explicit factors of $`\mathrm{}k`$ are included to take account of the relation $`J^\alpha t^\alpha =(\mathrm{}k)q`$ in the $`k\mathrm{}`$ limit and recover the classical expressions given in the Appendix (see eq. (4.18)). $$(\mathrm{}k)^2D_{\alpha \beta }[\mathrm{\Lambda }_{},t^\beta ]=\frac{1}{2}[\stackrel{}{\mu }\stackrel{}{t},t^\alpha ],$$ (4.23) for any $`\mathrm{rank}(g)`$-component vector $`\stackrel{}{\mu }`$. This leads to the following solutions labelled by $`\stackrel{}{\mu }`$ ($`\stackrel{}{\lambda }`$ is defined in (A.4)) $$(\mathrm{}k)^2D_{\alpha \beta }(\stackrel{}{\mu })=\frac{1}{2}\frac{\stackrel{}{\mu }\stackrel{}{\alpha }}{\stackrel{}{\lambda }\stackrel{}{\alpha }}\delta _{\alpha \beta }=mD_{\alpha \beta }^{(0)}(\stackrel{}{\mu }),$$ (4.24) which satisfy $$\mathrm{Res}\left(_2(z)\mathrm{\Phi }(w,\overline{w})\right)=\frac{(k+h_{g_0}^{})}{k^2}\stackrel{}{\mu }\stackrel{}{t},P(w,\overline{w}).$$ (4.25) Therefore, the quantities $$(\mathrm{}k)^2_2(\stackrel{}{\mu })=\frac{1}{2}\underset{\alpha >0}{}\frac{\stackrel{}{\mu }\stackrel{}{\alpha }}{\stackrel{}{\lambda }\stackrel{}{\alpha }}(J^\alpha J^\alpha )$$ (4.26) provide $`\mathrm{rank}(g)`$ linearly independent conserved densities of spin-2. The field $`\overline{}_2(\stackrel{}{\mu })`$ that satisfies the conservation law $`\overline{}_2=\overline{}_2`$ is obtained through the explicit evaluation of the integral (4.13); the result is $$(\mathrm{}k)^2\overline{}_2(\stackrel{}{\mu })=m^2(\mathrm{}k)^2\left(1+\frac{4h_{g_0}^{}\mathrm{\Psi }_g^2h_g^{}}{4k}\right)\stackrel{}{\mu }\stackrel{}{t},P.$$ (4.27) It has already been pointed out that these quantum conserved quantities can be understood as the renormalization of the classical ones calculated in the Appendix. In this case, the relationship is particularly simple: the spin-2 quantum conserved densities equal the normal ordered classical conserved ones of scale dimension 2, up to the multiplicative renormalization in (4.27). In more general cases, one has to add quantum $`O(1/k)O(\mathrm{}\beta ^2)`$ corrections; however, the number of them must be finite as a reflection of the fact that the perturbation is relevant and, therefore, the quantum theory is super-renormalizable. As expected, the stress-energy tensor is a particular example of a spin-2 conserved quantity. It is recovered from (4.26) for the particular choice $`\stackrel{}{\mu }=\stackrel{}{\lambda }`$: $$(\mathrm{}k)^2_2(\stackrel{}{\lambda })=\frac{1}{2}(J^\alpha J^\alpha )=(k+h_{g_0}^{})\left\{\frac{1}{2(k+h_{g_0}^{})}(J^\alpha J^\alpha )\right\}=(k+h_{g_0}^{})T_{z,z},$$ (4.28) in agreement with the Sugawara construction. Consequently, $$(\mathrm{}k)^2\overline{}_2(\stackrel{}{\lambda })=m^2(\mathrm{}k)^2\left(1+\frac{4h_{g_0}^{}\mathrm{\Psi }_g^2h_g^{}}{4k}\right)\mathrm{\Lambda }_{},P=(k+h_{g_0}^{})T_{z,\overline{z}}.$$ (4.29) #### 4.2.2 Quantum spin-4 conserved densities (spin-3 conserved quantities). Up to total $``$-derivatives, the most general form of a spin-4 normal ordered operator constructed with the WZNW currents $`J^\alpha `$ is $$_4(z)=R_{\alpha \beta \gamma \rho }(J^\alpha (J^\beta (J^\gamma J^\rho )))(z)+P_{\alpha \beta \gamma }(J^\alpha (J^\beta J^\gamma ))(z)+Q_{\alpha \beta }(J^\alpha ^2J^\beta )(z),$$ (4.30) with the following constraints: $`R_{\alpha \beta \gamma \rho }=R_{(\alpha \beta \gamma \rho )},Q_{\alpha \beta }=Q_{\beta \alpha },`$ $`P_{\alpha \beta \gamma }=P_{\beta \alpha \gamma },P_{\alpha \beta \gamma }+P_{\gamma \beta \alpha }+P_{\alpha \gamma \beta }=\mathrm{\hspace{0.17em}0},`$ (4.31) where the first one indicates that $`R_{\alpha \beta \gamma \rho }`$ is a totally symmetric tensor. Then, the residue of the simple pole in the OPE $`_4(z)\mathrm{\Phi }(w,\overline{w})`$ is $`\mathrm{Res}\left(_4(z)\mathrm{\Phi }(w,\overline{w})\right)`$ $`=`$ $`\mathrm{}\{(J^\alpha (J^\beta (J^\gamma \mathrm{\Omega }_{\alpha \beta \gamma },P)))+(J^\alpha (J^\beta M_{\alpha \beta },P))+`$ (4.32) $`+(^2J^\alpha T_\alpha ,P)\}(w,\overline{w}),`$ where $`\mathrm{\Omega }_{\alpha \beta \gamma }=\mathrm{\hspace{0.17em}4}R_{\alpha \beta \gamma \rho }[\mathrm{\Lambda }_{},t^\rho ],`$ (4.33) $`M_{\alpha \beta }=12\mathrm{}R_{\alpha \beta \gamma \rho }[[\mathrm{\Lambda }_{},t^\gamma ],t^\rho ]2(P_{\alpha \beta \gamma }P_{\beta \gamma \alpha })[\mathrm{\Lambda }_{},t^\gamma ],`$ (4.34) $`T_\alpha =\mathrm{\hspace{0.17em}2}\mathrm{}^2R_{\alpha \beta \gamma \rho }[[[\mathrm{\Lambda }_{},t^\beta ]t^\gamma ],t^\rho ]+\mathrm{}(P_{\alpha \beta \gamma }P_{\gamma \beta \alpha })[[\mathrm{\Lambda }_{},t^\beta ],t^\gamma ]+`$ $`+(2Q_{\alpha \rho }+\mathrm{}f^{\beta \gamma \alpha }P_{\beta \rho \gamma })[\mathrm{\Lambda }_{},t^\rho ].`$ (4.35) The condition that $`_4(z)`$ is a conserved quantity is equivalent to the existence of two tensors $`F_{\alpha \beta }`$ and $`R_\alpha `$ taking values in $`g_1`$ such that $$Res\left(I_4(z)\mathrm{\Phi }(w,\overline{w})\right)=(J^\alpha (J^\beta F_{\alpha \beta },P))+(J^\alpha R_\alpha ,P).$$ (4.36) We will assume that $`F_{\alpha \beta }`$ is totally symmetric because, otherwise, its antisymmetric part $`F_{[\gamma \beta ]}`$ can be absorbed in $`R_\alpha `$ through the transformation $$R_\alpha R_\alpha +\frac{\mathrm{}}{2}f^{\alpha \beta \gamma }F_{[\gamma \beta ]}.$$ (4.37) If we define $$\widehat{F}_{\alpha \beta }=\frac{F_{\alpha \beta }}{\mathrm{}(k+h_{g_o}^{})}\mathrm{and}\widehat{R}_\alpha =\frac{R_\alpha }{\mathrm{}(k+h_{g_o}^{})},$$ (4.38) eqs. (4.33)–(4.38) can be combined to obtain the following tensor relations $`12R_{\alpha \beta \gamma \rho }[\mathrm{\Lambda }_{},t^\rho ]=[\widehat{F}_{\alpha \beta },t^\gamma ]+[\widehat{F}_{\alpha \gamma },t^\beta ]+[\widehat{F}_{\gamma \beta },t^\alpha ],`$ (4.39) $`2(P_{\alpha \beta \gamma }P_{\beta \gamma \alpha })[\mathrm{\Lambda }_{},t^\gamma ]=\mathrm{\hspace{0.17em}2}F_{\alpha \beta }+[\widehat{R}_\alpha ,t^\beta ]+\mathrm{}\{[\widehat{F}_{\beta \rho },[t^\alpha ,t^\rho ]]+`$ $`+[[\widehat{F}_{\alpha \beta },t^\rho ],t^\rho ]+[[\widehat{F}_{\alpha \rho },t^\beta ],t^\rho ]+[[\widehat{F}_{\beta \rho },t^\alpha ],t^\rho ]\},`$ (4.40) $`2Q_{\alpha \beta }[\mathrm{\Lambda }_{},t^\beta ]=R_\alpha +\mathrm{}\left\{(P_{\alpha \beta \gamma }P_{\gamma \beta \alpha })[[\mathrm{\Lambda }_{},t^\beta ],t^\gamma ]+f^{\beta \gamma \alpha }P_{\beta \rho \gamma }[\mathrm{\Lambda }_{},t^\rho ]\right\}+`$ $`+2\mathrm{}^2\left\{R_{\alpha \beta \gamma \rho }[[[\mathrm{\Lambda }_{},t^\beta ],t^\gamma ],t^\rho ]+{\displaystyle \frac{1}{3}}f^{\xi \gamma \rho }f^{\beta \rho \alpha }[\widehat{F}_{\beta \xi },t^\gamma ]\right\},`$ (4.41) where (4.39) does not contain any explicit quantum correction at all, which means that the same relation will hold at the classical level. In fact, taking into account the explicit expressions for the classical densities given in the Appendix and (4.18), for each $`\mathrm{rank}(g)`$-component vector $`\stackrel{}{\mu }`$ there is a classical solution of (4.39)–(4.41) given by $$R_{\alpha \beta \gamma \rho }=\frac{(\mathrm{}k)^4}{m^3}R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\mu }),P_{\alpha \beta \gamma }=\frac{(\mathrm{}k)^3}{m^3}P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\mu }),Q_{\alpha \beta }=\frac{(\mathrm{}k)^2}{m^3}Q_{\alpha \beta }^{(0)}(\stackrel{}{\mu }),$$ (4.42) which satisfy the classical limit of those equations: $`12m^3R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\mu })[\mathrm{\Lambda }_{},t^\rho ]=(\mathrm{}k)^4\left\{[\widehat{F}_{\alpha \beta }^{(0)}(\stackrel{}{\mu }),t^\gamma ]+[\widehat{F}_{\alpha \gamma }^{(0)}(\stackrel{}{\mu }),t^\beta ]+[\widehat{F}_{\beta \gamma }^{(0)}(\stackrel{}{\mu }),t^\alpha ]\right\},`$ (4.43) $`2m^3\left(P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\mu })P_{\beta \gamma \alpha }^{(0)}(\stackrel{}{\mu })\right)[\mathrm{\Lambda }_{},t^\gamma ]=(\mathrm{}k)^3\left\{2(\mathrm{}k)\widehat{F}_{\alpha \beta }^{(0)}(\stackrel{}{\mu })+[\widehat{R}_\alpha ^{(0)}(\stackrel{}{\mu }),t^\beta ]\right\},`$ (4.44) $`2m^3Q_{\alpha \beta }^{(0)}(\stackrel{}{\mu })[\mathrm{\Lambda }_{},t^\beta ]=(\mathrm{}k)^3\widehat{R}_\alpha ^{(0)}(\stackrel{}{\mu }).`$ (4.45) Using eqs. (A.8)–(A.10), one can obtain the following explicit expressions for the classical limit of the corresponding tensors $`\widehat{F}_{\alpha \beta }(\stackrel{}{\mu })`$ and $`\widehat{R}_\alpha (\stackrel{}{\mu })`$ $`(\mathrm{}k)^3R_\alpha ^{(0)}(\stackrel{}{\mu })={\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\mu }}{(\stackrel{}{\lambda }\stackrel{}{\alpha })^2}}t^{\overline{\alpha }},`$ (4.46) $`(\mathrm{}k)^4F_{\alpha \beta }^{(0)}(\stackrel{}{\mu })={\displaystyle \frac{t^{\overline{\gamma }}}{4}}\{{\displaystyle \frac{1}{4}}({\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\mu }}{(\stackrel{}{\lambda }\stackrel{}{\alpha })^2}}f^{\overline{\alpha }\beta \overline{\gamma }}+{\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\mu }}{(\stackrel{}{\lambda }\stackrel{}{\beta })^2}}f^{\overline{\beta }\alpha \overline{\gamma }})+`$ $`+\left({\displaystyle \frac{f^{\gamma \alpha \beta }}{\stackrel{}{\lambda }\stackrel{}{\gamma }}}+{\displaystyle \frac{3}{4}}{\displaystyle \frac{\stackrel{}{\lambda }\stackrel{}{\gamma }f^{\overline{\alpha }\gamma \overline{\beta }}}{(\stackrel{}{\lambda }\stackrel{}{\alpha })(\stackrel{}{\lambda }\stackrel{}{\beta })}}\right)\left({\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}{\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\alpha }}}\right)`$ $`{\displaystyle \frac{\stackrel{}{\gamma }\stackrel{}{\mu }}{4\stackrel{}{\lambda }\stackrel{}{\gamma }}}({\displaystyle \frac{f^{\overline{\beta }\alpha \overline{\gamma }}}{\stackrel{}{\lambda }\stackrel{}{\beta }}}+{\displaystyle \frac{f^{\overline{\alpha }\beta \overline{\gamma }}}{\stackrel{}{\lambda }\stackrel{}{\alpha }}})\}{\displaystyle \frac{(\stackrel{}{\alpha }\stackrel{}{\mu })\delta _{\alpha \beta }}{2(\stackrel{}{\lambda }\stackrel{}{\alpha })^2}}\stackrel{}{\alpha }\stackrel{}{t},`$ (4.47) where $`\stackrel{}{\alpha }\stackrel{}{t}=\alpha ^At^A`$ is in the Cartan subalgebra of $`g`$. In particular, when $`\stackrel{}{\mu }=\stackrel{}{\lambda }`$ the previous equations reduce to $`(\mathrm{}k)^3\widehat{R}_\alpha ^{(0)}(\stackrel{}{\lambda })={\displaystyle \frac{1}{\stackrel{}{\lambda }\stackrel{}{\alpha }}}t^{\overline{\alpha }},`$ (4.48) $`(\mathrm{}k)^4\widehat{F}_{\alpha \beta }^{(0)}(\stackrel{}{\lambda })={\displaystyle \frac{\delta _{\alpha \beta }}{2\stackrel{}{\lambda }\stackrel{}{\alpha }}}\stackrel{}{\alpha }\stackrel{}{t}.`$ (4.49) Once we have solved the eqs. (4.39)–(4.41) in the classical limit, we will try to find the solutions for the full equations. Since they become quite complicated, we will only obtain a single one, which is enough to establish the quantum integrability of these models. Projecting eqs. (4.39)–(4.41) on the CSA we get $`\widehat{F}_{\alpha \beta }^{}{}_{}{}^{\overline{\gamma }}\gamma ^A+\widehat{F}_{\beta \gamma }^{}{}_{}{}^{\overline{\alpha }}\alpha ^A+\widehat{F}_{\alpha \gamma }^{}{}_{}{}^{\overline{\beta }}\beta ^A=0,`$ (4.50) $`2\mathrm{}(k+h_{g_o}^{}{\displaystyle \frac{\mathrm{\Psi }_g^2h_g^{}}{4}})\widehat{F}_{\alpha \beta }^{}{}_{}{}^{A}\mathrm{\hspace{0.17em}2}\mathrm{}(\widehat{F}_{\alpha \beta }^{}{}_{}{}^{B}\alpha ^B)\alpha ^A=\widehat{R}_{\alpha }^{}{}_{}{}^{\overline{\beta }}\beta ^A+`$ $`+\mathrm{}\{f^{\alpha \rho \delta }\widehat{F}_{\beta \rho }^{\overline{\delta }}+f^{\overline{\gamma }\alpha \overline{\delta }}\widehat{F}_{\beta \delta }^{}{}_{}{}^{\overline{\gamma }}+f^{\overline{\gamma }\beta \overline{\delta }}\widehat{F}_{\alpha \delta }^{}{}_{}{}^{\overline{\gamma }}\}\delta ^A,`$ (4.51) $`\mathrm{}(k+h_{g_0}^{})\widehat{R}_\alpha ^A={\displaystyle \frac{\mathrm{}}{2}}\widehat{R}_\beta ^{\overline{\xi }}f^{\overline{\xi }\alpha \overline{\beta }}\beta ^A+{\displaystyle \frac{\mathrm{}^2}{2}}\{f^{\alpha \rho \delta }\{\widehat{F}_{\beta \rho }^{\overline{\xi }}f^{\overline{\xi }\delta \overline{\beta }}+(\widehat{F}_{\beta \rho }^B\delta ^B)\delta _{\overline{\delta }\overline{\beta }}\}+`$ $`+f^{\beta \rho \delta }f^{\overline{\xi }\delta \overline{\beta }}\widehat{F}_{\alpha \rho }^{\overline{\xi }}+2f^{\overline{\xi }\rho \overline{\delta }}f^{\overline{\delta }\rho \overline{\beta }}\widehat{F}_{\alpha \beta }^{\overline{\xi }}+\mathrm{\hspace{0.17em}2}f^{\overline{\xi }\beta \overline{\delta }}f^{\overline{\delta }\rho \overline{\beta }}\widehat{F}_{\alpha \rho }^{\overline{\xi }}+\mathrm{\hspace{0.17em}2}f^{\overline{\xi }\alpha \overline{\delta }}f^{\overline{\delta }\rho \overline{\beta }}\widehat{F}_{\beta \rho }^{\overline{\xi }}+`$ $`+2f^{\overline{\alpha }\rho \overline{\beta }}(\widehat{F}_{\beta \rho }^B\alpha ^B)\}\beta ^A+2\mathrm{}^2f^{\overline{\rho }\gamma \overline{\beta }}\{(\widehat{F}_{\gamma \rho }^B\alpha ^B)\delta _{\overline{\alpha }\overline{\beta }}+\widehat{F}_{\alpha \gamma }^{\overline{\xi }}f^{\overline{\xi }\rho \overline{\beta }}+`$ $`+\widehat{F}_{\rho \gamma }^{\overline{\xi }}f^{\overline{\xi }\alpha \overline{\beta }}+\widehat{F}_{\rho \alpha }^{\overline{\xi }}f^{\overline{\xi }\gamma \overline{\beta }}\}\rho ^A{\displaystyle \frac{2\mathrm{}^2}{3}}f^{\xi \gamma \rho }f^{\beta \rho \alpha }\widehat{F}_{\beta \xi }^{\overline{\gamma }}\gamma ^A.`$ (4.52) where we have used the following definitions: $$\widehat{F}_{\alpha \beta }^{\overline{\gamma }}=\widehat{F}_{\alpha \beta },t^{\overline{\gamma }},\widehat{F}_{\alpha \beta }^A=\widehat{F}_{\alpha \beta },t^A,\widehat{R}_\alpha ^{\overline{\beta }}=\widehat{R}_\alpha ,t^{\overline{\beta }},\widehat{R}_\alpha ^A=\widehat{R}_\alpha ,t^A.$$ (4.53) Taking into account (4.50)–(4.52) together with the constraints (4.31) one can check that there is a particular solution for eq. (4.36) given by $`(\mathrm{}k)^4\widehat{F}_{\alpha \beta }=(\mathrm{}k)^4\widehat{F}_{\alpha \beta }^{(0)}(\stackrel{}{\lambda })={\displaystyle \frac{\delta _{\alpha \beta }}{\mathrm{\hspace{0.17em}2}\stackrel{}{\lambda }\stackrel{}{\alpha }}}\stackrel{}{\alpha }\stackrel{}{t}.`$ (4.54) $`(\mathrm{}k)^3\widehat{R}_\alpha =(\mathrm{}k)^3\left\{\widehat{R}_\alpha ^{(0)}(\stackrel{}{\lambda })+{\displaystyle \frac{1}{k}}\widehat{R}_\alpha ^{(1)}(\stackrel{}{\lambda })\right\},`$ (4.55) where $`\widehat{R}_\alpha ^{(0)}(\stackrel{}{\lambda })`$ is the classical solution (4.48) and $$\widehat{R}_\alpha ^{(1)}(\stackrel{}{\lambda })=\frac{\mathrm{\Psi }_g^2h_g^{}\mathrm{\hspace{0.17em}4}h_{g_0}^{}+\mathrm{\hspace{0.17em}4}\stackrel{}{\alpha }^2}{4\stackrel{}{\lambda }\stackrel{}{\alpha }}t^{\overline{\alpha }}.$$ (4.56) For this particular solution, eqs. (4.39)–(4.41) imply $$(\mathrm{}k)^4R_{\alpha \beta \gamma \rho }=m^3R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\lambda }),$$ (4.57) where $`R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\lambda })`$ is just the classical solution (A.11), which is a consequence of the absence of quantum corrections in eq. (4.39), $$(\mathrm{}k)^3P_{\alpha \beta \gamma }=m^3\left\{P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\lambda })+\frac{1}{k}P_{\alpha \beta \gamma }^{(1)}(\stackrel{}{\lambda })\right\},$$ (4.58) where $`P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\lambda })`$ is the classical solution given by (A.12) and the quantum correction $`P_{\alpha \beta \gamma }^{(1)}(\stackrel{}{\lambda })`$ is $`m^3P_{\alpha \beta \gamma }^{(1)}(\stackrel{}{\lambda })={\displaystyle \frac{1}{6\stackrel{}{\lambda }\stackrel{}{\gamma }}}\{{\displaystyle \frac{\mathrm{\Psi }_g^2h_g^{}4h_{g_0}^{}}{4}}\{{\displaystyle \frac{f^{\overline{\alpha }\beta \overline{\gamma }}}{\stackrel{}{\lambda }\stackrel{}{\alpha }}}+{\displaystyle \frac{f^{\overline{\beta }\alpha \overline{\gamma }}}{\stackrel{}{\lambda }\stackrel{}{\beta }}}\}+`$ $`+{\displaystyle \frac{3f^{\alpha \beta \gamma }}{2}}\{{\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\gamma }}{\stackrel{}{\lambda }\stackrel{}{\alpha }}}{\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\gamma }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}\}\},`$ (4.59) and finally $$(\mathrm{}k)^2Q_{\alpha \rho }=m^3\left\{Q_{\alpha \rho }^{(0)}(\stackrel{}{\lambda })+\frac{1}{k}Q_{\alpha \rho }^{(1)}(\stackrel{}{\lambda })+\frac{1}{k^2}Q_{\alpha \rho }^{(2)}(\stackrel{}{\lambda })\right\},$$ (4.60) where the first contribution is, again, the classical one given by (A.13) and the contributions $`Q_{\alpha \rho }^{(1)}(\stackrel{}{\lambda })`$ and $`Q_{\alpha \rho }^{(2)}(\stackrel{}{\lambda })`$ denote respectively first and second order quantum corrections whose explicit expressions are $`m^3Q_{\alpha \rho }^{(1)}(\stackrel{}{\lambda })={\displaystyle \frac{\mathrm{\Psi }_{g}^{}{}_{}{}^{2}h_g^{}8h_{g_0}^{}+4\stackrel{}{\alpha }^2}{8(\stackrel{}{\lambda }\stackrel{}{\alpha })^2}}\delta _{\alpha \rho }+{\displaystyle \frac{f^{\overline{\beta }\gamma \overline{\rho }}f^{\overline{\beta }\gamma \overline{\alpha }}}{4(\stackrel{}{\lambda }\stackrel{}{\alpha })(\stackrel{}{\lambda }\stackrel{}{\rho })}},`$ (4.61) $`m^3Q_{\alpha \rho }^{(2)}(\stackrel{}{\lambda })={\displaystyle \frac{5(\stackrel{}{\gamma }\stackrel{}{\beta })\left\{f^{\overline{\beta }\gamma \overline{\rho }}f^{\overline{\gamma }\beta \overline{\alpha }}+f^{\overline{\beta }\gamma \overline{\alpha }}f^{\overline{\gamma }\beta \overline{\rho }}\right\}}{48(\stackrel{}{\lambda }\stackrel{}{\alpha })(\stackrel{}{\lambda }\stackrel{}{\rho })}}+h_{g_0}^{}{\displaystyle \frac{\mathrm{\Psi }_{g}^{}{}_{}{}^{2}h_g^{}4h_{g_0}^{}+4\stackrel{}{\alpha }^2}{8(\stackrel{}{\lambda }\stackrel{}{\alpha })^2}}\delta _{\alpha \rho }+`$ $`+{\displaystyle \frac{f^{\overline{\beta }\gamma \overline{\rho }}f^{\overline{\beta }\gamma \overline{\alpha }}}{4(\stackrel{}{\lambda }\stackrel{}{\alpha })(\stackrel{}{\lambda }\stackrel{}{\rho })}}\left\{{\displaystyle \frac{\mathrm{\Psi }_{g}^{}{}_{}{}^{2}h_g^{}4h_{g_0}^{}}{4}}{\displaystyle \frac{\stackrel{}{\gamma }^2}{3}}{\displaystyle \frac{(\stackrel{}{\gamma }\stackrel{}{\beta })(\stackrel{}{\lambda }\stackrel{}{\beta })}{2\stackrel{}{\lambda }\stackrel{}{\gamma }}}\right\}`$ $`{\displaystyle \frac{2\stackrel{}{\rho }^4+(\stackrel{}{\beta }\stackrel{}{\rho })^2}{24(\stackrel{}{\lambda }\stackrel{}{\rho })^2}}\delta _{\alpha \rho }{\displaystyle \frac{f^{\beta \gamma \rho }f^{\beta \gamma \alpha }}{8\stackrel{}{\lambda }\stackrel{}{\gamma }}}\left\{{\displaystyle \frac{\stackrel{}{\gamma }\stackrel{}{\rho }}{\stackrel{}{\lambda }\stackrel{}{\rho }}}+{\displaystyle \frac{\stackrel{}{\gamma }\stackrel{}{\alpha }}{\stackrel{}{\lambda }\stackrel{}{\alpha }}}{\displaystyle \frac{\stackrel{}{\gamma }\stackrel{}{\beta }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}\right\}.`$ (4.62) The existence of this particular solution for eq. (4.36) means that, at least, there is a spin-4 quantum conserved density or, equivalently, a quantum conserved quantity of spin 3. Moreover, as explained below (4.12), there will be also a conserved quantity of spin $`3`$. Taking into account Parke’s results , this allows one to conclude that the Split models are quantum integrable and, hence, that they should admit a factorizable $`S`$-matrix. ## 5 The spectrum of the Split models. An important property of both the SSSG and the HSG theories is that they admit soliton solutions . Moreover, the semi-classical quantization of the solitons is expected to help in deducing the form of the exact $`S`$-matrix. The complete analysis of the soliton spectrum of the SSSG theories is beyond the scope of this paper and will be presented in subsequent publications. However, in this section we will discuss its main features and provide some explicit soliton solutions. We will restrict ourselves to the Split models, whose quantum integrability has been explicitly established in the previous section. They are associated with a type I symmetric space $`G/G_0`$ of maximal rank, i.e., $`\mathrm{rank}(G/G_0)=\mathrm{rank}(G)=r_g`$, and two regular elements $`\mathrm{\Lambda }_\pm g_1`$. This means that $`\mathrm{Ker}(\text{ad}_{\mathrm{\Lambda }_\pm })`$ is a Cartan subalgebra of $`g`$ contained in $`g_1`$, and we will use the realization of $`g_0`$ and $`g_1`$ given in the Appendix. Actually, the simplest Split model is the well known sine-Gordon model, which is recovered with the symmetric space $`SU(2)/SO(2)`$. In this case, $`G_0`$ is abelian and, hence, the coupling constant does not have to be quantized; this is the reason why this SSSG theory has not been mentioned in the tables 26. First of all, we have to obtain the vacuum manifold $`_0`$, which consists of the constant field configurations $`h_0`$ that minimise the potential in (2.3). This condition amounts to $$[\mathrm{\Lambda }_+,h_0^{}\mathrm{\Lambda }_{}h_0]=\mathrm{\hspace{0.25em}0},\mathrm{and}\stackrel{}{\alpha }\left(\mathrm{\Lambda }_+\right)\stackrel{}{\alpha }\left(h_0^{}\mathrm{\Lambda }_{}h_0\right)>0,$$ (5.1) for all roots $`\stackrel{}{\alpha }`$ of $`g`$, which requires that $`\mathrm{\Lambda }_+`$ and $`h_0^{}\mathrm{\Lambda }_{}h_0`$ belong to the same Weyl chamber of the Cartan subalgebra of $`g`$. Since they are regular, and taking into account that the conjugation $`h_0^{}h_0`$ permutes the Weyl chambers , we can assume that $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$ already belong to the same Weyl chamber without any loss of generality. Then (5.1) implies $`h_0^{}\mathrm{\Lambda }_{}h_0=\mathrm{\Lambda }_{}`$ and, therefore, $`h_0`$ has to be of the form $$h_0=\mathrm{e}^{\pi \stackrel{}{\mu }\stackrel{}{t}},$$ (5.2) where $`\stackrel{}{\mu }`$ either vanishes or belongs to the co-root lattice of $`G`$, which is the root lattice of the dual group $`G^{}`$, $`\mathrm{\Lambda }_R(G^{})`$. Recall that $`G^{}`$ is defined by requiring that its Lie algebra has roots which are the duals of the roots of $`g`$ defined by $`\stackrel{}{\alpha }^{}=2\stackrel{}{\alpha }/\stackrel{}{\alpha }^2`$. However, any element of the form (5.2) satisfies $`h_0^2=1`$, which implies that the vacuum manifold is given by $$_0=\{1,\mathrm{e}^{\pi \stackrel{}{\mu }\stackrel{}{t}}\stackrel{}{\mu }=\underset{i=1}{\overset{r_g}{}}n_i\stackrel{}{\alpha }_i^{},n_i=0,1\},$$ (5.3) where $`\stackrel{}{\alpha }_1,\mathrm{},\stackrel{}{\alpha }_{r_g}`$ are simple roots of $`g`$. Therefore, $`_0`$ is an abelian discrete group isomorphic to the coset $`\mathrm{\Lambda }_R(G^{})/2\mathrm{\Lambda }_R(G^{})`$. Moreover, using the normalization in the Appendix, one can check that $$\mathrm{e}^{\pi \stackrel{}{\alpha }^{}\stackrel{}{t}}=\mathrm{exp}\left(+\frac{2\pi t^\alpha }{\sqrt{\stackrel{}{\alpha }^2}}\right)=\mathrm{exp}\left(\frac{2\pi t^\alpha }{\sqrt{\stackrel{}{\alpha }^2}}\right),$$ (5.4) for any root of $`g`$, which emphasises that $`_0G_0`$. ### 5.1 Fundamental particles. They correspond to the fluctuations of the field $`h`$ around a vacuum configuration $`h_0_0`$. Let us take, $`h=h_0e^\mathrm{\Phi }`$, where $`\mathrm{\Phi }g_0`$. The linearized eqs. (2.6) are $$P(\mathrm{\Phi })=\mathrm{\hspace{0.17em}0}\mathrm{and}_\mu ^\mu \mathrm{\Phi }=4m^2[\mathrm{\Lambda }_+,[\mathrm{\Lambda }_{},\mathrm{\Phi }]],$$ (5.5) which show that the fundamental particles are associated with the non-vanishing eigenvalues of the mass-matrix $`4m^2[\mathrm{\Lambda }_+,[\mathrm{\Lambda }_{},\mathrm{\Phi }]]`$ on $`g_0`$. They correspond to the field configurations of the form $`\mathrm{\Phi }=\varphi t^\alpha `$, where $`\stackrel{}{\alpha }`$ is an arbitrary positive root of $`g`$, whose mass is $$m_\alpha =\mathrm{\hspace{0.25em}2}m\sqrt{\stackrel{}{\alpha }\left(\mathrm{\Lambda }_+\right)\stackrel{}{\alpha }\left(\mathrm{\Lambda }_{}\right)},$$ (5.6) which is real because of (5.1). Therefore, for each positive root $`\stackrel{}{\alpha }`$ of $`g`$, there is a fundamental particle described by a real field $`\varphi =\varphi (x,t)`$ whose mass is given by (5.6). It is worth noticing that (5.6) is the mass formula giving the spectrum of fundamental particles of the HSG theory associated with $`G`$; however, in that case the particles are described by complex fields. A very peculiar property of this spectrum that is shared with the HSG and all the other SSSG theories is that the mass formula (5.6) satisfies the following kind of inequalities $$m_{\alpha +\beta }m_\alpha +m_\beta ,$$ (5.7) which suggests that some of the fundamental particles might be unstable. For the HSG theories, this has been checked using perturbation theory and, consequently, their exact $`S`$-matrix exhibits resonance poles associated with the unstable particles . For the Split model related to $`SU(3)/SO(3)`$ it has also been checked that some of the fundamental particles actually decay ; however, more work is needed in order to establish this for the generality of the SSSG models and construct their exact $`S`$-matrix. ### 5.2 Soliton solutions. In order for a solution to eqs. (2.6) to have finite energy the field $`h`$ must tend to limits in $`_0`$ as $`x\pm \mathrm{}`$. So, $$h(+\mathrm{},t)h^{}(\mathrm{},t)_0,$$ (5.8) and its value is conserved as the system evolves in time. This means that, at fixed $`t`$, a solution $`h=h(t,x)`$ with finite energy is a path on the $`G_0`$ manifold connecting two elements in $`_0`$ and, since $`G_0`$ is not simply connected in general, there could be different solutions sharing the same value of $`h(+\mathrm{},t)h^{}(\mathrm{},t)`$. Therefore, each solution will be characterized by two topological ‘quantum numbers’. The first will be the value of (5.8), which is an element in $`_0`$ or, equivalently, in the discrete group $`\mathrm{\Lambda }_R(G^{})/2\mathrm{\Lambda }_R(G^{})`$, and the second will be an element in $`\pi _1(G_0)`$, the fundamental group of $`G_0`$, which can be found in table 1. In other words, the solitons of the Split models will be topological, like the solitons of the sine-Gordon equation or of the affine Toda equations with imaginary coupling constant. This is in contrast with the solitons of the HSG theories, which only carry Noether charges. On their side, the solitons of the Split models do not carry any Noether charge because $`g_0^0=\{0\}`$ in (2.2) ($`p=0`$). Nevertheless, for a generic SSSG theory $`p0`$ and the solitons will carry both topological and $`U(1)^p`$ Noether charges, which make them similar to the dyons in four-dimensional non-abelian gauge theories Taking into account that the sine-Gordon theory is the Split model corresponding to $`SU(2)/SO(2)`$, a number of explicit soliton solutions for the Split models can be obtained by embedding the sine-Gordon solitons in the regular $`SU(2)`$ subgroups of $`G`$. This method is widely used in the context of Yang-Mills theories based on arbitrary Lie groups to construct monopole or instanton solutions by embeddings of the $`SU(2)`$ spherically symmetric ’t-Hooft-Polyakov monopole or the self-dual $`SU(2)`$ Belavin-Polyakov-Schwartz-Tyupkin instanton . It has also been used to construct the soliton solutions of the affine Toda theories with imaginary coupling constant and, more recently, to construct the soliton solutions of the HSG theories starting with the Complex sine-Gordon solitons . For each positive root $`\stackrel{}{\alpha }`$ of $`G`$, let us consider the field configuration $`h=\mathrm{exp}(\varphi t^\alpha /\sqrt{\stackrel{}{\alpha }^2})`$, which trivially satisfies the constraints in (2.6). According to (2.3), its Lagrangian density is $$=\frac{1}{4\pi \beta ^2\stackrel{}{\alpha }^2}\left(\frac{1}{2}_\mu \varphi ^\mu \varphi +m_\alpha ^2\left(\mathrm{cos}\varphi \mathrm{\hspace{0.17em}1}\right)\right),$$ (5.9) which is just the Lagrangian density of the sine-Gordon model. This way, for each positive root $`\stackrel{}{\alpha }`$ of $`G`$, each soliton solution of the sine-Gordon equation provides a soliton solution for the Split model. Namely, the usual soliton and anti-soliton solutions allow one to construct two a priori different soliton solutions with mass $$M_{s,\overline{s}}(\stackrel{}{\alpha })=\frac{2}{\stackrel{}{\alpha }^2}\frac{1}{\pi \beta ^2}m_\alpha ,$$ (5.10) while the masses of the solitons associated with the breathers of the sine-Gordon equation are $$M_n(\stackrel{}{\alpha })=\mathrm{\hspace{0.17em}2}M_s(\stackrel{}{\alpha })\mathrm{sin}\left(\frac{\stackrel{}{\alpha }^2}{2}\frac{\pi \beta ^2}{2}n\right)\mathrm{\hspace{0.17em}2}M_s(\stackrel{}{\alpha }),$$ (5.11) and there is a different soliton for each integer $`n`$ such that $`n<2/(\beta ^2\stackrel{}{\alpha }^2)`$. As usual, in the weak coupling $`\beta ^20`$ limit we obtain $$M_n(\stackrel{}{\alpha })nm_\alpha 𝒪\left(\beta ^2\right)$$ (5.12) and, hence, the fundamental particle associated with $`\stackrel{}{\alpha }`$ becomes identified with the lightest breather with mass $`M_1(\stackrel{}{\alpha })`$. It is important to notice that the mass formulae (5.10) and (5.11) satisfy inequalities similar to (5.7), which again suggests that some of these solitons might be unstable. Actually, some examples of unstable solitons are already known in the Split model related to $`SU(3)/SO(3)`$ . The soliton ($`s`$) and antisoliton ($`\overline{s}`$) solutions of the sine-Gordon model satisfy $$\varphi _{s,\overline{s}}(+\mathrm{},t)\varphi _{s,\overline{s}}(\mathrm{},t)=\pm 2\pi ,$$ (5.13) which implies the following asymptotic behaviour for the fields $`h_s=\mathrm{exp}(\varphi _st^\alpha /\sqrt{\stackrel{}{\alpha }^2})`$ and $`h_{\overline{s}}=\mathrm{exp}(\varphi _{\overline{s}}t^\alpha /\sqrt{\stackrel{}{\alpha }^2})`$: $$h_s(+\mathrm{},t)h_s^{}(\mathrm{},t)=\mathrm{e}^{+\frac{2\pi t^\alpha }{\sqrt{\stackrel{}{\alpha }^2}}},h_{\overline{s}}(+\mathrm{},t)h_{\overline{s}}^{}(\mathrm{},t)=\mathrm{e}^{\frac{2\pi t^\alpha }{\sqrt{\stackrel{}{\alpha }^2}}},$$ (5.14) both in $`_0`$. Therefore, taking into account (5.4), their asymptotic behaviour is the same and it will not be possible to distinguish these two soliton configurations unless $`G_0`$ is not simply connected and $`h_s`$ and $`h_{\overline{s}}`$ are associated with different elements in $`\pi _1(G_0)`$. ## 6 Conclusions. In this paper we have studied some of the quantum properties of the massive SSSG theories constructed in . First of all, we have identified the perturbed conformal field theories corresponding to these theories when the symmetric space $`G/G_0`$ is of type I. This amounts to find which are the unperturbed CFT and the perturbing operator specified by the potential term in the classical action. Since the type I symmetric spaces are irreducible, the perturbation is given by a single spinless primary field whose conformal dimension has been calculated and can be found in tables 2 and 3. Actually, our calculation only depends on the algebraic structure of the symmetric space and, therefore, it provides the conformal dimension of the perturbations for the general class of SSSG theories constructed in . They are obtained from ours by substituting $`\mathrm{\Lambda }_+`$ and $`\mathrm{\Lambda }_{}`$ for two arbitrary elements $`T`$ and $`\overline{T}`$ in $`g_1`$, and $`g_0^0`$ for h, their simultaneous centralizer in $`g_0`$. The resulting SSSG theories are perturbations of the coset CFT related to $`G/H`$, where $`H`$ is the Lie group corresponding to h. However, as shown in , these theories will not exhibit a mass gap unless $`H`$ is either trivial or abelian. Among others, the resulting class of perturbed CFT’s include massive perturbations of WZNW models, which are related to symmetric spaces of maximal rank; we have named these theories ‘Split models’. In addition, there are new massive perturbations of parafermion theories different to those provided by the HSG theories . In particular, this class of theories includes the perturbations of the simplest $`\text{}_k`$ parafermions by their first and second thermal operators , which are related to $`G/G_0=Sp(2)/U(2)`$ and $`SU(3)/SO(3)`$ , respectively. Our second task was to investigate the quantum integrability of the SSSG theories. In view of the large variety of different types of perturbed CFT’s corresponding to the SSSG theories, we have restricted ourselves to give a detailed proof of the quantum integrability only for the Split models. Classically, they exhibit $`\mathrm{rank}(G)`$ conserved quantities for each odd spin $`\pm 1,\pm 3,\mathrm{}`$, and we have checked that there are at least two quantities of spin $`+3`$ and $`3`$ that remain conserved in the quantum theory after an appropriate renormalization. This implies, via the usual folklore, the factorization of their scattering matrices and, hence, their quantum integrability. This result, together with the integrability of the HSG theories whose simplest higher spin conserved quantities have spin $`\pm 2`$, lead us to conjecture that all the massive SSSG theories will be quantum integrable. The quantum integrability of the SSSG theories implies that they should admit a factorizable $`S`$-matrix and the next stage of analysis consists in establishing its form. As a first step towards this aim, we have illustrated the general properties of the spectrum of the SSSG theories by discussing the spectrum of fundamental particles and solitons of the Split models. Its main features, which are expected to be shared with all the other SSSG theories, are, first, that the fundamental particles become identified with some of the solitons in the semiclassical limit, like in the sine-Gordon and complex sine-Gordon theories. This makes us expect that the spectrum will be entirely solitonic in the general case. Second, since there are different vacua, these solitons are topological, in contrast with the solitons of the HSG theories which are not. Moreover, in the general case, they are expected to carry conserved Noether charges as well. Third, some of these solitons are expected to correspond to unstable particles in the quantum theory. Therefore, like in the HSG theories, only the stable solitons should correspond to asymptotic states while the unstable ones will produce resonance poles in the $`S`$-matrix . Finally, both the fundamental particles and the solitons are labelled by the roots of the Lie algebra of $`G`$, which is reminiscent of the spectrum of the HSG theories. Taking into account the form of the exact $`S`$-matrices of the HSG theories , the last feature suggests that the $`S`$-matrices of the SSSG theories could be somehow related to the ‘colour valued’ scattering matrices constructed in . These $`S`$-matrices are related to pairs $`\{\stackrel{~}{g}\|g\}`$ of simply laced Lie algebras, where $`\stackrel{~}{g}`$ governs the mass spectrum and the fusing rules, while $`g`$ provides the ‘colour’ quantum numbers. Using this construction, the $`S`$-matrix of the HSG theory corresponding to the Lie group $`G`$ at level $`k`$ is related to the pair $`\{A_{k1}\|g\}`$, where $`g`$ is the Lie algebra of $`G`$. Then, it is worthwhile to notice that the conjectured central charge of the ultraviolet CFT corresponding to the colour valued $`S`$-matrix specified by the pair $`\{D_k\|g\}`$ coincides with the central charge of the unperturbed CFT of the $`G/G_0`$ Split model ($`p=0`$) at level $`k`$, or level $`2k`$ if $`G=SU(m)`$, where $`G/G_0`$ is one of the symmetric spaces listed in (4.1) with simply laced $`G`$. However, additional work is needed in order to go beyond this numerical coincidence. Acknowledgments We would like to thank J. Sánchez Guillén, A. Fring and C. Korff for their valuable comments. This research is supported partially by CICYT (AEN99-0589), DGICYT (PB96-0960), and the EC Commission via a TMR Grant (FMRX-CT96-0012). ## Appendix A Classical conserved densities of the Split models. In this appendix we give the explicit expressions for the classical conserved densities with scale dimension 2 and 4 corresponding to the Split models, which have been obtained by solving eqs. (2.6). We will use the standard explicit realization of the basis $`t^a`$ of $`g`$ in terms of a Cartan basis of its complexification $`t^A=iH^A,A=\mathrm{\hspace{0.25em}1},\mathrm{},\mathrm{rank}(g),`$ 0.2truecm (A.1) $`t^\alpha =\sqrt{{\displaystyle \frac{1}{2E_\alpha ,E_\alpha }}}(E_\alpha E_\alpha ),t^{\overline{\alpha }}=i\sqrt{{\displaystyle \frac{1}{2E_\alpha ,E_\alpha }}}(E_\alpha +E_\alpha ),`$ where the $`H^A`$’s provide an orthonormal basis for the Cartan subalgebra with respect to the invariant bilinear form $`,`$, and the step operators $`E_\alpha `$ are normalized such that $$[H^A,E_\alpha ]=\stackrel{}{\alpha }(H^A)E_\alpha =\alpha ^AE_\alpha ,[E_\alpha ,E_\alpha ]=E_\alpha ,E_\alpha \alpha ^AH^A,$$ (A.2) with $`\stackrel{}{\alpha }`$ a positive root of $`g`$. The generators (A.1) are normalized in such a way that $`t^a,t^b=\delta ^{ab}`$. Moreover, the corresponding structure functions are totally antisymmetric and satisfy: $`f^{\overline{\alpha }\overline{\beta }\overline{\gamma }}=f^{\alpha \beta \overline{\gamma }}=f^{\alpha \overline{\alpha }\beta }=f^{\alpha \overline{\alpha }\overline{\beta }}=0,\stackrel{}{\alpha },\stackrel{}{\beta },\stackrel{}{\gamma }>0,`$ $`f^{ABC}=f^{AB\alpha }=f^{AB\overline{\alpha }}=0,\stackrel{}{\alpha }>0\mathrm{and}A,B=1,\mathrm{},\mathrm{rank}(g),`$ $`f^{\alpha \overline{\alpha }A}=\alpha ^A,\stackrel{}{\alpha }>0\mathrm{and}A=1,\mathrm{},\mathrm{rank}(g).`$ (A.3) Recall that the rank of a symmetric space $`G/G_0`$ is the dimension of any maximal abelian subspace contained in $`g_1`$. Therefore, if $`\mathrm{rank}(G/G_0)=\mathrm{rank}(G)`$, one can take $`t^Ag_1`$ for all $`A`$ and, hence, the subset of generators $`t^\alpha `$ provide a basis for the subalgebra $`g_0`$, while the $`t^{\overline{\alpha }}`$’s, together with the $`t^A`$’s, generate $`g_1`$. We will use this realization of $`g_0`$ and $`g_1`$ in our calculations for the Split models. Therefore, $$\mathrm{\Lambda }_{}=\lambda ^At^A=\stackrel{}{\lambda }\stackrel{}{t},$$ (A.4) and the potential $`q=q(x)`$ in the Lax operator (4.5) is of the form $`q=q^\alpha t^\alpha `$. Taking into account all this, the expressions for the conserved densities of scale-dimension 2 are $$\mu ^A_2^{(0)A}=\stackrel{}{\mu }\stackrel{}{}_2^{(0)}=D_{\alpha \beta }q^\alpha q^\beta ,\mathrm{with}D_{\alpha \beta }=\frac{1}{2m}\frac{\stackrel{}{\mu }\stackrel{}{\alpha }}{\stackrel{}{\lambda }\stackrel{}{\alpha }},$$ (A.5) where $`\stackrel{}{\mu }`$ is an arbitrary $`\mathrm{rank}(g)`$-component vector which allows one to write the $`\mathrm{rank}(g)`$ conserved quantities $`_2^{(0)A}`$ in a compact way. In particular, for $`\stackrel{}{\mu }=\stackrel{}{\lambda }`$, (A.5) reduces to $$m\lambda ^A_2^{(0)A}=\frac{1}{2}q^\alpha q^\alpha =\mathrm{\hspace{0.17em}2}\pi \beta ^2T_{},$$ (A.6) which is one of the components of the classical stress-energy tensor. For scale dimension 4, the $`\mathrm{rank}(g)`$-conserved densities can be written as: $$\stackrel{}{\mu }\stackrel{}{}_4^{(0)}=R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\mu })q^\alpha q^\beta q^\gamma q^\rho +P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\mu })q^\alpha q^\beta _{}q^\gamma +Q_{\alpha \beta }^{(0)}(\stackrel{}{\mu })q^\alpha _{}^2q^\beta ,$$ (A.7) where $`R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\mu })=({\displaystyle \frac{(\stackrel{}{\alpha }\stackrel{}{\gamma })(\stackrel{}{\gamma }\stackrel{}{\mu })}{8m^3(\stackrel{}{\lambda }\stackrel{}{\alpha })(\stackrel{}{\lambda }\stackrel{}{\gamma })^2}}\delta _{\alpha \beta }\delta _{\gamma \rho }{\displaystyle \frac{f^{\overline{\beta }\gamma \overline{\xi }}f^{\xi \rho \alpha }(\stackrel{}{\alpha }\stackrel{}{\mu })}{6m^3(\stackrel{}{\lambda }\stackrel{}{\beta })(\stackrel{}{\lambda }\stackrel{}{\xi })(\stackrel{}{\lambda }\stackrel{}{\alpha })}}+`$ $`+{\displaystyle \frac{f^{\overline{\rho }\alpha \overline{\xi }}f^{\overline{\xi }\overline{\gamma }\beta }}{24m^3(\stackrel{}{\lambda }\stackrel{}{\gamma })(\stackrel{}{\lambda }\stackrel{}{\rho })}}\{{\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}+{\displaystyle \frac{\stackrel{}{\xi }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\xi }}}\})_{\mathrm{sym}(\alpha \beta \gamma \rho )},`$ (A.8) $`P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\mu })=({\displaystyle \frac{f^{\gamma \alpha \beta }}{3m^3(\stackrel{}{\lambda }\stackrel{}{\gamma })^2}}{\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}+{\displaystyle \frac{f^{\overline{\alpha }\gamma \overline{\beta }}}{4m^3(\stackrel{}{\lambda }\stackrel{}{\alpha })(\stackrel{}{\lambda }\stackrel{}{\beta })}}{\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}`$ $`{\displaystyle \frac{f^{\overline{\beta }\alpha \overline{\gamma }}}{4m^3(\stackrel{}{\lambda }\stackrel{}{\gamma })(\stackrel{}{\lambda }\stackrel{}{\beta })}}\{{\displaystyle \frac{\stackrel{}{\gamma }\stackrel{}{\mu }}{\mathrm{\hspace{0.17em}3}\stackrel{}{\lambda }\stackrel{}{\gamma }}}+{\displaystyle \frac{\stackrel{}{\beta }\stackrel{}{\mu }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}\})_{\mathrm{sym}(\alpha \beta )},`$ (A.9) $`Q_{\alpha \beta }^{(0)}(\stackrel{}{\mu })={\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\mu }}{\mathrm{\hspace{0.17em}2}m^3(\stackrel{}{\lambda }\stackrel{}{\alpha })^3}}\delta _{\alpha \beta }.`$ (A.10) In these equations, $`\mathrm{sym}(\alpha \beta \gamma \rho )`$ and $`\mathrm{sym}(\alpha \beta )`$ means that the corresponding expressions have to be considered completely symmetrized in $`(\alpha \beta \gamma \rho )`$ or $`(\alpha \beta )`$, respectively. In particular, for $`\stackrel{}{\mu }=\stackrel{}{\lambda }`$ these expressions simplify to $`R_{\alpha \beta \gamma \rho }^{(0)}(\stackrel{}{\lambda })={\displaystyle \frac{1}{24m^3\stackrel{}{\lambda }\stackrel{}{\alpha }}}\left\{{\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\gamma }}{\stackrel{}{\lambda }\stackrel{}{\gamma }}}\left\{\delta _{\alpha \beta }\delta _{\gamma \rho }+\delta _{\alpha \rho }\delta _{\gamma \beta }\right\}+{\displaystyle \frac{\stackrel{}{\alpha }\stackrel{}{\beta }}{\stackrel{}{\lambda }\stackrel{}{\beta }}}\delta _{\alpha \gamma }\delta _{\beta \rho }\right\},`$ (A.11) $`P_{\alpha \beta \gamma }^{(0)}(\stackrel{}{\lambda })={\displaystyle \frac{1}{6m^3\stackrel{}{\lambda }\stackrel{}{\gamma }}}\left\{{\displaystyle \frac{f^{\overline{\beta }\alpha \overline{\gamma }}}{\stackrel{}{\lambda }\stackrel{}{\beta }}}+{\displaystyle \frac{f^{\overline{\alpha }\beta \overline{\gamma }}}{\stackrel{}{\lambda }\stackrel{}{\alpha }}}\right\},`$ (A.12) $`Q_{\alpha \beta }^{(0)}(\stackrel{}{\lambda })={\displaystyle \frac{\delta _{\alpha \beta }}{\mathrm{\hspace{0.17em}2}m^3(\stackrel{}{\lambda }\stackrel{}{\alpha })^2}}.`$ (A.13)
warning/0002/astro-ph0002448.html	ar5iv	text	# Observing high-redshift Supernovae in lensed galaxies ## 1 Introduction Observations of distant sources with known absolute luminosity (cosmological standard candles) are of primary importance to modern cosmology, since the relation between the apparent magnitude, luminosity and redshift of distant galaxies can be used to determine the Hubble constant $`H_0`$, the deceleration (or density) parameter $`q_0`$, and the cosmological constant $`\mathrm{\Lambda }`$. Observations of standard candles beyond $`z=0.05`$ (where peculiar velocities are small) can yield the value of $`H_0`$ with reasonably small uncertainty (e.g., Filippenko (1996); Hamuy et al. (1996)). Observations of standard candles at $`z>0.3`$ are being used for the determination of the fraction of the total energy of the Universe in matter $`\mathrm{\Omega }_M`$ and in some hitherto unknown form $`\mathrm{\Omega }_\mathrm{\Lambda }`$ (Riess et al. 1998b ; Perlmutter et al. (1999); Saini et al. (2000)). The study of the gravitational magnification of standard candles at even higher redshift will put tighter constraints on dark matter models of cosmogony (e.g., Kolatt & Bartelmann (1998); Marri & Ferrara (1998); Holz (1998); Metcalf (1999); Porciani & Madau (2000)). The work of Riess et al. (1998b) and Perlmutter et al. (1999) has shown that, if detected significantly earlier than the epoch of their peak luminosity, type Ia supernovae (SNe Ia) would be the most useful among cosmological candles at high redshift. However, the required integration times for good photometry and for obtaining spectra of such supernovae at redshift $`z1`$ are estimated to be tens of hours on a 10m telescope for $`0\stackrel{}{.}75`$ seeing (Goobar & Perlmutter (1995)). These observations would clearly be more favourable if these supernovae occur in galaxies magnified by gravitational lensing. The magnification due to lensing can be significant enough to make possible the detection of supernovae (SNe) in galaxies at high redshifts ($`z\mathrm{}>1`$). Narasimha & Chitre (1988) first pointed out that such events in giant luminous arcs (as in the A370 system) can be used as a test of the lens models. In the case of multiply imaged supernovae, Kovner & Paczyński (1988) deduce simple relations between the magnification of such a SN, the separation of images, and the differences between the arrival times of the event in different images. Indeed, such SNe would serve as a unique probe for not only the distribution of matter in the clusters, but also for studying the source galaxies themselves. Due to the increased flux produced by the magnification of the images, photometric and spectroscopic studies of very distant galaxies can become possible. This would enable us to obtain information, which would be otherwise unavailable, about the star formation process in the young galaxies (Mellier et al. (1991); Yee & de Robertis (1991)), the evolutionary status of AGN (Stickel et al. (1991)), and even the morphology of distant galaxies (Colley, Tyson & Turner (1996)). Indeed, one of the farthest known galaxies (at $`z`$=4.92, Franx et al. (1997)) would not have been detected had it not been for the $`10`$-fold magnification by the cluster CL1358+62 at $`z=0.33`$. In this paper we address the feasibility of detecting lensed SN events in high redshift galaxies which would be useful in the measurement of cosmological parameters. From a qualitative point of view such a study seems worthwhile for several reasons. For a typical magnification of 3–4 mag (Kovner & Paczyński (1988)) the study of lensed SNe stretches the usefulness of using them to characterize the distance ladder to further distances by a factor of 4–6, or, equivalently, results in a considerable decrease in the required duration of observation. Furthermore, although galaxies lensed into arcs are resolved in one direction due to stretching, a supernova in such a galaxy will remain a point source, hence the signal-to-noise ratio (SNR) of a lensed supernova in an arc will be superior to that of one in an unlensed galaxy. Finally, a cluster lens typically produces multiple images with time delays between them being up to several months, thereby making it possible to observe the same SN again, and measuring its light curve more accurately, particularly in its pre-peak phase. In the searches we propose, we do not have to be confined to one lensed SN at a time. In many known cases of gravitational lensing of background objects by galaxy clusters, several arcs and arclets can be found in an area of the sky typically imaged by a single CCD frame. In the case of Abell 2218, for instance, there are 30 observed arclets (Ebbels et al. (1998); Bézecourt et al. (1998)) with $`R23.5`$ and $`\mu _R25.5`$ between $`z=`$0.5 to 1.5, so clearly a lot of galaxies can be simultaneously monitored. This is also true of the cluster Abell~2390 at $`z=0.23`$, in which, in addition to the famous “straight arc” (triple image of a galaxy at $`z=0.913`$), there are at least 12 arclets ($`R<21`$) between $`z=`$0.4–1.3 in an area of $`2.7\times 2.7`$ arcmin<sup>2</sup> around it (Bézecourt & Soucail (1997)). In the same area, in the magnitude range $`21<R<23.5`$, there are four images of two galaxies at $`z=4.05`$ (Pelló et al. (1998); Frye and Broadhurst (1998)) lensed by the same cluster. Similarly, in a single HST WFPC field, which covers much less area ($``$5 arcmin<sup>2</sup>) than most CCD cameras on terrestrial telescopes, one finds 20 arclets brighter than $`m_{\mathrm{F675W}}`$=24 (corresponding to at least 15 independent galaxies) beyond $`z=0.7`$ in the cluster Abell~370 ($`z=0.37`$, Bézecourt et al. (1999)), and 14 arclets corresponding to 5 galaxies between $`z=`$0.5–1.5 in the cluster MS0440+0204 at $`z=0.197`$ (Gioia et al. (1998)). Not all of these, of course, would be magnified to the same degree, but on average they would be magnified, making them easier to be observed than in an unlensed case. Hereafter, we briefly review the usefulness of SN Ia and SN IIL in the determination of cosmological parameters, and in §3 quantify the effect of a cluster lens on the detectability of high-$`z`$ SNe. We illustrate our case with the cluster lens Abell 2218. In §4, we show that, in addition to this, for SNe detected in giant arcs, the signal-to-noise ratio for the detection of SNe is enhanced by an order of magnitude due to lensing. We summarize in §5 in the light of ongoing SN searches. ## 2 Standard Bombs: Supernovae of Type Ia & IIL Two subgroups of SNe seem to be relevant for cosmological use: SNe Ia and SNe IIL. SNe II are more frequent (by a factor of $``$4, van den Bergh & Tammann (1991)) in late spirals (Sbc-Sd), which are the most numerous among field galaxies. It is more likely that supernovae at $`z>1`$, where the dependence of the luminosity distance on models is the most sensitive (Fig. 1), will be of type-II (e.g., Madau et al. 1998). Though SNe II maxima in general are known to have a large spread in luminosities, about half of them (those that have “linear” light curves, SNe IIL, Young & Branch (1989); Cappellaro et al. (1997)) represent very good standard bombs (Gaskell (1992)), since they have a small intrinsic scatter ($`\sigma =0.3`$ mag) around their peak magnitude ($`M_B=17.05`$, Miller & Branch (1990), $`H_0=75`$). However, SNe Ia are much more luminous at their peak ($`M_B=18.95`$ for $`H_0=75`$), with a smaller scatter in peak magnitude, if corrected for the slope of their light curve (Perlmutter et al. (1999); Riess et al. 1998b ). The major problem in the use of high redshift SNe as standard candles lies in the identification of the type of the SN, and its photometric calibration. The identification of the type of the SN depends much on the shape of the light curve, which at high redshifts will be time dilated, making it easier to determine its shape. The magnified flux will render it easier to obtain a spectroscopic identification as well. It is therefore obvious that due to the considerable magnification, lensed SNe will be far more suitable candidates than unlensed SNe at the same redshift. ## 3 The Effect of Lens Magnification on detectability ### 3.1 Gravitational magnification We summarize briefly the essential results needed for this paper. Excellent reviews of gravitational lensing can be found elsewhere (e.g., Blandford & Narayan (1986); Schneider et al. (1992)). The basic equation which relates the angular coordinates of the source ($`\beta _1,\beta _2`$) to those of the image ($`\theta _1,\theta _2`$) is $$\beta =\theta \psi (\theta ).$$ (1) The dimensionless relativistic lens potential satisfies the two dimensional Poisson equation $`^2\psi (\theta )=2\kappa (\theta )`$, the convergence $`\kappa (\theta )=\mathrm{\Sigma }(\theta )/\mathrm{\Sigma }_{cr}`$, $`\mathrm{\Sigma }(\theta )`$ being the two dimensional surface mass density of the lens, and $`\mathrm{\Sigma }_{cr}=(c^2/4\pi G)(d_s/d_ld_{ls})`$ is the critical density. Here the distance between the observer and the source, that between the observer and the lens and that between the lens and the source are $`d_s`$, $`d_l`$ and $`d_{ls}`$ respectively. Gravitational lensing preserves the surface brightness of the light rays. The flux of light received by an observer is directly proportional to the solid angle subtended by the image at the observer. Since the solid angle of the image after lensing is, in general, different from that of the source the observer can receive more (or less) flux than in the unlensed situation. Thus galaxy clusters can act as gravitational telescopes by collecting light from the distant galaxies over a large area and sending it in our direction. The shape and size of the image are related to that of the source by the transformation matrix $`M_{ij}^1=\beta _i/\theta _j`$. This matrix is generally written in the form $$M^1=\left(\begin{array}{cc}1\kappa \gamma _1& \gamma _2\\ \gamma _2& 1\kappa +\gamma _1\end{array}\right),$$ (2) where $`\kappa `$ is the usual convergence and $`\gamma _1=1/2(^2\psi /\theta _1^2^2\psi /\theta _2^2),\gamma _2=^2\psi /\theta _1\theta _2`$ are the components of the shear. The magnification for a point source is given by the Jacobian of the inverse mapping $`f:\beta \theta `$, which is, in general, one to many. From equation (2) we find that the magnification $$\mu \mathrm{det}[M]=1/((1\kappa )^2\gamma ^2),$$ (3) where $`\gamma ^2=\gamma _1^2+\gamma _2^2`$, and is therefore different for different images. The set of all the points where det$`[M^1]`$ vanishes (singular points) in the source plane is called the caustic set, and the images of the caustic set are called the critical curves. A source of finite size (i.e., a galaxy) close to the caustic produces magnified images near the critical curves. Any point source within these galaxies will also be substantially magnified. For well-studied cluster lenses which have $`\mathrm{}>`$10 arclets, the model parameters of the lens can be constrained well enough so that the magnification of a SN occurring at any point on an arclet can be estimated reliably to an accuracy of $`\mathrm{}<0.5`$ mag. In this paper, we suggest the monitoring of such well-modelled lenses to look for cosmologically significant high-$`z`$ SNe. ### 3.2 The enhancement of detectable events A SN event at a redshift $`z1`$ is expected to have an apparent blue magnitude $`m_B25`$. If we are monitoring a system of arclets with an observing setup (telescope and detector) of limiting magnitude of detection $`m_{\mathrm{lim}}`$ (given an acceptable value of S/N), then we would like to estimate the probability that the SNe might be magnified by an amount $`\mathrm{\Delta }m=m_Bm_{\mathrm{lim}}`$. Given this probability and the SNe rates, one can obtain an estimate of the expected number of detectable events. Consider an area of the sky of a few arcmin<sup>2</sup> around the centre of a cluster at redshift $`z_L`$ being monitored with an array of CCDs. Within this region, let the number of SNe occurring per year in galaxies between redshifts $`z_s`$ and $`z_s+dz_s`$ be $$dN=N_0𝒩(z_s)dz_s$$ (4) where $$_{z_L}^{z_{\mathrm{max}}}𝒩(z_s)𝑑z_s=1.$$ If $`m_{\mathrm{lim}}`$ is the limiting magnitude, then the threshold magnification at a redshift $`z_s`$ for the source of unlensed magnitude $`m_s`$ to be detected will be $`\mu _0(z_s)=10^{(m_sm_{\mathrm{lim}})/2.5}`$. The number of these SNe being magnified by a ratio $`\mu >\mu _0(z_s)`$ is $$dN^{}=N_0𝒩(z_s)P(\mu >\mu _0(z_s))dz_s.$$ (5) where $`P(\mu >\mu _0(z_s))=`$ (6) $`{\displaystyle \frac{1}{\pi \beta _0^2}}{\displaystyle \frac{\mathrm{\Theta }[\|\mu (\theta )\|\mu _0(z_s)]\mathrm{\Theta }[\beta _0^2\beta (\theta )^2]}{\|\mu (\theta )\|}d^2\theta },`$ for a source at redshift $`z_s`$. Here $`\mathrm{\Theta }`$ is the Heaviside step function, $`\beta _0(z_s)`$ is the radius of the source (assumed circular) and the integral is performed over the field of the observed image. Since a single source can produce more than one image, the value of this quantity, for a given threshold magnification $`\mu _0`$, can be greater than one. Here we would quantify the enhancement in the detection of distant SNe by $$\mathrm{\Phi }_L(z)=_z^{z_{\mathrm{max}}}𝒩(z_s)P(\mu >\mu _0(z_s))𝑑z_s.$$ (7) This function represents the cumulative fraction of SNe that are observed, given the limiting magnitude $`m_{\mathrm{lim}}`$ of the observational setup, of the total number of SNe that occur between $`z`$ and $`z_{\mathrm{max}}`$ in the area of the sky that is being monitored. This depends upon the number density and redshift distribution of the host galaxies and the frequency of SN as a function of redshift. This should be compared with the quantity $$\mathrm{\Phi }_U(z)=_z^{z_{\mathrm{max}}}𝒩(z_s)\mathrm{\Theta }(m_{\mathrm{lim}}m(z_s))𝑑z_s$$ (8) which is the corresponding fraction that would be observed in the absence of the lens. For instance, since we assume the peak magnitude of a Type Ia SN to be $`m_B=25`$ at $`z=1`$, if $`m_{\mathrm{lim}}=25`$, the value of $`\mathrm{\Phi }_U(z>1)`$ will be zero whereas due to lensing $`\mathrm{\Phi }_L(z)`$ can be finite till $`z=z_{max}`$. ### 3.3 Example: the case of Abell 2218 To estimate the typical fraction of SNe which would be seen behind a cluster with a certain magnification, we consider the case of the well-studied cluster lens Abell 2218 ($`z=0.175`$), for which good published mass models exist. Here we use the model of Kneib et al. (1996), where the bimodal mass distribution is represented by two cluster-scale clumps of dark matter centred on the two brightest elliptical galaxies, their potentials modelled by the difference of two pseudo-isothermal elliptical mass distributions (PIEMDs), with an external truncation radius. In addition, the small-scale mass structure is represented by galaxy-sized lenses corresponding to the 34 brightest galaxies belonging to the cluster, modelled by similar functions with the appropriate parameters (velocity dispersion, core radius and truncation radius) scaled to the observed luminosities of the galaxies. There are 258 background galaxies in the magnitude range $`21.5<R<25`$, detected in the HST WFPC image of total area of 4.7 arcmin<sup>2</sup> analyzed by Kneib et al.. Of these, 35 spectroscopic redshifts are known, and another 18 redshifts are estimated in Ebbels et al. (1998). For other galaxies, random redshifts were chosen in the range $`0.175<z<2.5`$, from a distribution that conserves the number of galaxies per unit comoving volume in a flat, matter dominated universe. The redshift dependence of the frequencies $`𝒩(z_s)`$ of both SN Ia and SN II are taken from Madau et al. (1998), where the evolution of cosmic supernova rates with redshift is computed from estimates of the global history of star formation compiled from multi-wavelength observations of faint galaxies. We assume that redshift dependence of Type IIL frequencies to be the same as that of Type II SNe taken as a whole. Here the absolute value of the frequency is not important, since we are interested in comparing the number of SNe that would be detectable in the presence of a lens to that if the lens were not there. We take each of the background galaxies in our list, assuming their intrinsic sizes to be 10 Kpc, and map them back to the source plane, by means of the lens model and their measured/assumed redshift. The integrals (6) is performed in the image plane by summing over a fine grid, since the presence of the two $`\mathrm{\Theta }`$ functions in the integrand makes it difficult to evaluate them using Gaussian quadrature. We present the curves of $`\mathrm{\Phi }_L(z)`$ for the Abell 2218 HST field for both SN Ia and SN II in Figures 2 and 3 respectively, where we consider SNe in lensed galaxies in the redshift range $`0.175<z_s<2.5`$. For comparison we also give the corresponding values $`\mathrm{\Phi }_U(z)`$ for the unlensed case, to show the dramatic difference the presence of the lens makes. For example, for a limiting magnitude of $`B=25`$, 20% of all SNe beyond $`z>1`$ in the field of the cluster A2218 will be detected, none of which would have been detected if the lens were not present. ## 4 Signal to noise enhancement in giant arcs One of the uncertainties in the accurate photometry of a SN comes from the correction for emission from the host galaxy at the site of the SN. The signal-to-noise ratio (SNR) related to this uncertainty for SNe in lensed galaxies which form significantly elongated arcs or arclets will be better than the SNR in the unlensed galaxies. An enhanced SNR will in turn favour the detection of a SN and the measurement of its characteristics. For a seeing of 0.5–1 arcsec, a galaxy at $`z=1`$ would occupy (referring to the area enclosing 90% of the light) typically $``$20 pixels. On the other hand, a supernova in the galaxy being a point object would occupy the number of pixels covered by the seeing disk, i.e., 3–4 pixels. Here we calculate the enhancement factor $`\eta =\mathrm{𝑆𝑁𝑅}_{\mathrm{lensed}}/\mathrm{𝑆𝑁𝑅}_{\mathrm{unlensed}}`$ for the SN in the galaxy image. For the sake of simplicity, we assume that in the unlensed case the SN is not resolved. The total flux $`F_{tot}`$ that we would receive is the sum of flux from the SN and the galaxy, so the SNR for the detection of the unlensed SN is given by $$SNR\frac{F_{tot}F_{gal}}{\sqrt{F_{tot}+F_{gal}}},$$ (9) where $`F_{gal}`$ is measured long after the occurrence of the SN, so that the latter no longer contributes to the flux from the galaxy. In the lensed case the fluxes should be multiplied by an appropriate average magnification factor $`\mu `$. The lensed galaxy is stretched into an arc in one direction, allowing us to define a stretch factor $`s`$, which is the ratio of the angular size of the seeing disk to that of the arc. The amount of galaxy light contaminating the SN flux is $`s`$ times the unlensed value. Hence the SNR enhancement $`\eta `$ in the lensed case is given by $$\eta =\sqrt{\mu }\left(\frac{_{tot}sF_{gal}}{F_{tot}F_{gal}}\right)\frac{\sqrt{F_{tot}+F_{gal}}}{\sqrt{_{tot}+sF_{gal}}}$$ (10) where $`_{tot}=sF_{gal}+F_{SN}`$. The quantity in brackets is unity, since both numerator and the denominator are equal to $`F_{SN}`$. Denoting $`r=F_{SN}/F_{gal}`$ then the above formula simplifies to $$\eta =\sqrt{\mu }\frac{\sqrt{r+2}}{\sqrt{r+2s}}.$$ (11) For a typical case considered above, $`r0.2`$, $`s0.1`$, $`\mu 40`$, which would give $`\eta 15`$. This shows that an order-of-magnitude enhancement in S/N ratio is achievable in lensed SNe that appear in giant arcs. ## 5 Conclusions A considerable fraction of SNe in high redshift galaxies can be magnified by foreground galaxy clusters. This pushes the equivalent observational distance to SNe further by a factor of 2–3, and thus allows measurements of magnitudes and light-curves of high redshift SNe in significantly shorter observational periods. Even a small telescope with a limiting magnitude of $`m_{\mathrm{lim}}=24`$ will detect SNe Ia up to $`z1.4`$ in the field of a lens like A2218, while in the absence of such a lens, the same setup would not be able to detect SN Ia beyond $`z=0.7`$ (Figure 2). SNe of Type II will never be detected by such a setup beyond $`z=0.5`$ for $`m_{\mathrm{lim}}<26`$, while in the presence of a A2218-like lens, they could be detected up to $`z2`$. SNe occurring in giant arcs will have an additional gain of signal-to-noise of $``$15, making it easier to observe them. The magnification for lensed SNe does not change dramatically for sources between $`z=12`$, while the surface brightness of the host galaxy decreases as $`(1+z)^4`$, which in turn further improves the signal-to-noise ratio. Such SNe will also be multiply imaged, further constraining mass models from time-delay measurements. A number of projects have been searching for SNe at moderate and high redshifts. Neither of the high-$`z`$ SN search projects (Perlmutter et al. (1999); Schmidt et al. (1998)) has yet discovered a SN with a favourable geometry behind a cluster. The low-$`z`$ Abell cluster search (Riess et al. 1998a ) considers clusters out to $`z=0.08`$, again not ideal for these searches. In order to find high-$`z`$ SNe without having to wait for the commissioning of NGST, one needs to continually monitor clusters with known arcs with a favourable geometry (e.g., Abell 2218, Abell 370, CL0024+17) in the spirit of the SN Cosmology searches. Once a SN event is observed, one needs to calculate the magnification (weakly model-dependent) from the location of the SN in the arc. This will yield the intrinsic apparent magnitude of the SN from the light curve, which will enable us to calculate the distance to the supernova independent of its redshift. For a SN in the redshift range $`z=12`$, this method can even yield a reliable value of $`q_0`$, because we would be measuring $`H_0`$ in the relativistic regime. As Fig. 1 shows, in this redshift range, SNe Ia with measured distances can distinguish, for example, between $`\mathrm{\Lambda }=0`$ and $`\mathrm{\Lambda }0`$ models even if the lens model does not allow the estimation of the magnification to better than 0.5 mag. ###### Acknowledgements. We thank J.-P. Kneib for sending us the data for his observation and model of Abell 2218 in electronic form, and Rajaram Nityananda for useful conversations. YS acknowledges the hospitality of IUCAA, where this work began, and the hospitality of Astronomisches Institut, Ruhr Universität Bochum, where he was at the final stages of the work. He received partial financial support from the German Science Foundation (DFG) within Sonderforschungsbereich, a travel grant from the IAU, and a NATO Guest Fellowships grant (Ann. #219.29) from the Italian Consiglio Nazionale delle Ricerche. TDS thanks the University Grants Commission, India for providing the financial support (#2-5/93(II)-E.U. II) without which this work would not have been possible. This research has made use of NASA’s Astrophysics Data System Abstract Service.
warning/0002/cond-mat0002060.html	ar5iv	text	# 1 Introduction ## 1 Introduction The one-component log-gas, consisting of $`N`$ unit charges on a circle of circumference length $`L`$ interacting via the two-dimensional Coulomb potential $`\mathrm{\Phi }(\stackrel{}{r},\stackrel{}{r}^{})=\mathrm{log}\left\|\stackrel{}{r}\stackrel{}{r}^{}\right\|`$, is specified by the Boltzmann factor $$A_{N,\beta }\underset{1j<kN}{}\left\|e^{2\pi ix_k/L}e^{2\pi ix_j/L}\right\|^\beta ,0x_jL.$$ (1.1) The constant $`A_{N,\beta }`$, which plays no role in the calculation of distribution functions, results from scaling the radius of the circle out of the logarithmic potential, and also includes the particle-background and background-background interactions (a uniform neutralizing background is imposed for thermodynamic stability). The thermodynamic limit $`N,L\mathrm{}`$,$`N/L=\rho `$ (fixed) is taken, which gives an infinite system on a straight line with particle density $`\rho `$. This system was first studied because of its relation to the theory of random matrices . The thermodynamic functions were obtained. The pressure $`P`$ has the simple form $$\beta P=\left[1\left(\beta /2\right)\right]\rho $$ (1.2) at any inverse temperature $`\beta `$. However, exact (simple) forms for the correlation functions were obtained by the pioneers only for the special temperatures corresponding to $`\beta =1,2,4`$ (See Section 5). More recently, exact expressions for the two-body density were derived for arbitrary even integer $`\beta `$ and then for arbitrary rational $`\beta `$ . Unfortunately, these latter exact expressions are complicated multivariable integral representations which cannot be easily used as such for actual computations. The purpose of the present paper is to obtain explicit small $`k`$ expansions for the structure function (the Fourier transform of the two-body density). The log-gas is an example of a system interacting via the $`d`$-dimensional Coulomb system (here $`d=2`$) but confined to a domain of dimension $`d1`$. It therefore exhibits universal features — that is features independent of microscopic details such as any short range potential between charges or the number of charge species — characteristic of Coulomb systems in this setting . One universal feature is the existence of an algebraic tail in the leading non-oscillatory term of the large-distance asymptotic expansion of the charge-charge correlation function. For general charged systems in their conductive phase, interacting via the two-dimensional Coulomb potential in a one-dimensional domain, this is predicted to have the form $$\frac{1}{\beta \left(\pi r\right)^2},$$ (1.3) where $`r`$ is the distance. For the one-component log-gas, (1.3) can be verified for all $`\beta `$ rational . The verification is possible because the charge-charge correlation function (which for a one-component system is the same as the density-density correlation) is known explicitly for $`\beta `$ rational (see (2.3) below). In this work we further analyze the properties of the structure factor $`S(k;\beta )`$ (Fourier transform of the charge-charge correlation) for the one-component log-gas. In particular we are interested in the $`\beta `$ dependence of the coefficients in the small $`k`$ expansion of $`S(k;\beta )`$. The large distance behaviour (1.3) is equivalent to the small $`k`$ behaviour $$S(k;\beta )\frac{\left\|k\right\|}{\pi \beta }.$$ (1.4) Furthermore, by making use of the equivalence of the charge-charge and density-density correlation in the one-component log-gas, together with the exact equation of state the second order term in (1.4) has been predicted for general $`\beta `$ , giving $$S(k;\beta )\frac{\left\|k\right\|}{\pi \beta }+\frac{\left(\beta /21\right)k^2}{\left(\pi \beta \right)^2\rho }+O\left(\left\|k\right\|^3\right).$$ (1.5) Let $$f(k;\beta ):=\frac{\pi \beta }{\left\|k\right\|}S(k;\beta ),0<k<\mathrm{min}(2\pi \rho ,\pi \beta \rho )$$ (1.6) and define $`f`$ for $`k<0`$ by analytic continuation (we will see below that $`f(k;\beta )`$ is analytic for $`0\left\|k\right\|<\mathrm{min}(2\pi \rho ,\pi \beta \rho )`$). In Section 2 we use the exact result (2.3) below to derive the functional equation $$f(k;\beta )=f(\frac{2k}{\beta };\frac{4}{\beta })$$ (1.7) The simplest structure consistent with (1.7) is $$\frac{\pi \beta }{\left\|k\right\|}S(k;\beta )=1+\underset{j=1}{\overset{\mathrm{}}{}}p_j\left(\beta /2\right)\left(\frac{\left\|k\right\|}{\pi \beta \rho }\right)^j,\left\|k\right\|<\mathrm{min}(2\pi \rho ,\pi \beta \rho )$$ (1.8) where $`p_j\left(x\right)`$ is a polynomial of degree $`j`$ which satisfies the functional relation $$p_j\left(1/x\right)=\left(1\right)^jx^jp_j\left(x\right).$$ (1.9) Equivalently, (1.9) can be stated as requiring $`p_j\left(x\right)`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{j}{}}}a_{j,l}x^l,a_{j,l}=a_{j,jl}\left(j\mathrm{even}\right)`$ (1.10) $`p_j\left(x\right)`$ $`=`$ $`\left(x1\right){\displaystyle \underset{l=0}{\overset{j1}{}}}\stackrel{~}{a}_{j,l}x^l,\stackrel{~}{a}_{j,l}=\stackrel{~}{a}_{j,j1l}\left(j\mathrm{odd}\right).`$ (1.11) Inspection of (1.5) shows that the conjectured structure (1.8) is correct at order $`\left\|k\right\|`$ and furthermore gives $$p_1\left(x\right)=\left(x1\right),$$ (1.12) and thus $`\stackrel{~}{a}_{1,0}=1`$ in (1.11). In Section 3 we use (2.3) to verify that the structure (1.8) is correct at order $`k^2`$ and we compute $`p_2\left(x\right)`$ explicitly. In Section 4 we use an exact an exact evaluation of the two-particle distribution function for $`\beta `$ even to rederive the result of Section 3, and we also use this formula to verify the structure (1.8) at order $`k^4`$ and to compute $`p_4\left(x\right)`$ explicitly. Assuming the validity of (1.8) we see that $`p_j\left(x\right)`$ can be computed from knowledge of the coefficient of $`\left\|k\right\|^j`$ in $`S(k;\beta )`$, or the coefficient of $`\left\|k\right\|^j`$ in $`^pS(k;\beta )/\beta ^p`$ ($`pj`$), for an appropriate number of distinct values of $`\beta `$. Because the functional relation (1.7) has via (1.10) and (1.11) been made a feature of (1.8) the values of $`\beta `$ cannot be related by $`\beta 4/\beta `$. In Section 5 the known exact evaluation of $`S(k;\beta )`$ to leading order in $`\beta `$ is reviewed, as are the exact evaluations of $`S(k;2)`$ and $`S(k;4)`$. Also noted are the exact evaluations of $`S(k,1)`$ and $`S(k;\beta )`$ to leading order in $`1/\beta `$, which according to (1.7) are related to $`S(k;4)`$ and $`S(k;\beta )`$ to leading order in $`\beta `$ respectively. All of these exact evaluations are in terms of elementary functions, and so can be expanded to all orders in $`k`$. We then present the exact evaluation of $`S(k;\beta )/\beta `$ to leading order in $`\beta `$, as well as the exact evaluation of $`S(k;\beta )/\beta `$ evaluated at $`\beta =2`$ and $`\beta =4`$. The details of the latter two calculations are given in separate appendices. Again the final expressions can be expanded to high order in $`\left\|k\right\|`$. Using this data all polynomials in the expansion (1.8) up to and including the term with $`j=9`$ can be computed. This expansion is written out explicitly in the final section and some special features of the polynomials therein, relating to the sign of the coefficients and the zeros, are noted. A physical interpretation of the functional equation, based on an analogy with a quantum many body system, which identifies an equivalence between quasi-hole and quasi-particle states contributing to $`S(k;\beta )`$ for $`\left\|k\right\|`$ small enough is given. We end with some remarks on the possible occurence of a functional equation analogous to (1.7) in the two-dimensional one-component plasma. ## 2 The functional equation The Boltzmann factor (1.1) also has the physical interpretation as the absolute value squared of the exact ground state wave function, $`\|0`$ say, for the Calogero-Sutherland quantum many body Hamiltonian $$H=\underset{j=1}{\overset{N}{}}\frac{^2}{x_j^2}+\beta \left(\beta /21\right)\left(\frac{\pi }{L}\right)^2\underset{1j<kN}{}\frac{1}{\mathrm{sin}^2\pi \left(x_jx_k\right)/L}.$$ (2.1) This Hamiltonian describes quantum particles on a circle of circumference length $`L`$ interacting via the inverse square of the distance between the particles. In the thermodynamic limit $`N,L\mathrm{}`$, $`N/L=\rho `$ (fixed) the $`N`$ particle system becomes an infinite system on a line with particle density $`\rho `$. The ground state dynamical density-density correlation function $$\rho ^{\mathrm{dyn}.}(0,x;t):=0\left\|n\left(0\right)e^{iHt}n\left(x\right)e^{iHt}\right\|0,n\left(y\right):=\underset{j=1}{\overset{N}{}}\delta \left(yx_j\right)$$ (2.2) in the infinite system has been calculated exactly for all rational $`\beta `$ . The fact that $`\left(\|0\right)^2`$ is proportional to (1.1) tells us that at $`t=0`$ (2.2) is equal to $$\rho _{\left(2\right)}^T(0,x)+\rho \delta \left(x\right),$$ where $`\rho _{\left(2\right)}^T`$ is the truncated two-body density, for the log-gas system. Thus the exact evaluation of $$S(k,\beta ):=_{\mathrm{}}^{\mathrm{}}\left(\rho _{\left(2\right)}^T(0,x)+\rho \delta \left(x\right)\right)e^{ikx}𝑑x$$ for the log-gas follows from the exact evaluation of (2.2) for the quantum system. Taking $`\beta `$ to be rational and setting $$\beta /2:=p/q=:\lambda $$ where $`p`$ and $`q`$ are relatively prime integers, the latter exact result gives $$S(k;\beta )=\pi C_{p,q}\left(\lambda \right)\underset{i=1}{\overset{q}{}}_0^{\mathrm{}}𝑑x_i\underset{j=1}{\overset{p}{}}_0^1𝑑y_jQ_{p,q}^2F(q,p,\lambda \|\{x_i,y_j\})\delta \left(kQ_{p,q}\right),$$ (2.3) where $`C_{p,q}\left(\lambda \right)`$ $`:=`$ $`{\displaystyle \frac{\lambda ^{2p\left(q1\right)}\mathrm{\Gamma }^2\left(p\right)}{2\pi ^2p!q!}}{\displaystyle \frac{\mathrm{\Gamma }^q\left(\lambda \right)\mathrm{\Gamma }^p\left(1/\lambda \right)}{_{i=1}^q\mathrm{\Gamma }^2\left(p\lambda \left(i1\right)\right)_{j=1}^p\mathrm{\Gamma }^2\left(1\left(j1\right)/\lambda \right)}}`$ $`Q_{p,q}`$ $`:=`$ $`2\pi \rho \left({\displaystyle \underset{i=1}{\overset{q}{}}}x_i+{\displaystyle \underset{j=1}{\overset{p}{}}}y_j\right)`$ $`F(q,p,\lambda \|\{x_i,y_j\})`$ $`:=`$ $`{\displaystyle \frac{\underset{i<i^{}}{}\left\|x_ix_i^{}\right\|^{2\lambda }\underset{j<j^{}}{}\left\|y_jy_j^{}\right\|^{2/\lambda }}{_{i=1}^q_{j=1}^p\left(x_i+\lambda y_j\right)^2}}`$ (2.4) $`\times {\displaystyle \frac{1}{_{i=1}^q\left(x_i\left(x_i+\lambda \right)\right)^{1\lambda }_{j=1}^p\left(\lambda y_j\left(1y_j\right)\right)^{11/\lambda }}}`$ In the domain of integration of (2.3) the integration variables are all positive and because of the delta function are restricted to the hyperplane $$\underset{i=1}{\overset{q}{}}x_i+\underset{j=1}{\overset{p}{}}y_j=\frac{\left\|k\right\|}{2\pi \rho }.$$ We see immediately from these constraints that the restriction $`y_j<1`$ in the domain of integration is redundant for $$\left\|k\right\|<2\pi \rho .$$ (2.5) Thus assuming (2.5) we can extend the integration over $`y_j`$ to the region $`(0,\mathrm{})`$. Doing this and changing variables $`x_i\left\|k\right\|x_i`$ and $`y_j\left\|k\right\|y_j`$ we see that for $`\left\|k\right\|`$ in the region (2.5) $$S(k;\beta )=\pi \left\|k\right\|C_{p,q}\left(\lambda \right)\underset{i=1}{\overset{q}{}}_0^{\mathrm{}}𝑑x_i\underset{j=1}{\overset{p}{}}_0^{\mathrm{}}𝑑y_jQ_{p,q}^2\widehat{F}(q,p,\lambda \|\{x_i,y_j\};k)\delta \left(1Q_{p,q}\right),$$ (2.6) where $`\widehat{F}(q,p,\lambda \|\{x_i,y_j\};k)`$ $`=`$ $`{\displaystyle \frac{1}{_{i=1}^q\left(x_i\left(1+kx_i/\lambda \right)\right)^{1\lambda }_{j=1}^p\left(y_j\left(1ky_j\right)\right)^{11/\lambda }}}`$ (2.7) $`\times {\displaystyle \frac{\underset{i<i^{}}{}\left\|x_ix_i^{}\right\|^{2\lambda }\underset{j<j^{}}{}\left\|y_jy_j^{}\right\|^{2/\lambda }}{_{i=1}^q_{j=1}^p\left(x_i+\lambda y_j\right)^2}}`$ Notice that (2.7) is such that the integral in (2.6) is analytic for $$\left\|k\right\|<\mathrm{min}(2\pi \rho ,\pi \rho \beta ).$$ (2.8) Thus according to the definition (1.6) we read off that $$f(k;\beta )=2\pi ^2\lambda C_{p,q}\left(\lambda \right)\underset{i=1}{\overset{q}{}}_0^{\mathrm{}}𝑑x_i\underset{j=1}{\overset{p}{}}_0^{\mathrm{}}𝑑y_jQ_{p,q}^2\widehat{F}(q,p,\lambda \|\{x_i,y_j\};k)\delta \left(1Q_{p,q}\right).$$ (2.9) The functional equation (1.7) is a simple consequence of this exact formula. Thus we see that the integral in (2.9) is unchanged by the mapping $`\lambda 1/\lambda `$ (and thus $`pq`$) followed by $`kk/\lambda `$. The precise functional equation (1.7) follows provided we can show that $$C_{p,q}\left(\lambda \right)=\lambda ^{2pq2}C_{q,p}\left(1/\lambda \right),$$ which indeed readily follows from the definition of $`C_{p,q}\left(\lambda \right)`$ in (2). ## 3 Expanding $`f(k;\beta )`$ in terms of Dotsenko-Fateev type integrals Here we will develop a strategy based on the integral formula (2.9) to expand $`f(k,\beta )`$ at order $`k^2`$. This relies on our ability to compute certain generalizations of a limiting case of the Dotsenko-Fateev integral. This same method has been used in to compute the equivalent of $`f(k,\beta )`$ and its derivative at $`k=0`$. We first expand the integrand in (2.9) as a function of $`k`$. According to (2.7) we have $$\widehat{F}(q,p,\lambda \|\{x_i,y_j\};k)=G(q,p,\lambda \|\{x_i,y_j\})\left(1+\underset{\nu =1}{\overset{\mathrm{}}{}}H_\nu (q,p,\lambda \|\{x_i,y_j\})k^\nu \right)$$ where $`G(q,p,\lambda \|\{x_i,y_j\})`$ $`=`$ $`{\displaystyle \frac{\underset{i<i^{}}{}\left\|x_ix_i^{}\right\|^{2\lambda }\underset{j<j^{}}{}\left\|y_jy_j^{}\right\|^{2/\lambda }}{_{i=1}^q_{j=1}^p\left(x_i+\lambda y_j\right)^2_{i=1}^qx_i^{1\lambda }_{j=1}^py_j^{11/\lambda }}}`$ $`1+{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}H_\nu (q,p,\lambda \|\{x_i,y_j\})k^\nu `$ $`=`$ $`{\displaystyle \frac{1}{_{i=1}^q\left(1+kx_i/\lambda \right)^{1\lambda }_{j=1}^p\left(1ky_j\right)^{11/\lambda }}}.`$ (3.1) The coefficients $`H_\nu `$ are homogeneous polynomials in $`\{x_i,y_j\}`$ of degree $`\nu `$. Let us now introduce the notation $$I_{p,q,\lambda }\left[h\left(\{x_i,y_j\}\right)\right]:=\underset{i=1}{\overset{q}{}}_0^{\mathrm{}}𝑑x_i\underset{j=1}{\overset{p}{}}_0^{\mathrm{}}𝑑y_jQ_{p,q}^2G(q,p,\lambda \|\{x_i,y_j\})\delta \left(1Q_{p,q}\right)h\left(\{x_i,y_j\}\right).$$ (3.2) Because of the presence of the delta function the value of $`I_{p,q,\lambda }`$ is unchanged if $`Q_{p,q}^2`$ is replaced by $`Q_{p,q}^n`$ for any $`n`$. Doing this and also introducing the usual integral representation of the delta function, we see by a change of variables as detailed in that for $`h`$ homogeneous of degree $`\nu `$ $$I_{p,q,\lambda }\left[h\left(\{x_i,y_j\}\right)\right]=\frac{J_{p,q,\lambda ,n}\left[h\left(\{x_i,y_j\}\right)\right]}{\left(\nu +n1\right)!}=\frac{J_{p,q,\lambda }\left[h\left(\{x_i,y_j\}\right)\right]}{\left(\nu 1\right)!}$$ (3.3) where $$J_{p,q,\lambda ,n}\left[h\left(\{x_i,y_j\}\right)\right]:=\underset{i=1}{\overset{q}{}}_0^{\mathrm{}}𝑑x_i\underset{j=1}{\overset{p}{}}_0^{\mathrm{}}𝑑y_jQ_{p,q}^nG(q,p,\lambda \|\{x_i,y_j\})e^{Q_{p,q}}h\left(\{x_i,y_j\}\right)$$ and $`J_{p,q,\lambda }:=J_{p,q,\lambda ,0}`$. Recalling (2.9) and (3) we see that in terms of the notation (3.2) $$f(k;\beta )=C_{p,q}\left(\lambda \right)\left(I_{p,q,\lambda }\left[1\right]+\underset{\nu =1}{\overset{\mathrm{}}{}}I_{p,q,\lambda }\left[H_\nu (q,p,\lambda \|\{x_i,y_j\})\right]k^\nu \right).$$ (3.4) The definition of $`H_\nu `$ in (3) shows $$H_2(q,p,\lambda \|\{x_i,y_j\})=\frac{\left(\lambda 1\right)^2}{2\lambda ^2}\frac{Q_{p,q}^2}{\left(2\pi \rho \right)^2}\frac{\lambda 1}{2\lambda ^2}\left(\underset{i=1}{\overset{q}{}}x_i^2\lambda \underset{j=1}{\overset{p}{}}y_j^2\right),$$ (3.5) so to compute $`f(k,\beta )`$ at order $`k^2`$ our task is to evaluate $$I_{p,q,\lambda }\left[Q_{p,q}^2\right]\mathrm{and}I_{p,q,\lambda }\left[\underset{i=1}{\overset{q}{}}x_i^2\lambda \underset{j=1}{\overset{p}{}}y_j^2\right].$$ (3.6) Now because of the delta function in (3.2) $$I_{p,q,\lambda }\left[Q_{p,q}^2\right]=I_{p,q,\lambda }\left[1\right],$$ (3.7) and we know from that $$C_{p,q}\left(\lambda \right)I_{p,q,\lambda }\left[1\right]=1.$$ (3.8) Thus our remaining task is to compute the second expression in (3.6) or equivalently, using (3.3), to compute $$J_{p,q,\lambda }\left[\underset{i=1}{\overset{q}{}}x_i^2\lambda \underset{j=1}{\overset{p}{}}y_j^2\right]=qJ_{p,q,\lambda }\left[x_i^2\right]\lambda pJ_{p,q,\lambda }\left[y_j^2\right]$$ (3.9) where the second equality, valid for any $`1iq`$ and $`1jp`$, follows from the symmetry of the integrand. For this purpose we first note formulas for $`J_{p,q,\lambda }\left[h\right]`$ in the cases $`h=x_i^2`$ and $`h=y_j^2`$. The formulas are $`J_{p,q,\lambda }\left[x_i^2\right]`$ $`=`$ $`{\displaystyle \frac{\left(2p\lambda +1\right)}{2\pi \rho }}J_{p,q,\lambda }\left[x_i\right]{\displaystyle \frac{p}{\pi \rho }}J_{p,q,\lambda }\left[{\displaystyle \frac{x_i^2}{x_i+\lambda y_j}}\right],`$ (3.10) $`J_{p,q,\lambda }\left[y_j^2\right]`$ $`=`$ $`{\displaystyle \frac{\left(2q1/\lambda +1\right)}{2\pi \rho }}J_{p,q,\lambda }\left[y_j\right]{\displaystyle \frac{\lambda q}{\pi \rho }}J_{p,q,\lambda }\left[{\displaystyle \frac{x_j^2}{x_i+\lambda y_j}}\right].`$ (3.11) The derivation of (3.10) and (3.11) uses a technique based on the fundamental theorem of calculus. It was first used by Aomoto in the context of the Selberg integral, and has been adapted in to the case of the Dotsenko-Fateev integral. Let us give the details of the derivation of (3.10) (the derivation of (3.11) is similar). From the definition (3) we see that $$\frac{}{x_i}G(q,p,\lambda \|\{x_i,y_j\})=\left(\frac{\lambda 1}{x_i}2\underset{j=1}{\overset{p}{}}\frac{1}{x_i+\lambda y_j}+2\lambda \underset{\genfrac{}{}{0pt}{}{i^{}=1}{i^{}i}}{\overset{q}{}}\frac{1}{x_ix_i^{}}\right)G(q,p,\lambda \|\{x_i,y_j\}).$$ Thus $`0`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{q}{}}}{\displaystyle _0^{\mathrm{}}}𝑑x_i{\displaystyle \underset{j=1}{\overset{p}{}}}{\displaystyle _0^{\mathrm{}}}𝑑y_j{\displaystyle \frac{}{x_i}}\left(x_i^2G(q,p,\lambda \|\{x_i,y_j\})e^{Q_{p,q}}\right)`$ $`=`$ $`\left(\lambda +1\right)J_{p,q,\lambda }\left[x_i\right]2{\displaystyle \underset{j=1}{\overset{p}{}}}J_{p,q,\lambda }\left[{\displaystyle \frac{x_i^2}{x_i+\lambda y_j}}\right]+2\lambda {\displaystyle \underset{\genfrac{}{}{0pt}{}{i^{}=1}{i^{}i}}{\overset{q}{}}}J_{p,q,\lambda }\left[{\displaystyle \frac{x_i^2}{x_ix_i^{}}}\right]2\pi \rho J_{p,q,\lambda }\left[x_i^2\right]`$ $`=`$ $`\left(\lambda +1\right)J_{p,q,\lambda }\left[x_i\right]2pJ_{p,q,\lambda }\left[{\displaystyle \frac{x_i^2}{x_i+\lambda y_j}}\right]+2\lambda \left(q1\right)J_{p,q,\lambda }\left[{\displaystyle \frac{x_i^2}{x_ix_i^{}}}\right]2\pi \rho J_{p,q,\lambda }\left[x_i^2\right]`$ where the first equality follows from the fundamental theorem of calculus, while the final equality, valid for any $`j=1,\mathrm{},p`$ and any $`i^{}=1,\mathrm{},q`$, ($`i^{}i`$) follows by the symmetry of the integrand with respect to $`\left\{x_i\right\}`$ and $`\left\{y_j\right\}`$. The symmetry of the integrand with respect to $`\left\{x_i\right\}`$ also gives $$J_{p,q,\lambda }\left[\frac{x_i^2}{x_ix_i^{}}\right]=J_{p,q,\lambda }\left[\frac{x_i^{}^2}{x_i^{}x_i}\right]$$ so we have $$J_{p,q,\lambda }\left[\frac{x_i^2}{x_ix_i^{}}\right]=\frac{1}{2}\left(J_{p,q,\lambda }\left[\frac{x_i^2}{x_ix_i^{}}\right]+J_{p,q,\lambda }\left[\frac{x_i^{}^2}{x_i^{}x_i}\right]\right)=J_{p,q,\lambda }\left[x_i\right].$$ Substituting in (3) implies (3.10). From (3.10) and (3.11) we see that $`qJ_{p,q,\lambda }\left[x_i^2\right]\lambda pJ_{p,q,\lambda }\left[y_j^2\right]`$ $`={\displaystyle \frac{q\left(2p\lambda +1\right)}{2\pi \rho }}J_{p,q,\lambda }\left[x_i\right]{\displaystyle \frac{\lambda p\left(2q1/\lambda +1\right)}{2\pi \rho }}J_{p,q,\lambda }\left[y_j\right]{\displaystyle \frac{pq}{\pi \rho }}J_{p,q,\lambda }\left[{\displaystyle \frac{x_i^2\lambda ^2y_j^2}{x_i+\lambda y_j}}\right]`$ $`={\displaystyle \frac{q\left(\lambda +1\right)}{2\pi \rho }}J_{p,q,\lambda }\left[x_i\right]{\displaystyle \frac{\lambda p\left(1/\lambda +1\right)}{2\pi \rho }}J_{p,q,\lambda }\left[y_j\right]={\displaystyle \frac{\lambda 1}{\left(2\pi \rho \right)^2}}J_{p,q,\lambda }\left[Q_{p,q}\right]={\displaystyle \frac{\lambda 1}{\left(2\pi \rho \right)^2}}J_{p,q,\lambda }\left[1\right]`$ Recalling (3.5), the results (3.3), (3.7), (3.9) and (3) give that $$I_{p,q,\lambda }[H_2(q,p,\lambda \|\left\{x_i\right\},\left\{y_j\right\}]=\frac{1}{\left(2\pi \rho \right)^2}\frac{\left(\lambda 1\right)^2}{\lambda ^2}J_{p,q,\lambda }\left[1\right].$$ Use of (3.8) then gives that the term proportional to $`k^2`$ in (3.4) is equal to $$\left(\lambda 1\right)^2\left(\frac{k}{2\pi \lambda \rho }\right)^2=\left(\beta /21\right)^2\left(\frac{k}{\pi \beta \rho }\right)^2.$$ (3.14) It follows from this that the structure (1.8) is valid at order $`k^2`$ with $$p_2\left(x\right)=\left(x1\right)^2.$$ (3.15) ## 4 Large-$`x`$ expansion of $`\rho _{(2)}^T(0,x)`$ We have already remarked that the large-$`x`$ expansion (1.3) of the charge-charge correlation, or what is the same thing for the one-component log-gas, the large-$`x`$ expansion of $`\rho _{\left(2\right)}^T(0,x)`$, is equivalent to the small-$`k`$ behaviour (1.4) of $`S(k;\beta )`$. More generally the expansion $$\rho _{\left(2\right)}^T(0,x)\underset{x\mathrm{}}{}\underset{n=1}{\overset{\mathrm{}}{}}\frac{c_n}{x^{2n}}$$ (4.1) is equivalent to the expansion $$S(k;\beta )\underset{k0}{}\pi \underset{n=1}{\overset{\mathrm{}}{}}\frac{\left(1\right)^nc_n}{\left(2n1\right)!}\left\|k\right\|^{2n1}$$ (4.2) where the expansion (4.2) contains the terms singular in $`k`$ (i.e. of odd order in $`\left\|k\right\|`$) only. This follows using the Fourier transform $$_{\mathrm{}}^{\mathrm{}}\frac{e^{ikx}}{x^{2n}}𝑑x=\pi \frac{\left(1\right)^n\left\|k\right\|^{2n1}}{\left(2n1\right)!}$$ from the theory of generalized functions (see e.g. ). From the equivalence between (4.1) and (4.2) we see the fact, following from (3.14), that the term proportional to $`\left\|k\right\|^3`$ in the small $`k`$ expansion of $`S(k;\beta )`$ is equal to $$\rho \left(\beta /21\right)^2\left(\frac{\left\|k\right\|}{\pi \beta \rho }\right)^3$$ is equivalent to the statement that the term proportional to $`1/x^4`$ in the large $`x`$ expansion of $`\rho _{\left(2\right)}^T(0,x)`$ is equal to $$\rho ^26\beta \left(\beta /21\right)^2\left(\frac{1}{\pi \beta \rho x}\right)^4.$$ (4.3) In this section we will derive (4.3) directly. We will also calculate the $`O\left(1/x^6\right)`$ term and so explicitly determine the $`O\left(\left\|k\right\|^5\right)`$ term in (4.2). The starting point for our calculation is an exact $`\beta `$-dimensional integral formula for the two-particle distribution $`\rho _{\left(2\right)}(0,x)`$ valid for $`\beta `$ even. With $$S_n(a,b,c):=\underset{j=0}{\overset{n1}{}}\frac{\mathrm{\Gamma }\left(a+1+jc\right)\mathrm{\Gamma }\left(b+1+jc\right)\mathrm{\Gamma }\left(1+\left(j+1\right)c\right)}{\mathrm{\Gamma }\left(a+b+2+\left(N+j1\right)c\right)\mathrm{\Gamma }\left(1+c\right)}$$ the formula gives that in the thermodynamic limit $`\rho _{\left(2\right)}(0,x)`$ $`=`$ $`\rho ^2\left(\beta /2\right)^\beta {\displaystyle \frac{\left(\left(\beta /2\right)!\right)^3}{\beta !\left(3\beta /2\right)!}}{\displaystyle \frac{e^{\pi i\beta \rho x}\left(2\pi \rho x\right)^\beta }{S_\beta (12/\beta ,12/\beta ,2/\beta )}}`$ (4.4) $`\times {\displaystyle _{[0,1]}}du_1\mathrm{}du_\beta {\displaystyle \underset{j=1}{\overset{\beta }{}}}e^{2\pi i\rho xu_j}u_j^{1+2/\beta }(1u_j)^{1+2/\beta }{\displaystyle \underset{j<k}{}}\|u_ku_j\|^{4/\beta }.`$ In a previous analysis it has been shown that the non-oscillatory large-$`x`$ behaviour is determined by the integrand in the vicinity of the endpoints 0 and 1, with the requirement that $`\beta /2`$ of the integration variables are in the vicinity of the endpoint 0, while the remaining $`\beta /2`$ integration variables are in the vicinity of the endpoint 1. Thus we write $`u_{\beta /2+j}=1v_j`$ ($`j=1,\mathrm{},\beta /2`$) (this introduces a combinatorial factor $`\beta `$ choose $`\beta /2`$ to account for the different ways of so partitioning the integration variables) and then expand the integrand (excluding the exponential factors which involve $`x`$) in terms of the “small” variables $`u_j,v_j`$ ($`j=1,\mathrm{},\beta /2`$). In particular we must expand $$\underset{j=1}{\overset{\beta /2}{}}\left(1u_j\right)^{1+2/\beta }\left(1v_j\right)^{1+2/\beta }\underset{l,l^{}=1}{\overset{\beta /2}{}}\left(1u_lv_l^{}\right)^{4/\beta }.$$ (4.5) The function (4.5) is a symmetric function of the variables $`\left\{u_j\right\}`$ and $`\left\{v_j\right\}`$ separately. Let $`\left\{q_\kappa \right\}_\kappa `$ be a polynomial basis for symmetric functions with $`\kappa `$ denoting a partition (ordered set of non-negative integers) of no more than $`\beta /2`$ parts, and suppose furthermore that $`q_\kappa `$ is homogeneous of order $`\left\|\kappa \right\|:=\kappa _1+\mathrm{}+\kappa _{\beta /2}`$. Then we can write $`{\displaystyle \underset{j=1}{\overset{\beta /2}{}}}\left(1u_j\right)^{1+2/\beta }\left(1v_j\right)^{1+2/\beta }{\displaystyle \underset{l,l^{}=1}{\overset{\beta /2}{}}}\left(1u_lv_l^{}\right)^{4/\beta }`$ (4.6) $`={\displaystyle \underset{\kappa ,\mu }{}}w_{\kappa ,\mu }q_\kappa (u_1,\mathrm{},u_{\beta /2})q_\mu (v_1,\mathrm{},v_{\beta /2}).`$ Substituting (4.6) in (4.4), then following the procedure of , which involves extending the range of integration to $`u_j(0,\mathrm{})`$, $`v_j(0,\mathrm{})`$ and changing variables $`u_j2\pi i\rho xu_j`$, $`v_j2\pi i\rho xv_j`$ making use in the process of the fact that $`q_\kappa `$ is homogeneous of degree $`\left\|\kappa \right\|`$, we obtain the non-oscillatory terms in the large-$`x`$ asymptotic expansion of $`\rho _{\left(2\right)}(0,x)`$. This reads $`\rho _{\left(2\right)}(0,x)`$ $``$ $`\rho ^2\left({\displaystyle \genfrac{}{}{0pt}{}{\beta }{\beta /2}}\right)\left(\beta /2\right)^\beta {\displaystyle \frac{\left(\left(\beta /2\right)!\right)^3}{\beta !\left(3\beta /2\right)!}}{\displaystyle \frac{1}{S_\beta (12/\beta ,12/\beta ,2/\beta )}}`$ $`\times `$ $`{\displaystyle \underset{\kappa ,\mu }{}}w_{\kappa ,\mu }{\displaystyle \frac{K_{\beta ,\kappa }K_{\beta ,\mu }}{i^{\left\|\lambda \right\|\left\|\mu \right\|}\left(2\pi \rho x\right)^{\left\|\kappa \right\|+\left\|\mu \right\|}}}`$ (4.7) where $$K_{\beta ,\kappa }:=_{[0,\mathrm{})^{\beta /2}}𝑑u_1\mathrm{}𝑑u_{\beta /2}\underset{l=1}{\overset{\beta /2}{}}u_l^{1+2/\beta }e^{u_l}\underset{j<k}{}\left\|u_ku_j\right\|^{4/\beta }q_\kappa (u_1,\mathrm{},u_{\beta /2}).$$ (4.8) The symmetry $`w_{\kappa ,\mu }=w_{\mu ,\kappa }`$ evident from (4.6) implies terms in (4) with $`\left\|\kappa \right\|+\left\|\mu \right\|`$ odd cancel. Therefore the sum in (4) can be restricted to partitions such that $`\left\|\kappa \right\|+\left\|\mu \right\|`$ is even, which means the asymptotic expansion only contains inverse even powers of $`x`$. To proceed further we must be able to compute the expansion coefficients $`w_{\kappa ,\mu }`$ as well as the integrals $`K_{\beta ,\kappa }`$. For the former task it is convenient to choose $`q_\kappa `$ equal to the monomial symmetric polynomial $`m_\kappa `$, which is defined as the symmetrization of the monomial $`x_1^{\kappa _1}\mathrm{}x_{\beta /2}^{\kappa _{\beta /2}}`$ normalized so that the coefficient of $`x_1^{\kappa _1}\mathrm{}x_{\beta /2}^{\kappa _{\beta /2}}`$ is unity. First, we have the well known expansion $$\underset{j=1}{\overset{n}{}}\left(1u_j\right)^a=\underset{\mathrm{}\left(\kappa \right)n}{}a_\kappa m_\kappa (u_1,\mathrm{},u_n)$$ (4.9) where $$a_\kappa =\underset{p=1}{\overset{\mathrm{}\left(\kappa \right)}{}}a_{\kappa _p},a_k:=\frac{\left(a\right)_k}{k!}$$ (4.10) with $`\mathrm{}\left(\kappa \right)`$ denoting the length of $`\kappa `$ (i.e. number of non-zero parts). We can therefore immediately expand the first product in (4.6) in terms of monomial symmetric polynomials. Consider next the expansion of the double product in (4.6). Making use of the formulas $`\left(1x\right)^a`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\left(a\right)_n}{n!}}x^n`$ (4.11) $`{\displaystyle \underset{j=1}{\overset{n}{}}}\left({\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}a_kt_j^k\right)`$ $`=`$ $`{\displaystyle \underset{\mathrm{}\left(\kappa \right)n}{}}a_0^{N\mathrm{}\left(\kappa \right)}a_\kappa m_\kappa \left(\left\{t_j\right\}\right)`$ (4.12) where $`a_\kappa `$ is specified by the first equality in (4.10), we see that $$\underset{j=1}{\overset{\beta /2}{}}\left(1u_jv\right)^{4/\beta }=\underset{\mathrm{}\left(\kappa \right)\beta /2}{}\left(1v\right)^{2\left\|\kappa \right\|}c_\kappa m_\kappa (u_1,\mathrm{},u_{\beta /2})$$ (4.13) where $$c_\kappa =\underset{p=1}{\overset{\mathrm{}\left(\kappa \right)}{}}c_{\kappa _p},c_k:=\frac{\left(4/\beta \right)_k}{k!}.$$ Expanding the factor $`\left(1v\right)^{2\left\|\kappa \right\|}`$ we can rewrite (4.13) as $$\underset{j=1}{\overset{\beta /2}{}}\left(1u_jv\right)^{4/\beta }=\underset{n=0}{\overset{\mathrm{}}{}}w_n(u_1,\mathrm{},u_{\beta /2};\beta )v^n$$ for appropriate symmetric functions $`w_n`$. Replacing $`v`$ by $`v_j^{}`$ and forming the product over $`j^{}`$ using (4.12) we obtain $$\underset{j,j^{}=1}{\overset{\beta /2}{}}\left(1u_jv_j^{}\right)^{4/\beta }=\underset{\mathrm{}\left(\kappa \right)\beta /2}{}w_0^{\beta /2\mathrm{}\left(\kappa \right)}w_\kappa m_\kappa (v_1,\mathrm{},v_{\beta /2})$$ where $`w_\kappa :=_{p=1}^{\mathrm{}\left(\kappa \right)}w_{\kappa _p}`$. The final step is to expand $`w_0^{\beta /2\mathrm{}\left(\kappa \right)}w_\kappa `$ in terms of $`\left\{m_\mu \right\}`$ and so obtain the expansion $$\underset{j,j^{}=1}{\overset{\beta /2}{}}\left(1u_jv_j^{}\right)^{4/\beta }=\underset{\mu ,\kappa }{}t_{\mu ,\kappa }m_\mu (u_1,\mathrm{},u_{\beta /2})m_\kappa (v_1,\mathrm{},v_{\beta /2}).$$ (4.14) The practical implementation of this procedure requires the use of computer algebra. We work with arbitrary (positive integer) values of $`\beta /2`$. Furthermore, we only include terms with $`\left\|\mu \right\|+\left\|\kappa \right\|6`$ throughout since according to (4) these terms suffice for the evaluation of the coefficients of $`1/x^{2n}`$, $`n3`$. Having obtained the coefficients $`t_{\mu ,\kappa }`$ in (4.14), we multiply the series (4.14) with the two series of the form (4.9) representing the first two products in (4.5), expressing the answer in the form of (4.6), and so determining the coefficients $`w_{\kappa ,\mu }`$. Again this step requires computer algebra. With $`w_{\kappa ,\mu }`$ in (4.6) determined, it remains to compute the multiple integral (4.8) with $`q_\mu =m_\mu `$. For this task we introduce a further basis of symmetric functions, namely the Jack polynomials $`\left\{P_\kappa ^{\left(\beta /2\right)}(u_1,\mathrm{},u_{\beta /2})\right\}`$. The Jack polynomials $`P_\kappa ^{\left(2/\beta \right)}(z_1,\mathrm{},z_N)`$ with $`z_j:=e^{2\pi ix_j/L}`$, when muliplied by the ground state wave function $`\|0`$, are the eigenfunctions of the Calogero-Sutherland Schrödinger operator (2.1) . Each polynomial is homogeneous of degree $`\left\|\kappa \right\|`$ and has the expansion $$P_\kappa ^{\left(\alpha \right)}(z_1,\mathrm{},z_N)=m_\kappa +\underset{\mu <\kappa }{}a_{\kappa \mu }m_\mu $$ (4.15) where $`<`$ is the dominance partial ordering for partitions: $`\mu <\kappa `$ if $`\left\|\kappa \right\|=\left\|\mu \right\|`$ with $`\kappa \mu `$ and $`_{i=1}^p\mu _i_{i=1}^p\kappa _i`$ for each $`p=1,\mathrm{},N`$. The coefficients $`a_{\kappa \mu }`$ can be calculated by recurrence . The significance of the Jack polynomial basis is that we have the explicit integral evaluation $$\frac{1}{W_{a\alpha N}}\underset{l=1}{\overset{N}{}}_0^{\mathrm{}}𝑑t_lt_l^ae^{t_l}P_\kappa ^{\left(\alpha \right)}(t_1,\mathrm{},t_N)\underset{j<k}{}\left\|t_kt_j\right\|^{2/\alpha }=P_\kappa ^{\left(\alpha \right)}\left(1^N\right)\left[a+\left(N1\right)/\alpha +1\right]_\kappa ^{\left(\alpha \right)},$$ (4.16) which is a limiting case of an integration formula due to Macdonald , Kadell and Kaneko . In (4.16) $`W_{a\alpha N}`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{N}{}}}{\displaystyle _0^{\mathrm{}}}𝑑t_lt_l^ae^{t_l}{\displaystyle \underset{j<k}{}}\left\|t_kt_j\right\|^{2/\alpha }={\displaystyle \underset{j=0}{\overset{N1}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(1+\left(j+1\right)/\alpha \right)\mathrm{\Gamma }\left(a+1+j/\alpha \right)}{\mathrm{\Gamma }\left(1+1/\alpha \right)}},`$ $`\left[u\right]_\kappa ^{\left(\alpha \right)}`$ $`:=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(u\left(j1\right)/\alpha +\kappa _j\right)}{\mathrm{\Gamma }\left(u\left(j1\right)/\alpha \right)}}`$ and $`P_\kappa ^{\left(\alpha \right)}\left(1^N\right)`$ denotes $`P_\kappa ^{\left(\alpha \right)}(x_1,\mathrm{},x_N)`$ evaluated at $`x_1=\mathrm{}=x_N=1`$. To make use of (4.16) we must first express the monomial symmetric polynomials $`m_\kappa `$ in terms of $`\left\{P_\mu ^{2/\beta }\right\}_{\mu \kappa }`$, which can be done using computer algebra from knowledge of the expansion (4.15). Substituting in (4.16) allows the integrals $`K_{\kappa \beta }`$ to be computed. After completing this procedure all terms in (4) for $`\left\|\kappa \right\|+\left\|\mu \right\|6`$ are known explicitly. Performing the sum and simplifying we obtain $$\rho _{\left(2\right)}(0,x)\rho ^2\left(1\frac{1}{\beta \left(\pi \rho x\right)^2}+\frac{3\left(\beta 2\right)^2}{2\beta ^3\left(\pi \rho x\right)^4}\frac{15\left(\beta 2\right)^2\left(\beta ^23\beta +4\right)}{2\beta ^5\left(\pi \rho x\right)^6}+\mathrm{}\right).$$ (4.17) Note that this agrees with the known form (1.3) for the term $`O\left(1/x^2\right)`$, and the form (4.3) for the term $`O\left(1/x^4\right)`$. The term $`O\left(1/x^6\right)`$, due to the equivalence between (4.1) and (4.2), implies the term $`O\left(\left\|k\right\|^5\right)`$ in the small-$`k`$ expansion of $`S(k,\beta )`$ is equal to $$\left(\beta /21\right)^2\left(\left(\beta /2\right)^2\frac{3}{2}\left(\beta /2\right)+1\right)\left(\frac{\left\|k\right\|}{\pi \beta }\right)^5$$ (4.18) This is of the form of the conjecture (1.8) with $$p_4\left(x\right)=\left(x1\right)^2\left(x^2\frac{3}{2}x+1\right).$$ (4.19) ## 5 $`S(k;\beta )`$ for special $`\beta `$ Let us assume the validity of (1.8). The coefficients specifying the polynomials $`p_j\left(x\right)`$ therein can be determined from knowledge of the coefficient of $`\left\|k\right\|^{j+1}`$ in $`S(k;\beta )`$ or $`^pS(k;\beta )/\beta ^p`$ ($`pj`$) at special values of $`\beta `$. Now in the context of random matrix theory $`S(k;\beta )`$ has been evaluated in terms of elementary functions for $`\beta =1,2`$ and 4. The results are $`S(k;1)`$ $`=`$ $`\{\begin{array}{cc}\frac{\left\|k\right\|}{\pi }\frac{\left\|k\right\|}{2\pi }\mathrm{log}\left(1+\frac{\left\|k\right\|}{\pi \rho }\right),\hfill & \left\|k\right\|2\pi \rho \hfill \\ 2\frac{\left\|k\right\|}{2\pi }\mathrm{log}\left(\frac{1+\left\|k\right\|/\pi \rho }{1+\left\|k\right\|/\pi \rho }\right),\hfill & \left\|k\right\|2\pi \rho \hfill \end{array}`$ (5.3) $`S(k;2)`$ $`=`$ $`\{\begin{array}{cc}\frac{\left\|k\right\|}{2\pi },\hfill & \left\|k\right\|2\pi \rho \hfill \\ 1,\hfill & \left\|k\right\|2\pi \rho \hfill \end{array}`$ (5.6) $`S(k;4)`$ $`=`$ $`\{\begin{array}{cc}\frac{\left\|k\right\|}{4\pi }\frac{\left\|k\right\|}{8\pi }\mathrm{log}\left\|1\frac{\left\|k\right\|}{2\pi \rho }\right\|,\hfill & \left\|k\right\|4\pi \rho \hfill \\ 1,\hfill & \left\|k\right\|4\pi \rho \hfill \end{array}`$ (5.9) Recalling the definition (1.6) of $`f(k,\beta )`$ we read off $`f(k;1)`$ $`=`$ $`1{\displaystyle \frac{1}{2}}\mathrm{log}\left(1+{\displaystyle \frac{k}{\pi \rho }}\right)`$ (5.10) $`f(k;2)`$ $`=`$ $`1`$ (5.11) $`f(k;4)`$ $`=`$ $`1{\displaystyle \frac{1}{2}}\mathrm{log}\left(1{\displaystyle \frac{k}{2\pi \rho }}\right)`$ (5.12) The exact evaluation (5.11) implies that for all $`j`$ $`p_j\left(x\right)`$ contains a factor of $`\left(x1\right)`$. In the case of $`j`$ odd this gives no new information since the factor $`\left(x1\right)`$ was already deduced as a consequence of the functional equation (1.9). On the other hand, in the case $`j`$ even this fact together with the functional equation (1.9) implies $$p_j\left(x\right)=\left(x1\right)^2\underset{l=0}{\overset{j2}{}}b_{j,l}x^l,b_{j,l}=b_{j,j2l}\left(j\mathrm{even}\right).$$ (5.13) Consider now the constraints on the coefficients in (5.13) and (1.11) which follow from (5.10) and (5.12). As (5.10) and (5.12) are related by the functional equation (1.7), and this is built into the structures (5.13) and (1.11), only one of these exact evaluations gives distinct information on $`p_j\left(x\right)`$. For definiteness consider (5.10). We see that $$\left[k^j\right]f(k;1)=\frac{1}{2}\frac{\left(1\right)^j}{j\left(\pi \rho \right)^j},j1$$ (5.14) where the notation $`\left[k^j\right]`$ denotes the coefficient of $`k^j`$. Recalling (1.8), (5.13) and (1.11) this implies, for $`j`$ even, $$\frac{1}{j}=\frac{1}{2}\left(\left(1+2^{\left(j2\right)}\right)b_{j,0}+\left(2^1+2^{\left(j3\right)}\right)b_{j,1}+\mathrm{}+\left(2^{j/2+2}+2^{j/2}\right)b_{j,j/22}+2^{j/2+1}b_{j,j/21}\right),$$ (5.15) while for $`j`$ odd $`{\displaystyle \frac{1}{j}}`$ $`=`$ $`((1+2^{\left(j1\right)})\stackrel{~}{a}_{j,0}+(2^1+2^{\left(j2\right)})\stackrel{~}{a}_{j,1}+\mathrm{}+(2^{\left(j1\right)/2+1}+2^{\left(j1\right)/21})\stackrel{~}{a}_{j,\left(j1\right)/21}`$ (5.16) $`+2^{\left(j1\right)/2}\stackrel{~}{a}_{j,\left(j1\right)/2}).`$ In the case $`j=1`$ (5.16) gives $`\stackrel{~}{a}_{j,0}=1`$ which reclaims (1.12), while in the case $`j=2`$ (5.15) gives $`b_{j,0}=1`$ which reclaims (3.15). The exact form of $`S(k;\beta )`$ in the weak coupling scaling limit $`\beta 0`$, $`k0`$, $`k/\beta `$ fixed is also available. Introducing the dimensionless Fourier transforms $$\stackrel{~}{S}(k;\beta ):=\rho _{\mathrm{}}^{\mathrm{}}\left(\rho _{\left(2\right)}^T(0,x)+\rho \delta \left(x\right)\right)e^{i\rho xk}𝑑x,\stackrel{~}{\mathrm{\Phi }}\left(k\right):=\rho _{\mathrm{}}^{\mathrm{}}\mathrm{\Phi }\left(x\right)e^{i\rho xk}𝑑x$$ where $`\mathrm{\Phi }\left(x\right):=\mathrm{log}\left\|x\right\|`$ is the pair potential of the log-gas (thus the integral in the definition of the $`\stackrel{~}{\mathrm{\Phi }}\left(k\right)`$ is to be interpreted as a generalized function) we have $$\stackrel{~}{S}(k;\beta )1\frac{\beta \stackrel{~}{\mathrm{\Phi }}\left(k\right)}{1+\beta \stackrel{~}{\mathrm{\Phi }}\left(k\right)}.$$ (5.17) Since $$\stackrel{~}{\mathrm{\Phi }}\left(k\right)=\frac{\pi }{\left\|k\right\|},$$ (5.18) and noting $`\stackrel{~}{S}(k;\beta )=S(k\rho ;\beta )/\rho `$ we thus have that in the weak coupling scaling limit $$S(k,\beta )=\rho \left(1\frac{1}{1+\left\|k\right\|/\pi \beta \rho }\right).$$ (5.19) Expanding (5.19) in the form (1.8) and recalling (5.13) and (1.11) we deduce $$\stackrel{~}{a}_{j,0}=1\mathrm{and}b_{j,0}=1$$ (5.20) for all $`j`$. Using (5.20) in (5.15) and (5.16) gives that in the case $`j=3`$, $`\stackrel{~}{a}_{j,1}=\frac{11}{6}`$, and in the case $`j=4`$, $`b_{j,1}=\frac{3}{2}`$. The latter result reclaims (4.18) while the former result together with (5.20) gives $$p_3\left(x\right)=\left(x1\right)\left(1\frac{11}{6}x+x^2\right).$$ (5.21) An alternative way to derive (5.20) is to consider the $`\beta \mathrm{}`$ low temperature limit. In this limit the system behaves like an harmonic crystal, for which we have available the analytic formula <sup>2</sup><sup>2</sup>2The denominator of the exponent in (3.10) of contains a spurious factor of $`\pi ^2`$ which is corrected in (5.22) $$\rho _{\left(2\right)}^{\left(\mathrm{har}\right)}(x;0)=\rho ^2\underset{\genfrac{}{}{0pt}{}{p=\mathrm{}}{p0}}{\overset{\mathrm{}}{}}\left(\frac{\beta }{4\pi f\left(p\right)}\right)^{1/2}e^{\beta \left(p\rho x\right)^2/4f\left(p\right)}$$ (5.22) where $$f\left(p\right)=\frac{1}{\pi ^2}_0^{1/2}\frac{1\mathrm{cos}2\pi pt}{tt^2}𝑑t.$$ Taking the Fourier transform gives for $`\left\|k\right\|<2\pi \rho `$ $`S^{\left(\mathrm{har}\right)}(k;\beta )`$ $`=`$ $`\rho {\displaystyle \underset{p=\mathrm{}}{\overset{\mathrm{}}{}}}\left(e^{k^2f\left(p\right)/\beta \rho ^2}1\right)e^{ikp/\rho }`$ (5.23) $`\underset{\beta \mathrm{}}{}`$ $`\rho {\displaystyle \frac{k^2}{\beta \rho ^2}}{\displaystyle \underset{p=\mathrm{}}{\overset{\mathrm{}}{}}}f\left(p\right)e^{ikp/\rho }={\displaystyle \frac{\left\|k\right\|/\pi \beta }{1\left\|k\right\|/2\pi \rho }}.`$ This formula maps to the weak coupling result (5.17) under the action of the functional equation (1.7) and so implies (5.20). ## 6 Perturbation about $`\beta =0`$ The formula (5.19) is just the first term in a systematic weak coupling renormalized Mayer series expansion in $`\beta `$. In the case of the two-dimensional one-component plasma, low order terms of this expansion have recently been analyzed by Kalinay et al. . Results from that study can readily be transcribed to the case of the one-component log-gas. Formally, the renormalized Mayer series expansion is for the dimensionless free energy $`\beta \overline{F}^{\mathrm{ex}}`$ (in our $`\beta \overline{F}^{\mathrm{ex}}`$ is written $`\beta \overline{F}^{\mathrm{ex}}`$), and one computes the direct correlation function via the functional differentiation formula $$c(0,x)=\frac{\delta ^2\left(\beta \overline{F}^{\mathrm{ex}}\right)}{\delta \rho _{\left(1\right)}\left(0\right)\delta \rho _{\left(1\right)}\left(x\right)}.$$ (6.1) The Ornstein-Zernicke relation gives that the dimensionless Fourier transform of the direct correlation function, $`\stackrel{~}{c}(k,\beta )`$ say, is related to the dimensionless structure function $`\stackrel{~}{S}(k;\beta )`$ by $$\stackrel{~}{c}(k;\beta )=1\frac{1}{\stackrel{~}{S}(k;\beta )}$$ (6.2) so expanding $`\stackrel{~}{c}(k,\beta )`$ about $`\beta =0`$ with $`k/\beta `$ fixed is equivalent to expanding $`\stackrel{~}{S}(k;\beta )`$ about $`\beta =0`$ with $`k/\beta `$ fixed. Now, transcribing the results of we read off that the weak coupling diagrammatic expansion of $`c(x_1,x_2)`$ starts as $$c(x_1,x_2)=\beta \mathrm{\Phi }(x_1,x_2)+\frac{1}{2!}\left(K(x_1,x_2)\right)^2+\mathrm{}$$ (6.3) where $$K(x_1,x_2)=\beta \pi _{\mathrm{}}^{\mathrm{}}\frac{e^{ik\left(x_1x_2\right)}}{\left\|k\right\|+\kappa }𝑑k$$ (6.4) with $`\kappa =\beta \pi \rho `$. This implies $$\stackrel{~}{c}(k;\beta )=\frac{\beta \pi }{\left\|k\right\|}+\frac{1}{2}\rho _{\mathrm{}}^{\mathrm{}}\frac{dl}{2\pi }\frac{\beta \pi }{\left\|l\right\|+\kappa }\frac{\beta \pi }{\left\|\rho kl\right\|+\kappa }.$$ (6.5) The integral is straightforward (consider separately the ranges of $`l`$ such that $`l>0`$ ($`l<0`$) and $`\rho kl>0`$ ($`\rho kl<0`$)). In terms of $`k^{}:=\rho k/\kappa =k/\pi \beta `$, $$\stackrel{~}{c}(k;\beta )=\frac{1}{\left\|k^{}\right\|}+\beta \frac{1+\left\|k^{}\right\|}{\left\|k^{}\right\|\left(2+\left\|k^{}\right\|\right)}\mathrm{log}\left(1+\left\|k^{}\right\|\right)+O\left(\beta ^2\right),$$ (6.6) or equivalently using (6.2) $$S(k;\beta )=\rho \frac{\left\|k/\kappa \right\|}{1+\left\|k/\kappa \right\|}+\beta \rho \frac{\left\|k/\kappa \right\|}{\left(1+\left\|k/\kappa \right\|\right)\left(2+\left\|k/\kappa \right\|\right)}\mathrm{log}\left(1+\left\|k/\kappa \right\|\right)+O\left(\beta ^2\right).$$ (6.7) Notice that the leading order term in (6.7) reproduces (5.17). The exact result (6.7) gives the explicit value of the coefficient of $`x`$ in the polynomial $`p_j\left(x\right)`$. Thus recalling (1.11) and (5.13) we have $`{\displaystyle \frac{1}{2}}\left(b_{j,1}2\right)`$ $`=`$ $`\left[x^j\right]{\displaystyle \frac{1}{\left(1+x\right)\left(2+x\right)}}\mathrm{log}\left(1+x\right),\left(j\mathrm{even}\right)`$ $`{\displaystyle \frac{1}{2}}\left(1\stackrel{~}{a}_{j,1}\right)`$ $`=`$ $`\left[x^j\right]{\displaystyle \frac{1}{\left(1+x\right)\left(2+x\right)}}\mathrm{log}\left(1+x\right),\left(j\mathrm{odd}\right).`$ (6.8) Furthermore, a simple calculation gives $$\left[x^j\right]\frac{1}{\left(1+x\right)\left(2+x\right)}\mathrm{log}\left(1+x\right)=\left(1\right)^j\underset{q=1}{\overset{j}{}}\frac{1}{q}\left(12^{qj}\right)$$ (6.9) so we have for example $$\stackrel{~}{a}_{5,1}=\frac{91}{30},b_{6,1}=\frac{31}{15},\stackrel{~}{a}_{7,1}=\frac{1607}{420},b_{8,1}=\frac{263}{84},\stackrel{~}{a}_{9,1}=\frac{791}{180}.$$ (6.10) Substituting $`\stackrel{~}{a}_{5,1}`$ from (6.10) and $`\stackrel{~}{a}_{5,0}`$ from (5.20) in (5.15) shows $`\stackrel{~}{a}_{5,2}=\frac{62}{15}`$. Similarly, the value of $`b_{6,1}`$ above allows us to deduce that $`b_{6,2}=\frac{13}{4}`$. Thus we have $`p_5\left(x\right)`$ $`=`$ $`\left(x1\right)\left(x^4{\displaystyle \frac{91}{30}}x^3+{\displaystyle \frac{62}{15}}x^2{\displaystyle \frac{91}{30}}x+1\right)`$ $`p_6\left(x\right)`$ $`=`$ $`\left(x1\right)^2\left(x^4{\displaystyle \frac{37}{15}}x^3+{\displaystyle \frac{13}{4}}x^2{\displaystyle \frac{37}{15}}x+1\right).`$ (6.11) We remark that according to the conjecture (1.8), the expansion of $`S(k,\beta )`$ about $`\beta =0`$ should have the structure $$S(k,\beta )=f_0\left(k/\kappa \right)+\beta f_1\left(k/\kappa \right)+\beta ^2f_2\left(k/\kappa \right)+\mathrm{}$$ (6.12) where $$f_j\left(u\right)=u^j\left(c_{j,0}+c_{j,1}u+\mathrm{}\right).$$ (6.13) Consideration of the analysis of reveals that the structure (6.12) will indeed result from the weak coupling expansion, however the structure (6.13) is not immediately evident. (Of course the explicit form $`f_2`$ as revealed by (6.7) exhibits this structure.) ## 7 Perturbation about $`\beta =2`$ and $`\beta =4`$ A feature of the couplings $`\beta =1,2`$ and 4 is that the $`n`$-particle distribution functions are known for each $`n=2,3,\mathrm{}`$ . Introducing the dimensionless distribution $$g(x_1,\mathrm{},x_n):=\rho _{\left(n\right)}(x_1,\mathrm{},x_n)/\rho ^n$$ we can use our knowledge of $`g(x_1,\mathrm{},x_n)`$ for $`n=2,3`$ and 4 at these specific $`\beta `$ to expand $`g(x_1,x_2)`$ about $`\beta =\beta _0`$ to first order in $`\beta \beta _0`$. Thus with $`\mathrm{\Phi }(x_1,x_2):=\mathrm{log}\left\|x_1x_2\right\|`$ we have $`g(x_1,x_2;\beta )`$ $`=`$ $`g(x_1,x_2)+(\beta \beta _0)\{g(x_1,x_2)\mathrm{\Phi }(x_1,x_2)`$ (7.1) $`2\rho {\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(g(x_1,x_2,x_3)g(x_1,x_2)\right)\mathrm{\Phi }(x_1,x_3)𝑑x_3`$ $`{\displaystyle \frac{1}{2}}\rho ^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}(g(x_1,x_2,x_3,x_4)g(x_1,x_2)g(x_3,x_4)g(x_1,x_2,x_3)`$ $`g(x_1,x_2,x_4)+2g(x_1,x_2))\mathrm{\Phi }(x_3,x_4)dx_3dx_4\}+O((\beta \beta _0)^2)`$ where on the right hand side the dimensionless distributions are evaluated at $`\beta =\beta _0`$. Here we will compute this first order correction, and the corresponding first order correction for $`S(k;\beta )`$, in the cases $`\beta _0=2`$ and $`\beta _0=4`$ (we do not consider $`\beta _0=1`$ because of its relation to $`\beta _0=4`$ via the functional equation (1.7)). Now, in the case $`\beta _0=2`$ we have $$g(x_1,\mathrm{},x_n)=det\left[P_2(x_j,x_k)\right]_{j,k=1,\mathrm{},n},P_2(x,y):=\frac{\mathrm{sin}\pi \rho \left(xy\right)}{\pi \rho \left(xy\right)}$$ (7.2) while in the case $`\beta _0=4`$ $$g(x_1,\mathrm{},x_n)=\mathrm{qdet}\left[P_4(x_j,x_k)\right]_{j,k=1,\mathrm{},n}$$ (7.3) where $$P_4(x_j,x_k)=\left[\begin{array}{cc}\frac{\mathrm{sin}2\pi \rho x_{jk}}{2\pi \rho x_{jk}}& \mathrm{Si}\left(2\pi \rho x_{jk}\right)\\ \frac{1}{2\pi \rho }\frac{d}{dx_{jk}}\left(\frac{\mathrm{sin}2\pi \rho x_{jk}}{2\pi \rho x_{jk}}\right)& \frac{\mathrm{sin}2\pi \rho x_{jk}}{2\pi \rho x_{jk}}\end{array}\right]$$ (7.4) with $`x_{jk}:=x_jx_k`$ and Si$`\left(x\right)`$ denoting the complimentary sine integral, defined in terms of the sine integral si$`\left(x\right)`$ by $$\mathrm{Si}\left(x\right)=_0^x\frac{\mathrm{sin}t}{t}𝑑t=\frac{\pi }{2}+\mathrm{si}\left(x\right),\mathrm{si}\left(x\right):=_x^{\mathrm{}}\frac{\mathrm{sin}t}{t}𝑑t.$$ (7.5) In (7.3) qdet denotes quaternion determinant, which can be defined as $$\mathrm{qdet}\left[P_4(x_j,x_k)\right]_{j,k=1,\mathrm{},n}=\underset{PS_n}{}\left(1\right)^{nl}\underset{1}{\overset{l}{}}\left(P_4(x_a,x_b)P_4(x_b,x_c)\mathrm{}P_4(x_d,x_a)\right)^{\left(0\right)}$$ (7.6) where the superscript $`\left(0\right)`$ denotes the operation $`\frac{1}{2}\mathrm{Tr}`$, $`P`$ is any permutation of the indicies $`(1,\mathrm{},n)`$ consisting of $`l`$ exclusive cycles of the form $`\left(abc\mathrm{}da\right)`$ and $`\left(1\right)^{nl}`$ is equal to the parity of $`P`$. Note that this reproduces the definition of an ordinary determinant in the case that $`P_4`$ is a multiple of the identity. The task now is to substitute (7.2) in the case $`\beta _0=2`$ and (7.3) in the case $`\beta _0=4`$, and to compute the integrals. Consider first the case $`\beta _0=2`$. After some calculation (see Appendix A) we find $`g(0,x;\beta )`$ $`=`$ $`1\left({\displaystyle \frac{\mathrm{sin}\pi \rho x}{\pi \rho x}}\right)^2+(\beta 2)\{{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\mathrm{sin}\pi \rho x}{\pi \rho x}}\right)^2{\displaystyle \frac{\mathrm{sin}2\pi \rho x}{2\pi \rho x}}+\mathrm{ci}\left(2\pi \rho x\right)`$ (7.7) $`+{\displaystyle \frac{1}{2\left(\pi \rho x\right)^2}}(\left(\mathrm{log}2\pi \rho \right\|x\|+C)\mathrm{cos}2\pi \rho x\mathrm{ci}\left(2\pi \rho x\right))\}+O((\beta 2)^2)`$ where $`C`$ denotes Euler’s constant while $$\mathrm{ci}\left(x\right)=C+\mathrm{log}\left\|x\right\|+_0^x\frac{\mathrm{cos}t1}{t}𝑑t=_x^{\mathrm{}}\frac{\mathrm{cos}t}{t}𝑑t$$ (7.8) denotes the cosine integral. From this we can compute (again see Appendix A) that up to terms $`O\left(\left(\beta 2\right)^2\right)`$ $$S(k;\beta )=\{\begin{array}{cc}\frac{\left\|k\right\|}{2\pi }+\left(\beta 2\right)\rho \left\{\frac{1}{2}\mathrm{log}\left(1\frac{k^2}{\left(2\pi \rho \right)^2}\right)+\frac{\left\|k\right\|}{4\pi \rho }\mathrm{log}\frac{2\pi \rho +\left\|k\right\|}{2\pi \rho \left\|k\right\|}\frac{\left\|k\right\|}{4\pi \rho }\right\},\hfill & \left\|k\right\|<2\pi \rho ,\hfill \\ \rho +\left(\beta 2\right)\rho \left\{\frac{1}{2}\mathrm{log}\frac{\left\|k\right\|+2\pi \rho }{\left\|k\right\|2\pi \rho }+\frac{\left\|k\right\|}{4\pi \rho }\mathrm{log}\left(1\frac{\left(2\pi \rho \right)^2}{k^2}\right)\frac{\pi \rho }{\left\|k\right\|}\right\},\hfill & \left\|k\right\|>2\pi \rho .\hfill \end{array}$$ (7.9) Let us consider the consequence of (7.9) in regards to the expansion (1.8). For $`\left\|k\right\|<2\pi \rho `$ we observe that all terms but the one proportional to $`\left\|k\right\|`$ are even in $`k`$. This is consistent with $`p_j\left(x\right)`$ having the quadratic factor $`\left(x1\right)^2`$ for $`j`$ odd (recall (5.13)), but only a linear factor for $`j`$ even (recall (1.11)). Moreover, we can use (7.9) to derive a linear equation for the coefficients $`\left\{\stackrel{~}{a}_j\right\}`$. First we differentiate (7.9) with respect to $`\beta `$, set $`\beta =2`$ and expand about $`k=0`$ to obtain $$\frac{S(k;\beta )}{\beta }\|_{\beta =2}=\frac{1}{2}\frac{\left\|k\right\|}{2\pi \rho }+\underset{j=1}{\overset{\mathrm{}}{}}\frac{1}{2j\left(2j1\right)}\left(\frac{\left\|k\right\|}{2\pi \rho }\right)^{2j},\left\|k\right\|<2\pi \rho .$$ Recalling (1.8) and (5.13) this in turn implies $$\frac{1}{2j\left(2j1\right)}=\frac{1}{2}\left(2\stackrel{~}{a}_{2j1,0}+2\stackrel{~}{a}_{2j1,1}+\mathrm{}+2\stackrel{~}{a}_{2j1,j2}+\stackrel{~}{a}_{2j1,j1}\right).$$ (7.10) In the case $`j=4`$ we deduce from this equation, (5.20), (6.10) and (5.16) that $$p_7\left(x\right)=\left(x1\right)\left(1\frac{1607}{420}x+\frac{2011}{280}x^2\frac{911}{105}x^3+\frac{2011}{280}x^4\frac{1607}{420}x^5+x^6\right).$$ (7.11) Consider now the case $`\beta _0=4`$. Due to $`P_4`$ in (7.3) being a $`2\times 2`$ matrix, the calculation required to compute (7.1) is more lengthy and tedius than in the case $`\beta _0=2`$, although the the common structure of $`n`$-point distributions means the two cases are analogous. Some details are given in Appendix B. Our final expression for $`g(x_1,x_2;\beta )`$ is given by (B.8). We find its Fourier transform can be computed explicitly in terms of elementary functions, together with the dilogarithm $$\mathrm{dilog}\left(x\right):=_1^x\frac{\mathrm{log}t}{1t}𝑑t.$$ (7.12) Explicitly, with $`\rho =1`$ for notational convenience, up to terms $`O\left(\left(\beta 4\right)^2\right)`$ $$S(k,\beta )=S(k,4)+\left(\beta 4\right)\left(\frac{\pi }{\left\|k\right\|}+\widehat{B}_0\left(k\right)+2\widehat{B}_1\left(k\right)4\widehat{B}_3\left(k\right)+2\widehat{B}_5\left(k\right)+\widehat{B}_6\left(k\right)\widehat{B}_7\left(k\right)\right)$$ (7.13) where $`\widehat{B}_0\left(k\right)`$ $`=`$ $`{\displaystyle \frac{3}{2}}+{\displaystyle \frac{3\left\|k\right\|}{8\pi }}+{\displaystyle \frac{\left\|k\right\|}{4\pi }}\mathrm{log}\left({\displaystyle \frac{4\pi +\left\|k\right\|}{\left\|k\right\|}}\right)+\left(C+{\displaystyle \frac{1}{2}}\mathrm{log}\left(16\pi ^2k^2\right)\right)\left(1{\displaystyle \frac{\left\|k\right\|}{4\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right)`$ (7.14) $`+{\displaystyle \frac{\left\|k\right\|}{16\pi }}\left(\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)\mathrm{dilog}\left({\displaystyle \frac{4\pi +\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\mathrm{log}\left({\displaystyle \frac{4\pi +\left\|k\right\|}{\left\|k\right\|}}\right)+g_1\left(k\right)\right)`$ $`+{\displaystyle \frac{2\pi \left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|,\left\|k\right\|<4\pi ,`$ $`\widehat{B}_0\left(k\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{\left\|k\right\|+4\pi }{\left\|k\right\|4\pi }}\right)+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left({\displaystyle \frac{k^216\pi ^2}{k^2}}\right)+{\displaystyle \frac{\left\|k\right\|}{16\pi }}(\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{\left\|k\right\|+2\pi }}\right)`$ (7.15) $`+\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{\left\|k\right\|2\pi }}\right)\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|+4\pi }{\left\|k\right\|+2\pi }}\right)\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|4\pi }{\left\|k\right\|2\pi }}\right)),\|k\|>4\pi ,`$ $$\widehat{B}_1\left(k\right)=\{\begin{array}{cc}\frac{\pi }{\left\|k\right\|}\left(1\frac{\left\|k\right\|}{4\pi }+\frac{\left\|k\right\|}{8\pi }\mathrm{log}\left\|1\frac{\left\|k\right\|}{2\pi }\right\|\right),\hfill & \left\|k\right\|<4\pi \hfill \\ 0,\hfill & \left\|k\right\|>4\pi ,\hfill \end{array}$$ (7.16) $`\widehat{B}_3\left(k\right)`$ $`=`$ $`{\displaystyle \frac{3}{2}}+{\displaystyle \frac{3\left\|k\right\|}{8\pi }}+C\left(1{\displaystyle \frac{\left\|k\right\|}{4\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right)+\left({\displaystyle \frac{1}{8}}{\displaystyle \frac{3\left\|k\right\|}{32\pi }}\right)\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|`$ $`+{\displaystyle \frac{\left\|k\right\|}{64\pi }}\left(\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right)^2+{\displaystyle \frac{1}{8\pi }}\left(4\pi \left\|k\right\|\right)\mathrm{log}\left(4\pi \left\|k\right\|\right){\displaystyle \frac{1}{8\pi }}\left\|k\right\|\mathrm{log}\left\|k\right\|+{\displaystyle \frac{1}{2}}\mathrm{log}4\pi `$ $`+{\displaystyle \frac{\left\|k\right\|}{32\pi }}(\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)+{\displaystyle \frac{\pi ^2}{12}}\mathrm{dilog}\left({\displaystyle \frac{4\pi }{2\pi +\left\|k\right\|}}\right)\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{2\pi }}\right)\mathrm{dilog}\left({\displaystyle \frac{4\pi \left\|k\right\|}{2\pi }}\right)`$ $`+2\mathrm{log}\left(2\pi \right)\mathrm{log}\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\|+\mathrm{log}(2\pi +\|k\left\|\right)\mathrm{log}\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\|+g_2\left(k\right)),\|k\|<4\pi ,`$ $`\widehat{B}_3\left(k\right)`$ $`=`$ $`0,\left\|k\right\|>4\pi ,`$ (7.17) $$\widehat{B}_5\left(k\right)=\{\begin{array}{cc}\widehat{B}_3\left(k\right)\frac{\left\|k\right\|}{128\pi }\left(\mathrm{log}\left\|1\frac{\left\|k\right\|}{2\pi }\right\|\right)^2+\frac{\left\|k\right\|}{32\pi }g_3\left(k\right),\hfill & \left\|k\right\|<4\pi \hfill \\ 0,\hfill & \left\|k\right\|>4\pi \hfill \end{array}$$ (7.18) $`\widehat{B}_6\left(k\right)`$ $`=`$ $`{\displaystyle \frac{3}{2}}+{\displaystyle \frac{3\left\|k\right\|}{8\pi }}{\displaystyle \frac{\left\|k\right\|}{16\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|+\left(C+\mathrm{log}\left(4\pi \left\|k\right\|\right)\right)\left(1{\displaystyle \frac{\left\|k\right\|}{4\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right)`$ $`+{\displaystyle \frac{\left\|k\right\|}{32\pi }}({\displaystyle \frac{\pi ^2}{3}}\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{2\pi }}\right)\mathrm{log}\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\|\mathrm{log}\left({\displaystyle \frac{\left\|k\right\|}{2\pi }}\right)2\mathrm{d}\mathrm{i}\mathrm{l}\mathrm{o}\mathrm{g}\left({\displaystyle \frac{4\pi \left\|k\right\|}{2\pi }}\right)`$ $`2\mathrm{log}\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\|\mathrm{log}\left({\displaystyle \frac{4\pi \left\|k\right\|}{2\pi }}\right)+g_4\left(k\right)),\|k\|<4\pi ,`$ $`\widehat{B}_6\left(k\right)`$ $`=`$ $`0,\left\|k\right\|>4\pi ,`$ (7.19) $$\widehat{B}_7\left(k\right)=\{\begin{array}{cc}\frac{\pi }{\left\|k\right\|}\left(1\frac{\left\|k\right\|}{4\pi }+\frac{\left\|k\right\|}{8\pi }\mathrm{log}\left\|1\frac{\left\|k\right\|}{2\pi }\right\|\right)^2,\hfill & \left\|k\right\|<4\pi \hfill \\ 0,\hfill & \left\|k\right\|>4\pi ,\hfill \end{array}$$ (7.20) with $`g_1\left(k\right)`$ $`=`$ $`\{\begin{array}{cc}\mathrm{dilog}\left(\frac{4\pi \left\|k\right\|}{2\pi \left\|k\right\|}\right)\mathrm{dilog}\left(\frac{2\pi }{2\pi \left\|k\right\|}\right)\frac{\pi ^2}{6}\mathrm{log}\left(1\frac{\left\|k\right\|}{2\pi }\right)\mathrm{log}\left(\frac{4\pi \left\|k\right\|}{2\pi \left\|k\right\|}\right),\hfill & \left\|k\right\|<2\pi \hfill \\ \mathrm{dilog}\left(\frac{\left\|k\right\|}{\left\|k\right\|2\pi }\right)\mathrm{dilog}\left(\frac{2\pi }{\left\|k\right\|2\pi }\right)\frac{\pi ^2}{6}+\mathrm{log}\left(\frac{2\pi }{\left\|k\right\|2\pi }\right)\mathrm{log}\left(\frac{\left\|k\right\|}{\left\|k\right\|2\pi }\right),\hfill & 2\pi <\left\|k\right\|<4\pi ,\hfill \end{array}`$ (7.23) $`g_2\left(k\right)`$ $`=`$ $`\{\begin{array}{cc}\mathrm{dilog}\left(\frac{4\pi \left\|k\right\|}{2\pi \left\|k\right\|}\right)\frac{\pi ^2}{6}+\mathrm{log}\left(2\pi \left\|k\right\|\right)\mathrm{log}\left(1\frac{\left\|k\right\|}{2\pi }\right),\hfill & \left\|k\right\|<2\pi \hfill \\ \mathrm{dilog}\left(\frac{2\pi }{\left\|k\right\|2\pi }\right)+\mathrm{log}\left(4\pi \left\|k\right\|\right)\mathrm{log}\left(\frac{\left\|k\right\|}{2\pi }1\right),\hfill & 2\pi <\left\|k\right\|<4\pi \hfill \end{array}`$ (7.26) $`g_3\left(k\right)`$ $`=`$ $`\{\begin{array}{cc}\frac{1}{2}\left(\frac{\pi ^2}{6}2\mathrm{d}\mathrm{i}\mathrm{l}\mathrm{o}\mathrm{g}\left(\frac{2\pi }{\left\|k\right\|}\right)\mathrm{log}\left(\frac{2\pi }{\left\|k\right\|}\right)\mathrm{log}\left(\frac{\left(2\pi \left\|k\right\|\right)^2}{2\pi \left\|k\right\|}\right)\right),\hfill & \left\|k\right\|<2\pi ,\hfill \\ \frac{1}{2}\left(\mathrm{dilog}\left(\frac{\left\|k\right\|2\pi }{2\pi }\right)\mathrm{dilog}\left(\frac{2\pi }{\left\|k\right\|}\right)+\mathrm{log}\left(\frac{\left\|k\right\|2\pi }{2\pi }\right)\mathrm{log}\left(\frac{\left\|k\right\|}{2\pi }\right)\right),\hfill & 2\pi <\left\|k\right\|<4\pi ,\hfill \end{array}`$ (7.29) $`g_4\left(k\right)`$ $`=`$ $`\{\begin{array}{c}\frac{\pi ^2}{6}\mathrm{dilog}\left(\frac{2\pi }{2\pi \left\|k\right\|}\right)+2\mathrm{d}\mathrm{i}\mathrm{l}\mathrm{o}\mathrm{g}\left(\frac{4\pi \left\|k\right\|}{2\pi \left\|k\right\|}\right)+2\mathrm{log}\left(\frac{2\pi }{2\pi \left\|k\right\|}\right)\mathrm{log}\left(\frac{4\pi \left\|k\right\|}{2\pi \left\|k\right\|}\right),\left\|k\right\|<2\pi \hfill \\ \mathrm{dilog}\left(\frac{\left\|k\right\|}{\left\|k\right\|2\pi }\right)2\mathrm{d}\mathrm{i}\mathrm{l}\mathrm{o}\mathrm{g}\left(\frac{2\pi }{\left\|k\right\|2\pi }\right)+\mathrm{log}\left(\frac{2\pi }{\left\|k\right\|2\pi }\right)\mathrm{log}\left(\frac{\left\|k\right\|}{\left\|k\right\|2\pi }\right),\mathrm{\hspace{0.25em}\hspace{0.25em}2}\pi <\left\|k\right\|<4\pi .\hfill \end{array}`$ (7.32) The above formula for $`S(k;\beta )`$ in the case $`\left\|k\right\|<2\pi `$ (recall here $`\rho =1`$) can be used to expand $`S(k;\beta )/\beta `$ about $`k=0`$. For this task we use computer algebra, which gives the result $`{\displaystyle \frac{S(k;\beta )}{\beta }}\|_{\beta =4}`$ $`=`$ $`{\displaystyle \frac{\left\|k\right\|}{16\pi }}+{\displaystyle \frac{\left\|k\right\|^3}{256\pi ^3}}+{\displaystyle \frac{5k^4}{3072\pi ^4}}+{\displaystyle \frac{3\left\|k\right\|^5}{4096\pi ^5}}+{\displaystyle \frac{27k^6}{81920\pi ^6}}`$ (7.34) $`+{\displaystyle \frac{37\left\|k\right\|^7}{245760\pi ^7}}+{\displaystyle \frac{1273k^8}{18350080\pi ^8}}+{\displaystyle \frac{887\left\|k\right\|^9}{27525120\pi ^9}}+{\displaystyle \frac{4423k^{10}}{293601280\pi ^{10}}}`$ $`+{\displaystyle \frac{1949\left\|k\right\|^{11}}{275251200\pi ^{11}}}+\mathrm{}`$ This allows us to deduce a further equation for $`\left\{b_{8,j}\right\}_{j=0,\mathrm{},4}`$ and $`\left\{\stackrel{~}{a}_{9,j}\right\}_{j=0,\mathrm{},4}`$, which in combination with (7.10), (6.10), (5.20), (5.15) and (5.16) implies $`p_8\left(x\right)`$ $`=`$ $`\left(x1\right)^2\left(1{\displaystyle \frac{263}{84}}x+{\displaystyle \frac{1697}{315}}x^2{\displaystyle \frac{6337}{1008}}x^3+{\displaystyle \frac{1697}{315}}x^4{\displaystyle \frac{263}{84}}x^5+x^6\right)`$ (7.35) $`p_9\left(x\right)`$ $`=`$ $`\left(x1\right)\left(1{\displaystyle \frac{791}{180}}x+{\displaystyle \frac{73603}{7560}}x^2{\displaystyle \frac{7355}{504}}x^3+{\displaystyle \frac{2231}{135}}x^4{\displaystyle \frac{7355}{504}}x^5+{\displaystyle \frac{73603}{7560}}x^6{\displaystyle \frac{791}{180}}x^7+x^8\right)`$ ## 8 Conclusion Collecting together the evaluations (1.12), (3.15), (5.21), (4.19), (6), (7.11), (7.35) and (LABEL:p9), and substituting in (1.8) we have that for $`\left\|k\right\|<\mathrm{min}(2\pi \rho ,\pi \beta \rho )`$ $`{\displaystyle \frac{\pi \beta }{\left\|k\right\|}}S(k;\beta )=`$ (8.1) $`1`$ $`+\left(x1\right)y`$ $`+\left(x1\right)^2y^2`$ $`+\left(x1\right)\left(x^2{\displaystyle \frac{11}{6}}x+1\right)y^3`$ $`+\left(x1\right)^2\left(x^2{\displaystyle \frac{3}{2}}x+1\right)y^4`$ $`+\left(x1\right)\left(x^4{\displaystyle \frac{91}{30}}x^3+{\displaystyle \frac{62}{15}}x^2{\displaystyle \frac{91}{30}}x+1\right)y^5`$ $`+\left(x1\right)^2\left(x^4{\displaystyle \frac{37}{15}}x^3+{\displaystyle \frac{13}{4}}x^2{\displaystyle \frac{37}{15}}x+1\right)y^6`$ $`+\left(x1\right)\left(x^6{\displaystyle \frac{1607}{420}}x^5+{\displaystyle \frac{2011}{280}}x^4{\displaystyle \frac{911}{105}}x^3+{\displaystyle \frac{2011}{280}}x^2{\displaystyle \frac{1607}{420}}x+1\right)y^7`$ $`+\left(x1\right)^2\left(x^6{\displaystyle \frac{263}{84}}x^5+{\displaystyle \frac{1697}{315}}x^4{\displaystyle \frac{6337}{1008}}x^3+{\displaystyle \frac{1697}{315}}x^2{\displaystyle \frac{263}{84}}x+1\right)y^8`$ $`+\left(x1\right)\left(x^8{\displaystyle \frac{791}{180}}x^7+{\displaystyle \frac{73603}{7560}}x^6{\displaystyle \frac{7355}{504}}x^5+{\displaystyle \frac{2231}{135}}x^4{\displaystyle \frac{7355}{504}}x^3+{\displaystyle \frac{73603}{7560}}x^2{\displaystyle \frac{791}{180}}x+1\right)y^9`$ $`+O\left(y^{10}\right)`$ where $`x=\beta /2`$ and $`y=\left\|k\right\|/\pi \beta \rho `$. With the coefficient of $`y^j`$ denoted $`p_j\left(x\right)`$ as has been throughout, we recall from our workings above that $`p_0\left(x\right)`$, $`p_1\left(x\right)`$, $`p_2\left(x\right)`$ and $`p_4\left(x\right)`$ have been calculated for general values of $`\beta `$. In all other cases the calculation has relied on the assumption that the $`p_j\left(x\right)`$ are indeed polynomials. On this point we remark that in such cases, excluding $`j=8`$ and 9, we have more data points than is necessary to uniquely specify $`p_j\left(x\right)`$, assuming it is a polynomial, and our extra data points are consistent with the explicit forms presented in (8.1). We remark that the structure exhibited by (8.1) is familiar from the study of exactly solvable two-dimensional lattice models . In this field one encounters two-variable generating functions $`G(x,y)`$ say with series expansions of the form $$G(x,y)=\underset{n=0}{\overset{\mathrm{}}{}}H_n\left(x\right)y^n$$ (8.2) in which $`H_n\left(x\right)`$ is a rational function, and furthermore the denominator polynomial in $`H_n\left(x\right)`$ only has a small number of (typically no more than two) distinct zeros. For example, the two-dimensional Ising model with couplings $`J_1`$ ($`J_2`$) between bonds in the horizontal (vertical) direction and $`x:=\mathrm{exp}\left(4J_1/k_BT\right)`$, $`y:=\mathrm{exp}\left(4J_2/k_BT\right)`$ has for its spontateous magnetisation the celebrated exact expression (see e.g. ) $$M(x,y)=\left(1\frac{16xy}{\left(1x\right)^2\left(1y\right)^2}\right)^{1/8}.$$ (8.3) When written in the form (8.2) one finds $$H_n\left(x\right)=\frac{2xP_n\left(x\right)}{\left(1x\right)^n}$$ (8.4) where $`P_n\left(x\right)`$ is a polynomial of degree $`2n2`$ which satisfies the functional relation $$P_n\left(x\right)=x^{2n2}P_n\left(1/x\right).$$ (8.5) As emphasized in , the exact solution (8.3) can be uniquely determined by the functional form (8.4), together with the functional (inversion) relation (8.5) and the symmetry relation $`M(x,y)=M(y,x)`$. For the structure function of the log-gas we have no analogue of the symmetry relation and so cannot characterize (2.6) this way. One immediate feature of the polynomials $`p_j\left(x\right)`$ in (8.1) is that for $`j`$ even the polynomial $`p_j\left(x\right)`$ has all coefficients positive, while for $`j`$ odd the polynomial $`p_j\left(x\right)`$ has all coefficients negative. Another general feature of the $`p_j\left(x\right)`$ in (8.1), obtained from numerical computation, is that all the zeros lie on the unit circle in the complex $`x`$-plane. This can be rigorously determined numerically because the symmetry (1.9) implies that if $`x_0`$ is a zero of $`p_j\left(x\right)`$, then so is $`1/x_0`$, which will be the complex conjugate of $`x_0`$ if and only if $`\left\|x_0\right\|=1`$. The quantum many body interpretation of (1.1) allows us to give a physical interpretation to the functional relation (1.7). As the functional relation is derived from the integral representation (2.9), it is appropriate to recall the physical interpretation of that formula. In (2.9), with $`\beta /2=p/q`$, there are $`q`$ integrals over $`x_i(0,\mathrm{})`$ and $`p`$ integrals over $`y_j(0,1)`$. The variables $`x_i`$ can be interpreted as being rapidities of quasi-particle excitations, while the $`y_j`$ are rapidities of quasi-hole excitations. Thus the transformation $`\beta 4/\beta `$ is equivalent to interchanging $`p`$ and $`q`$ and thus the quasi-holes and quasi-particles. In (2.9) this does not lead to an integral of the same functional form as before; although the functional form of the integrand is conserved, apart from a renormalization of $`k`$, the domain of integration is different for $`\left\{x_i\right\}`$ and $`\left\{y_j\right\}`$. But with $`k`$ restricted as in (2.8) both sets of variables can take any value in $`(0,\mathrm{})`$. The quasi-particles and quasi-holes play an identical role and the functional equation results. It is of interest to consider the small $`k`$ expansion of $`S(k;\mathrm{\Gamma })`$, $`\mathrm{\Gamma }:=q^2/k_BT`$ ($`q`$ = charge), for the two-dimensional one-component plasma. As mentioned earlier, this has recently been the object of study of Kalinay et al. . They obtain results which imply $$\frac{2\pi \mathrm{\Gamma }}{k^2}S(k;\mathrm{\Gamma })=1+\left(\frac{\mathrm{\Gamma }}{4}1\right)\frac{k^2}{2\pi \mathrm{\Gamma }\rho }+\left(\frac{\mathrm{\Gamma }}{4}\frac{3}{2}\right)\left(\frac{\mathrm{\Gamma }}{4}\frac{2}{3}\right)\left(\frac{k^2}{2\pi \mathrm{\Gamma }\rho }\right)^2+O\left(k^6\right).$$ (8.6) where $`k:=\left\|\stackrel{}{k}\right\|`$. The structure of (8.6) bears a striking resemblence to (8.1) with $`\mathrm{\Gamma }/4`$ corresponding to $`x`$ and $`k^2/2\pi \mathrm{\Gamma }\rho `$ to $`y`$. In particular with $`g(x,y):=\left(2\pi \mathrm{\Gamma }/k^2\right)S(k;\mathrm{\Gamma })`$, the expansion (8.6) to the given order is such that $$g(x,y)=g(\frac{1}{x};yx).$$ (8.7) Furthermore, writing $$g(x,y)=1+\underset{l=1}{\overset{\mathrm{}}{}}u_l\left(x\right)y^l$$ (8.8) we have $`u_1\left(x\right)=\left(x1\right)`$, $`u_2\left(x\right)=\left(x3/2\right)\left(x2/3\right)`$ so $`u_j\left(x\right)`$ is a monic $`j`$th degree polynomial for $`j2`$. However we can demonstrate that this analogy breaks down for the $`l=3`$ term in (8.8). To demonstrate this fact, suppose instead that the functional equation (8.8) was valid at order $`l=3`$ in (8.8) and $`u_3\left(x\right)`$ is a monic polynomial. Then $`u_3`$ must be of the form $$u_3\left(x\right)=\left(x1\right)\left(x^2+ax+1\right).$$ (8.9) From the definition of $`g(x,y)`$ we can check that this is equivalent to the statement that $$\frac{1}{\rho }\left(\frac{\pi \mathrm{\Gamma }\rho }{2}\right)^4_{𝐑^2}r^8S(r;\mathrm{\Gamma })𝑑\stackrel{}{r}=\left(4!\right)^2\left(x1\right)\left(x^2+ax+1\right).$$ (8.10) But as noted in , it follows from the perturbation expansion of that $`{\displaystyle \frac{1}{\rho }}\left({\displaystyle \frac{\pi \mathrm{\Gamma }\rho }{2}}\right)^4{\displaystyle _{𝐑^2}}r^8S(r;\mathrm{\Gamma })𝑑\stackrel{}{r}`$ $`=`$ $`4!+\left(\mathrm{\Gamma }2\right)4!\left({\displaystyle \underset{k=0}{\overset{4}{}}}{\displaystyle \frac{2^k1}{k+1}}2\right)+O\left(\left(\mathrm{\Gamma }2\right)^2\right)`$ (8.11) $`=`$ $`4!+\left(\mathrm{\Gamma }2\right)4!{\displaystyle \frac{17}{4}}+O\left(\left(\mathrm{\Gamma }2\right)^2\right).`$ The term in (8.11) proportional to $`\mathrm{\Gamma }2`$ is incompatible with (8.10) which gives instead $$\left(\mathrm{\Gamma }2\right)4!\frac{18}{4}$$ independent of the value of $`a`$. Indeed in evidence is presented which indicates $`u_3\left(x\right)`$ is an infinite series in $`x`$, although we have no way of determining if the functional equation (8.7) also breaks down at this order. ## Acknowledgements The work of PJF and DSM was supported by the Australian Research Council. ## Appendix A In this appendix some details of the derivation of (7.7) and (7.9) will be given. To simplify notation we take $`\rho =1`$ throughout. The first step is to substitute (7.2) and (7.1) and simplify by expanding out the determinant and cancelling terms where possible. This shows that up to terms $`O\left(\left(\beta 2\right)^2\right)`$ $`g_2(x_1,x_2;\beta )`$ $`=1\left(P_2(x_1,x_2)\right)^2+(\beta 2)\{(1\left(P_2(x_1,x_2)\right)^2)\mathrm{\Phi }(x_1,x_2)`$ $`2{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(\left(P_2(x_2,x_3)\right)^2\left(P_2(x_1,x_3)\right)^2+2P_2(x_1,x_2)P_2(x_2,x_3)P_2(x_3,x_1)\right)\mathrm{\Phi }(x_1,x_3)𝑑x_3`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}(4P_2(x_1,x_3)P_2(x_3,x_4)P_2(x_4,x_1)4P_2(x_1,x_2)P_2(x_2,x_3)P_2(x_3,x_4)P_2(x_4,x_1)`$ $`2P_2(x_1,x_3)P_2(x_3,x_2)P_2(x_2,x_4)P_2(x_4,x_1)+2\left(P_2(x_1,x_3)\right)^2\left(P_2(x_2,x_4)\right)^2)\mathrm{\Phi }(x_3,x_4)dx_3dx_4\}.`$ The convolution structure $$_{\mathrm{}}^{\mathrm{}}f\left(y_1x\right)g\left(xy_2\right)𝑑x$$ often occurs in the above integrals. Such an integral can be transformed by introducing the Fourier transforms $`\widehat{f}`$ ($`\widehat{g}`$) according to the formula $$_{\mathrm{}}^{\mathrm{}}f\left(y_1x\right)g\left(xy_2\right)𝑑x=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}\widehat{f}\left(l\right)\widehat{g}\left(l\right)e^{il\left(y_1y_2\right)}𝑑l.$$ (A.2) Making use of this formula typically leads to simplifications. For example, consider the first integral in (Appendix A). Starting with the Fourier transform $$_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{sin}^2\pi x}{\left(\pi x\right)^2}e^{ikx}𝑑x=\{\begin{array}{cc}1\frac{\left\|k\right\|}{2\pi },\hfill & \left\|k\right\|<2\pi \hfill \\ 0,\hfill & \left\|k\right\|>2\pi \hfill \end{array}$$ and (5.18), application of (A.2) gives $$A_1\left(x_{12}\right):=_{\mathrm{}}^{\mathrm{}}\left(P_2(x_2,x_3)\right)^2\mathrm{\Phi }(x_3,x_1)𝑑x_3=_{2\pi }^{2\pi }\left(1\frac{\left\|k\right\|}{2\pi }\right)\frac{\pi }{\left\|k\right\|}\mathrm{cos}kx_{12}\frac{dk}{2\pi }.$$ (A.3) This expression is indeed simpler than the original, but it suffers from being ill-defined, due to the singularity at the origin. However its derivative is well-defined, and can furthermore be evaluated in terms of elementary functions giving $$\frac{d}{dx}A_1\left(x\right)=\frac{\mathrm{sin}2\pi x}{2\pi x^2}\frac{1}{x}.$$ (A.4) Also, we have $$A_1\left(0\right)=_{\mathrm{}}^{\mathrm{}}𝑑x\frac{\mathrm{sin}^2\pi x}{\left(\pi x\right)^2}\mathrm{log}\left\|x\right\|=C+\mathrm{log}2\pi 1,$$ (A.5) where $`C`$ denotes Euler’s constant. Together (A.4) and (A.5) imply $$A_1\left(x\right)=\frac{\mathrm{sin}2\pi x}{2\pi x}+\mathrm{ci}\left(2\pi x\right)\mathrm{log}\left\|x\right\|$$ (A.6) where ci$`\left(x\right)`$ denotes the cosine integral (7.8). The other six integrals in (Appendix A) yield to similar techniques. We find $`A_2`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_2(x_1,x_3)\right)^2\mathrm{\Phi }(x_1,x_3)𝑑x_3=A_1\left(0\right)`$ $`A_3\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_2(x_2,x_3)P_2(x_3,x_1)\mathrm{\Phi }(x_3,x_1)𝑑x_3={\displaystyle _\pi ^\pi }\left(C+\mathrm{log}\left(\pi +k\right)\right)\mathrm{cos}kx_{12}{\displaystyle \frac{dk}{2\pi }}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(C+\mathrm{log}2\pi \mathrm{log}\left\|x_{12}\right\|+\mathrm{ci}\left(2\pi x_{12}\right)\right){\displaystyle \frac{\mathrm{sin}\pi x_{12}}{\pi x_{12}}}{\displaystyle \frac{1}{2}}\left(\mathrm{si}\left(2\pi x_{12}\right)+{\displaystyle \frac{\pi }{2}}\right){\displaystyle \frac{\mathrm{cos}\pi x_{12}}{\pi x_{12}}}`$ $`A_4`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_2(x_1,x_3)P_2(x_3,x_4)P_2(x_4,x_1)\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4=A_3\left(0\right)=A_1\left(0\right),`$ $`A_5\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_2(x_2,x_3)P_2(x_3,x_4)P_2(x_4,x_1)\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4=A_3\left(x_{12}\right),`$ $`A_6\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_2(x_1,x_3)P_2(x_3,x_2)P_2(x_2,x_4)P_2(x_4,x_1)\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4`$ $`=`$ $`{\displaystyle \frac{1}{2\left(\pi x_{12}\right)^2}}\left(C+\mathrm{log}2\pi +\mathrm{cos}2\pi x_{12}\left(\mathrm{log}\left\|x_{12}\right\|\mathrm{ci}\left(2\pi x_{12}\right)\right)\mathrm{sin}2\pi x_{12}\left(\mathrm{si}\left(2\pi x_{12}\right)+{\displaystyle \frac{\pi }{2}}\right)\right),`$ $`A_7\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_2(x_1,x_3)\right)^2\left(P_2(x_2,x_4)\right)^2\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4={\displaystyle _{2\pi }^{2\pi }}{\displaystyle \frac{dk}{2\pi }}\left(1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right)^2{\displaystyle \frac{\pi }{\left\|k\right\|}}\mathrm{cos}kx_{12}`$ (A.7) $`=`$ $`\mathrm{log}\left\|x_{12}\right\|{\displaystyle \frac{1\mathrm{cos}2\pi x_{12}}{\left(2\pi x_{12}\right)^2}}{\displaystyle \frac{\mathrm{sin}2\pi x_{12}}{2\pi x_{12}}}+\mathrm{ci}\left(2\pi x_{12}\right),`$ where si$`\left(x\right)`$ denotes the sine integral defined in (7.5). Of the results (Appendix A), the evaluation of $`A_6`$ is the most difficult, so it is appropriate to give details in that case also. We observe that $`A_6`$ consists of the convolution of $`P_2(x_1,x_3)P_2(x_3,x_2)`$ regarded as a function of $`x_3`$, and $`\mathrm{\Phi }(x_3,x_4)`$, and $`P_2(x_4,x_1)P_2(x_2,x_4)`$ regarded as a function of $`x_4`$. It simplifies the calculation to take as the origin in both integrations the centre of the interval between particle 1 and particle 2, which is achieved by the change of variables $`x_3x_3+\left(x_1+x_2\right)/2`$, $`x_4x_4+\left(x_1+x_2\right)/2`$. Use of (A.2) then shows $$A_6\left(x_{12}\right)=_{\mathrm{}}^{\mathrm{}}\frac{dk}{2\pi }\left(\widehat{V}(k,x_{12})\right)^2\frac{\pi }{\left\|k\right\|},$$ (A.8) $$\widehat{V}(k,x_{12}):=_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{sin}\pi x_{13}}{\pi x_{13}}\frac{\mathrm{sin}\pi x_{32}}{\pi x_{32}}\mathrm{cos}kx_3dx_3=\{\begin{array}{cc}\frac{1}{\pi x_{12}}\mathrm{sin}\left(\pi \frac{\left\|k\right\|}{2}\right)x_{12},\hfill & \left\|k\right\|<2\pi \hfill \\ 0,\hfill & \left\|k\right\|>2\pi \hfill \end{array}$$ (A.9) where the second equality in (A.9) follows after further use of (A.2). Thus $$A_6\left(x\right)=_{2\pi }^{2\pi }\frac{dk}{2\pi }\left(\frac{\mathrm{sin}\left(\pi \left\|k\right\|/2\right)x}{\pi x}\right)^2\frac{\pi }{\left\|k\right\|}.$$ (A.10) As in (A.10), this integrand is ill-defined. To proceed further, we write $$A_6\left(x\right)=A_6^{\left(1\right)}\left(x\right)+A_6^{\left(2\right)}\left(x\right)$$ where $`A_6^{\left(1\right)}\left(x\right)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dk}{2\pi }}\left\{\left({\displaystyle \frac{\mathrm{sin}\left(\pi \left\|k\right\|/2\right)x}{\pi x}}\right)^2\left({\displaystyle \frac{\mathrm{sin}\pi x}{\pi x}}\right)^2\right\}{\displaystyle \frac{\pi }{\left\|k\right\|}}`$ $`A_6^{\left(2\right)}\left(x\right)`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{sin}\pi x}{\pi x}}\right)^2{\displaystyle _{2\pi }^{2\pi }}{\displaystyle \frac{dk}{2\pi }}{\displaystyle \frac{\pi }{\left\|k\right\|}}.`$ The integral defining $`A_6^{\left(1\right)}`$ is well defined and can be computed by elementary means. The integral defining $`A_6^{\left(2\right)}`$ is singular. It coincides with the singular part of $`A_1\left(0\right)`$ (recall (A.3)), and so from (A.5) we have $$A_6^{\left(1\right)}\left(0\right)=\left(\frac{\mathrm{sin}\pi x}{\pi x}\right)^2\left(C+\mathrm{log}2\pi \right).$$ Collecting together the above evaluations of $`A_1`$$`A_7`$ and substituting as appropriate in (Appendix A) gives (7.7). The next task is to evaluate the Fourier transform. Now the evaluations of $`A_1`$ and $`A_7`$ are given as Fourier integrals, so their Fourier transform is immediate: $`\mathrm{FT}A_1\left(x\right)`$ $`=`$ $`{\displaystyle \frac{\pi }{\left\|k\right\|}}{\displaystyle \frac{1}{2}},\left\|k\right\|<2\pi `$ $`\mathrm{FT}A_7\left(x\right)`$ $`=`$ $`{\displaystyle \frac{\pi }{\left\|k\right\|}}1+{\displaystyle \frac{\left\|k\right\|}{4\pi }},\left\|k\right\|<2\pi ,`$ (A.11) while for $`\left\|k\right\|>2\pi `$ $$\mathrm{FT}A_1\left(x\right)=\mathrm{FT}A_7\left(x\right)=0.$$ (A.12) We can check that the constants $`A_2`$ and $`A_4`$ cancel when substituted in (Appendix A), and so play no further part in the calculation. Of the remaining terms, consider first the first term proportional to $`\beta 2`$ in (Appendix A), $`A_0\left(x\right)`$ say. Making use of (A.2) we see that $$\mathrm{FT}A_0\left(x\right)=\frac{\pi }{\left\|k\right\|}+_{2\pi }^{2\pi }\frac{dl}{2\pi }\frac{\pi }{\left\|lk\right\|}\left(1\frac{\left\|l\right\|}{2\pi }\right).$$ (A.13) For $`\left\|k\right\|<2\pi `$ minor manipulation allows the singular part $$_{2\pi }^{2\pi }\frac{dl}{2\pi }\frac{\pi }{\left\|l\right\|}=C+\mathrm{log}2\pi $$ (A.14) to be separated, while the remaining convergent integrals are elementary. We thus find that for $`\left\|k\right\|<2\pi `$ $$\mathrm{FT}A_0\left(x\right)=\frac{\pi }{\left\|k\right\|}+\left\{C+\mathrm{log}2\pi +\frac{1}{2}\mathrm{log}\left(1\left(\frac{k}{2\pi }\right)^2\right)\right\}\left(1\frac{\left\|k\right\|}{2\pi }\right)1+\frac{\left\|k\right\|}{2\pi }+\frac{\left\|k\right\|}{2\pi }\mathrm{log}\left(\frac{2\pi +\left\|k\right\|}{\left\|k\right\|}\right).$$ (A.15) For $`\left\|k\right\|>2\pi `$ the integrals in (A.13) are convergent and also elementary. In this case we find $$\mathrm{FT}A_0\left(x\right)=\frac{\pi }{\left\|k\right\|}+\frac{1}{2}\mathrm{log}\frac{\left\|k\right\|+2\pi }{\left\|k\right\|2\pi }+\frac{\left\|k\right\|}{4\pi }\mathrm{log}\left(1\frac{4\pi ^2}{k^2}\right).$$ (A.16) To compute the Fourier transform of $`A_6`$, we begin by making use of (A.2) in (A.8) thereby obtaining $$\mathrm{FT}A_6=_{\mathrm{}}^{\mathrm{}}\frac{dl}{2\pi }_{\mathrm{}}^{\mathrm{}}\frac{dk_1}{2\pi }\widehat{V}(l,k_1)\widehat{V}(l,\left(k_1k\right))\frac{\pi }{\left\|l\right\|}$$ where $$\widehat{V}(l,k):=_{\mathrm{}}^{\mathrm{}}𝑑x\widehat{V}(l,x)e^{ikx}=\chi _{\left\|k\right\|<\pi \left\|l\right\|/2}$$ with the equality in the latter formula following from the explicit form (A.9) of $`\widehat{V}(l,x)`$ and then computation of the resulting integral, and where $`\chi _T=1`$ for $`T`$ true and $`\chi _T=0`$ otherwise. Thus $`\mathrm{FT}A_6\left(x\right)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dl}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dk_1}{2\pi }}\chi _{\left\|k_1\right\|<\pi \left\|l\right\|/2}\chi _{\left\|k_1k\right\|<\pi \left\|l\right\|/2}{\displaystyle \frac{\pi }{\left\|l\right\|}}`$ (A.19) $`=`$ $`\{\begin{array}{cc}\left(1\frac{\left\|k\right\|}{2\pi }\right)\left(C+\mathrm{log}2\pi \right)\left(1\frac{\left\|k\right\|}{2\pi }\right)\mathrm{log}\frac{2\pi }{2\pi \left\|k\right\|}\frac{1}{2\pi }\left(2\pi \left\|k\right\|\right),\hfill & \left\|k\right\|<2\pi \hfill \\ 0,\hfill & \left\|k\right\|>2\pi \hfill \end{array}`$ where use has been made of the generalized integral evaluation (A.14). The final Fourier transform to consider is $`\mathrm{FT}{\displaystyle \frac{\mathrm{sin}\pi x}{\pi x}}A_3\left(x\right)=\mathrm{FT}{\displaystyle \frac{\mathrm{sin}\pi x}{\pi x}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{dk_1}{2\pi }}\left(C+{\displaystyle \frac{1}{2}}\mathrm{log}\left(\pi +k_1\right)+{\displaystyle \frac{1}{2}}\mathrm{log}\left(\pi k_1\right)\right)e^{ik_1x}`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑l\chi _{l[\pi ,\pi ]}\chi _{l[\pi +k,\pi +k]}\left(C+{\displaystyle \frac{1}{2}}\mathrm{log}\left(\pi +lk\right)+{\displaystyle \frac{1}{2}}\mathrm{log}\left(\pi \left(lk\right)\right)\right)`$ where to obtain the equality use has been made of (A.2). Evaluating the integral gives $$\mathrm{FT}\frac{\mathrm{sin}\pi x}{\pi x}A_3\left(x\right)=\left(C+\mathrm{log}2\pi \right)\left(1\frac{\left\|k\right\|}{2\pi }\right)\frac{\left\|k\right\|}{4\pi }\mathrm{log}\frac{\left\|k\right\|}{2\pi }+\frac{1}{2}\left(1\frac{\left\|k\right\|}{2\pi }\right)\mathrm{log}\left(1\frac{\left\|k\right\|}{2\pi }\right)$$ (A.21) for $`\left\|k\right\|<2\pi `$, while for $`\left\|k\right\|>2\pi `$ $$\mathrm{FT}\frac{\mathrm{sin}\pi x}{\pi x}A_3\left(x\right)=0.$$ (A.22) Substituting the above results as appropriate in the Fourier transform of (Appendix A) gives the result (7.9). ## Appendix B In this appendix we outline some details of the calculation of (7.1) in the case $`\beta _0=4`$ and show how this leads to (7.13). Because (7.2) and (7.3) formally have the same structure upon expansion (recall the definition of qdet (7.6)), the formula (Appendix A) formally maintains its structure when generalized to the case $`\beta _0=4`$. Thus we have $`g_2(x_1,x_2;\beta )`$ (B.1) $`=1\left(P_4(x_1,x_2)P_4(x_2,x_1)\right)^{\left(0\right)}+(\beta 4)\{(1\left(P_4(x_1,x_2)P_4(x_2,x_1)\right)^{\left(0\right)})\mathrm{\Phi }(x_1,x_2)`$ $`2{\displaystyle _{\mathrm{}}^{\mathrm{}}}(\left(P_4(x_2,x_3)P_4(x_3,x_2)\right)^{\left(0\right)}\left(P_4(x_1,x_3)P_4(x_3,x_1)\right)^{\left(0\right)}`$ $`+2\left(P_4(x_1,x_2)P_4(x_2,x_3)P_4(x_3,x_1)\right)^{\left(0\right)}\left)\mathrm{\Phi }\right(x_1,x_3)dx_3`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}(4\left(P_4(x_1,x_3)P_4(x_3,x_4)P_4(x_4,x_1)\right)^{\left(0\right)}4\left(P_4(x_1,x_2)P_4(x_2,x_3)P_4(x_3,x_4)P_4(x_4,x_1)\right)^{\left(0\right)}`$ $`2\left(P_4(x_1,x_3)P_4(x_3,x_2)P_4(x_2,x_4)P_4(x_4,x_1)\right)^{\left(0\right)}`$ $`+2\left(P_4(x_1,x_3)P_4(x_3,x_1)\right)^{\left(0\right)}\left(P_4(x_2,x_4)P_4(x_4,x_2)\right)^{\left(0\right)})\mathrm{\Phi }(x_3,x_4)dx_3dx_4\}+O((\beta 4)^2).`$ We treat each of the seven distinct integrals in (B.1) in an analogous way to their counterparts in (Appendix A), although extra working is involved due to $`P_4`$ being a matrix rather than a scalar. The final results are $`B_1\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_4(x_2,x_3)P_4(x_3,x_2)\right)^{\left(0\right)}\mathrm{\Phi }(x_1,x_3)𝑑x_3`$ $`=`$ $`{\displaystyle _{4\pi }^{4\pi }}{\displaystyle \frac{dk}{2\pi }}\left(1{\displaystyle \frac{\left\|k\right\|}{4\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right){\displaystyle \frac{\pi }{\left\|k\right\|}}\mathrm{cos}kx_{12}`$ $`=`$ $`\mathrm{log}\left\|x_{12}\right\|{\displaystyle \frac{\mathrm{sin}4\pi x_{12}}{4\pi x_{12}}}+\mathrm{ci}\left(4\pi x_{12}\right){\displaystyle \frac{\mathrm{cos}2\pi x_{12}}{4\pi x_{12}}}\mathrm{Si}\left(2\pi x_{12}\right)`$ $`B_2`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_4(x_1,x_3)P_4(x_3,x_1)\right)^{\left(0\right)}\mathrm{\Phi }(x_1,x_3)𝑑x_3=B_1\left(0\right)=C+\mathrm{log}4\pi {\displaystyle \frac{3}{2}}`$ $`B_3\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_4(x_2,x_3)P_4(x_3,x_1)\mathrm{\Phi }(x_1,x_3)𝑑x_3`$ (B.3) $`=`$ $`\left[\begin{array}{cc}\frac{1}{4}f_1\left(x_{12}\right)\frac{1}{4}f_3\left(x_{12}\right)& \frac{1}{4}_0^{x_{12}}\left(f_1\left(t\right)+f_2\left(t\right)\right)𝑑t\\ \frac{1}{4}f_1^{}\left(x_{12}\right)+\frac{1}{4}f_3^{}\left(x_{12}\right)& \frac{1}{4}f_1\left(x_{12}\right)+\frac{1}{4}f_2\left(x_{12}\right)\end{array}\right]`$ $`B_4`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_4(x_1,x_3)P_4(x_3,x_4)P_4(x_4,x_1)\right)^{\left(0\right)}\mathrm{\Phi }(x_1,x_3)𝑑x_3𝑑x_4=B_1\left(0\right)`$ $`B_5\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_4(x_2,x_3)P_4(x_3,x_4)P_4(x_4,x_1)\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4`$ (B.5) $`=`$ $`\left[\begin{array}{cc}\frac{1}{4}f_1\left(x_{12}\right)+\frac{1}{8}f_2\left(x_{12}\right)\frac{1}{8}f_3\left(x_{12}\right)& _0^{x_{12}}\left(\frac{1}{4}f_1\left(t\right)+\frac{1}{8}f_2\left(t\right)\frac{1}{8}f_3\left(t\right)\right)𝑑t\\ \left(\frac{1}{4}f_1^{}\left(x_{12}\right)+\frac{1}{8}f_2^{}\left(x_{12}\right)\frac{1}{8}f_3^{}\left(x_{12}\right)\right)& \frac{1}{4}f_1\left(x_{12}\right)+\frac{1}{8}f_2\left(x_{12}\right)\frac{1}{8}f_3\left(x_{12}\right)\end{array}\right]`$ $`B_6\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_4(x_1,x_3)P_4(x_3,x_2)P_4(x_2,x_4)P_4(x_4,x_3)\right)^{\left(0\right)}\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4`$ $`=`$ $`{\displaystyle _{4\pi }^{4\pi }}{\displaystyle \frac{1}{2\left\|k\right\|}}\{\left(g_1(k,x_{12})\right)^2+\mathrm{cos}(kx_{12}/2)g_1(k,x_{12})({\displaystyle \frac{\left\|k\right\|}{4\pi }}g_2(k,x_{12}){\displaystyle \frac{ik}{4\pi }}g_3(k,x_{12}))`$ $`+\mathrm{cos}\left(kx_{12}\right)({\displaystyle \frac{\left\|k\right\|}{8\pi }}g_2(k,x_{12}){\displaystyle \frac{ik}{8\pi }}g_3(k,x_{12}))^2({\displaystyle \frac{\mathrm{sin}\left(\left\|k\right\|x_{12}/2\right)}{4\pi }}g_2(k,x_{12})`$ $`+{\displaystyle \frac{\mathrm{cos}\left(kx_{12}/2\right)}{4\pi }}g_3(k,x_{12})\left)\right({\displaystyle \frac{(4\pi \|k\left\|\right)\mathrm{cos}\left((2\pi \|k\|/2)\right)x_{12})}{4\pi x_{12}}}`$ $`{\displaystyle \frac{\mathrm{sin}\left((2\pi \|k\|/2)\right)x_{12})}{2\pi x_{12}^2}})\}dk`$ $`B_7\left(x_{12}\right)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}\left(P_4(x_1,x_3)P_4(x_3,x_1)\right)^{\left(0\right)}\left(P_4(x_2,x_4)P_4(x_4,x_2)\right)^{\left(0\right)}\mathrm{\Phi }(x_3,x_4)𝑑x_3𝑑x_4`$ (B.6) $`=`$ $`{\displaystyle _{4\pi }^{4\pi }}\left(1{\displaystyle \frac{\left\|k\right\|}{4\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{ln}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right)^2{\displaystyle \frac{\pi }{\left\|k\right\|}}\mathrm{cos}kx_{12}{\displaystyle \frac{dk}{2\pi }}`$ $`=`$ $`\mathrm{log}\left\|x\right\|{\displaystyle \frac{\mathrm{sin}4\pi x}{4\pi x}}+\mathrm{ci}\left(4\pi x\right){\displaystyle \frac{\mathrm{Si}\left(2\pi x\right)\mathrm{sin}2\pi x}{8\pi ^2x^2}}{\displaystyle \frac{\mathrm{Si}\left(2\pi x\right)\mathrm{cos}2\pi x}{4\pi x}}`$ $`+{\displaystyle _{4\pi }^{4\pi }}\left\|k\right\|\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|^2\mathrm{cos}kx{\displaystyle \frac{dk}{128\pi ^2}}`$ where $`f_1\left(x\right)`$ $`:=`$ $`{\displaystyle _{2\pi }^{2\pi }}\left(2C+\mathrm{ln}\left(4\pi ^2k^2\right)\right)\mathrm{cos}kx{\displaystyle \frac{dk}{2\pi }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}2\pi x}{\pi x}}\left(C+\mathrm{ln}4\pi \mathrm{ln}\left\|x\right\|+\mathrm{ci}\left(4\pi x\right)\right){\displaystyle \frac{\mathrm{cos}2\pi x}{\pi x}}\mathrm{Si}\left(4\pi x\right)`$ $`f_2\left(x\right)`$ $`:=`$ $`{\displaystyle \frac{1}{\pi x}}\mathrm{Si}\left(2\pi x\right),f_3\left(x\right):={\displaystyle \frac{\mathrm{sin}2\pi x}{\pi x}}`$ $`g_1(k,x)`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}\left(\left(2\pi \left\|k\right\|/2\right)x\right)}{2\pi x}}`$ $`g_2(k,x)`$ $`=`$ $`\mathrm{ci}\left(\left(2\pi \left\|k\right\|\right)x\right)\mathrm{ci}\left(2\pi x\right)`$ $`g_3(k,x)`$ $`=`$ $`\mathrm{Si}\left(\left(2\pi \left\|k\right\|\right)x\right)+\mathrm{Si}\left(2\pi x\right)`$ (B.7) When substituted in (B.1), the constant terms $`B_2`$ and $`B_4`$ cancel, and we obtain the formula $`g(x_1,x_2;\beta )`$ $`=`$ $`1\left(P_4(x_1,x_2)P_4(x_2,x_1)\right)^{\left(0\right)}`$ (B.8) $`+(\beta 4)\{(1+\left(P_4(x_1,x_2)P_4(x_2,x_1)\right)^{\left(0\right)}\mathrm{\Phi }(x_1,x_2)`$ $`+2B_1\left(x_{12}\right)4\left(P_4\left(x_{12}\right)B_3\left(x_{12}\right)\right)^{\left(0\right)}+2\left(P_4\left(x_{12}\right)B_5\left(x_{12}\right)\right)^{\left(0\right)}`$ $`+B_6\left(x_{12}\right)B_7\left(x_{12}\right)\}+O((\beta 4)^2).`$ We will demonstrate the close analogy with the $`\beta =2`$ calculation of Appendix A by giving the derivation of the integral formula in (Appendix B) for $`B_1\left(x_{12}\right)`$. As in the derivation of the integral formula (A.3) for $`A_1\left(x_{12}\right)`$, our strategy is to use the convolution formula (A.2). However here the Fourier transform of $`P_4(x_1,x_2)P_4(x_2,x_1)`$ is not immediate. What is immediate is the Fourier transform of $`P_4(x_1,x_2)`$. Thus from the definition (7.4) we see that $`\mathrm{FT}P_4(x_1,x_2)`$ $`:=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}P_4(x_1,x_2)e^{ikx_{12}}𝑑x_{12}`$ $`=`$ $`\{\begin{array}{cc}\left[\begin{array}{cc}1/2& i/2k\\ ik/2& 1/2\end{array}\right],\hfill & \left\|k\right\|<2\pi \hfill \\ \left[\begin{array}{cc}0& 0\\ 0& 0\end{array}\right],\hfill & \left\|k\right\|>2\pi .\hfill \end{array}`$ Use of (A.2) then shows that for $`\left\|k\right\|<4\pi `$ $`\mathrm{FT}P_4(x_1,x_2)P_4(x_2,x_1)`$ $`:=`$ $`{\displaystyle _{2\pi }^{2\pi }}\left[\begin{array}{cc}1/2& i/2l\\ il/2& 1/2\end{array}\right]\left[\begin{array}{cc}1/2& i/2\left(kl\right)\\ i\left(kl\right)/2& 1/2\end{array}\right]\chi _{\left\|kl\right\|<2\pi }{\displaystyle \frac{dl}{2\pi }}`$ (B.16) $`=`$ $`\left[\begin{array}{cc}1\frac{\left\|k\right\|}{4\pi }+\frac{\left\|k\right\|}{8\pi }\mathrm{log}\left\|1\frac{\left\|k\right\|}{2\pi }\right\|& 0\\ 0& 1\frac{\left\|k\right\|}{4\pi }+\frac{\left\|k\right\|}{8\pi }\mathrm{log}\left\|1\frac{\left\|k\right\|}{2\pi }\right\|\end{array}\right],`$ (B.19) while for $`\left\|k\right\|>4\pi `$ $$\mathrm{FT}P_4(x_1,x_2)P_4(x_2,x_1)=0.$$ (B.20) The results (B.16) and (B.20) are the analogue of (Appendix A) in the working leading to the evaluation of $`A_1\left(x_{12}\right)`$. The integral formula for $`B_1\left(x_{12}\right)`$ in (Appendix B) now follows from (B.16), (B.20) and (5.18) upon a further application of (A.2). The Fourier transform of (B.8) can be computed explicitly. The final result has already been stated in (7.13). This is obtained through the intermediate results $$\mathrm{FT}B_j\left(x\right)=\widehat{B}_j\left(k\right)\mathrm{for}j=0,1,3,5,6,7$$ with $$B_0\left(x_{12}\right):=\left(P_4(x_1,x_2)P_4(x_2,x_1)\right)^{\left(0\right)}\mathrm{\Phi }(x_1,x_2)$$ and the $`\widehat{B}_j`$ specified by (7.14)–(7.20). We will illustrate the working by giving some details of the computation of $`\widehat{B}_0\left(k\right)`$ for $`\left\|k\right\|<4\pi `$. Using (B.16) and (5.18) we see from (A.2) that $`\mathrm{FT}B_0\left(x_{12}\right)`$ $`=`$ $`{\displaystyle _{4\pi }^{4\pi }}\left(1{\displaystyle \frac{\left\|l\right\|}{4\pi }}+{\displaystyle \frac{\left\|l\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|l\right\|}{2\pi }}\right\|\right){\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}`$ (B.21) $`=`$ $`\left(1{\displaystyle \frac{\left\|k\right\|}{4\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\right){\displaystyle _{4\pi }^{4\pi }}{\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}`$ $`+{\displaystyle _{4\pi }^{4\pi }}\left({\displaystyle \frac{\left\|k\right\|\left\|l\right\|}{4\pi }}\right){\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}+{\displaystyle _{4\pi }^{8\pi }}\left({\displaystyle \frac{\left\|l\right\|\left\|k\right\|}{8\pi }}\right)\mathrm{log}\left\|1{\displaystyle \frac{\left\|l\right\|}{2\pi }}\right\|{\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}`$ $`+{\displaystyle \frac{\left\|k\right\|}{8\pi }}{\displaystyle _{4\pi }^{4\pi }}\mathrm{log}\left\|{\displaystyle \frac{2\pi \left\|l\right\|}{2\pi \left\|k\right\|}}\right\|{\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}`$ where the second equality, which follows from minor manipulation of the first integral, is motivated by the desire to separate the singular integral. Thus in the second equality of (B.21) only the first integral is singular. It is essentially the same as the first singular integral in (A.13), and is evaluated as $$_{4\pi }^{4\pi }\frac{\pi }{\left\|kl\right\|}\frac{dl}{2\pi }=C+\frac{1}{2}\mathrm{log}\left(16\pi ^2k^2\right),\left\|k\right\|4\pi .$$ (B.22) The second integral in the second equality of (B.21) also appears in the evaluation of (A.13). An elementary calculation shows $$_{4\pi }^{4\pi }\frac{\left\|k\right\|\left\|l\right\|}{4\pi }\frac{\pi }{\left\|kl\right\|}\frac{dl}{2\pi }=1+\frac{\left\|k\right\|}{4\pi }+\frac{\left\|k\right\|}{4\pi }\mathrm{log}\left(\frac{4\pi +\left\|k\right\|}{\left\|k\right\|}\right).$$ (B.23) To evaluate the third integral in (B.21) we suppose without loss of generality that $`k>0`$ and write $`{\displaystyle _{4\pi }^{4\pi }}{\displaystyle \frac{\left\|l\right\|\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|l\right\|}{2\pi }}\right\|{\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}={\displaystyle _0^{2\pi }}\left({\displaystyle \frac{1}{16\pi }}+{\displaystyle \frac{k}{8\pi \left(lk\right)}}\right)\mathrm{log}\left\|1+{\displaystyle \frac{l}{2\pi }}\right\|dl`$ (B.24) $`{\displaystyle _0^k}\mathrm{log}\left\|1{\displaystyle \frac{\left\|l\right\|}{2\pi }}\right\|{\displaystyle \frac{dl}{16\pi }}+{\displaystyle _k^{4\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|l\right\|}{2\pi }}\right\|{\displaystyle \frac{dl}{16\pi }}.`$ The only non-elementary integral is the second term of the first integral. This can be computed by checking from the definition (7.12) that for $`4\pi <l<0`$ $$\frac{d}{dl}\left(\mathrm{dilog}\left(\frac{kl}{k+2\pi }\right)+\mathrm{log}\left\|1+\frac{l}{2\pi }\right\|\mathrm{log}\left(\frac{kl}{k+2\pi }\right)\right)=\frac{1}{lk}\mathrm{log}\left\|1+\frac{l}{2\pi }\right\|.$$ (B.25) In total we therefore have $`{\displaystyle _{4\pi }^{4\pi }}{\displaystyle \frac{\left\|l\right\|\left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|l\right\|}{2\pi }}\right\|{\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}`$ (B.26) $`={\displaystyle \frac{1}{2}}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}+{\displaystyle \frac{\left\|k\right\|}{8\pi }}\left(\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)\mathrm{dilog}\left({\displaystyle \frac{4\pi +\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)\right)+{\displaystyle \frac{2\pi \left\|k\right\|}{8\pi }}\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|.`$ To evaluate the final integral in (B.21), a similar approach to that leading to the evaluation (B.26) is adopted. Minor complications arise because of the need to modify the formula (B.25) for $`l>k`$. We find $`{\displaystyle _{4\pi }^{4\pi }}{\displaystyle \frac{\left\|k\right\|}{8\pi }}\mathrm{log}\left\|{\displaystyle \frac{2\pi \left\|l\right\|}{2\pi \left\|k\right\|}}\right\|{\displaystyle \frac{\pi }{\left\|kl\right\|}}{\displaystyle \frac{dl}{2\pi }}`$ $`={\displaystyle \frac{\left\|k\right\|}{16\pi }}\left\{\mathrm{dilog}\left({\displaystyle \frac{4\pi +\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)\mathrm{dilog}\left({\displaystyle \frac{\left\|k\right\|}{2\pi +\left\|k\right\|}}\right)\mathrm{log}\left\|1{\displaystyle \frac{\left\|k\right\|}{2\pi }}\right\|\mathrm{log}\left({\displaystyle \frac{4\pi +\left\|k\right\|}{\left\|k\right\|}}\right)+g_1\left(k\right)\right\}`$ where $`g_1`$ is defined in (7.23). Substituting (B.22)–(Appendix B) in (B.21) gives the result (7.14).
warning/0002/hep-ph0002223.html	ar5iv	text	# 1 Introduction ## 1 Introduction In 1995 the KARMEN collaboration discovered an anomaly in the time distribution of the charged and neutral current events induced by neutrinos from $`\pi ^+`$ and $`\mu ^+`$ decays at rest . This anomaly may be explained by the production of a new neutral particle in pion decay $`\pi ^+\mu ^++x^0,`$ (1) with the mass 33.9 MeV, barely permitted by the phase space, so that this particle moves with non-relativistic velocity. Its subsequent neutrino-producing decays could be the source of the delayed neutrinos observed in the experiment. The anomaly was recently confirmed by the same group with better statistics and substantially reduced cosmic-ray background. A search for $`x^0`$ particles produced in the rare pion decay (1) was performed in 1995 by an experiment at PSI . It gave an upper limit for the branching ratio of $`\mathrm{BR}(\pi ^+\mu ^++x^0)<2.6\times 10^8`$ at 95% CL. Several candidates for $`x^0`$ have been proposed in the literature. In ref. the authors considered a sterile neutrino, $`x^0\nu _s`$. Their conclusion was that the sterile-neutrino hypothesis is compatible with all laboratory constraints, but they noted possible problems with astrophysics and cosmology. Further laboratory constraints on this model were investigated in ref. , concluding that mixings of $`\nu _s`$ with $`\nu _e`$ and $`\nu _\mu `$ must be very small, while a mixing with $`\nu _\tau `$ was permitted. In this case $`\nu _s`$ would predominantly decay through neutral-current interactions into $`\nu _\tau +\mathrm{}+\overline{\mathrm{}}`$, where $`\mathrm{}`$ is any light lepton, $`\mathrm{}=\nu _e`$, $`\nu _\mu `$, $`\nu _\tau `$, or $`e^{}`$. The lifetime of $`\nu _s`$ with respect to this decay was estimated to be in the range $`10^3\mathrm{sec}<\tau _{\nu _s}<150\mathrm{sec}.`$ (2) The lower bound on $`\tau _{\nu _s}`$, which comes from the experimental bound on the $`\nu _\tau `$ mass (large mixing with $`\nu _s`$ makes it too heavy), was discussed in ref. . In ref. it was suggested that $`x^0\stackrel{~}{\chi }`$ could be the lightest supersymmetric particle, photino or zino, that decayed through the channel $`\stackrel{~}{\chi }\gamma +\nu _\mu .`$ (3) The new data , however, do not agree with the predictions of the model so that recently a new version of supersymmetric model was considered , according to which the light neutralino decayed through a three body channel $`\stackrel{~}{\chi }e^++e^{}+\nu _{\mu ,\tau }.`$ (4) In the paper a new decay mode of muons was proposed as a source of the anomaly: $`\mu ^+e^++S,`$ (5) where $`S`$ is a scalar boson with mass 103.9 MeV. A search for these decays was reported in ref. where the upper limit $`\mathrm{BR}(\mu ^+e^++S)<5.7\times 10^4`$ was obtained. This limit, though it does not exclude the model, makes it more complicated. Recent data by the NOMAD collaboration permit to strengthen the bound on the mixing of $`\nu _s`$ with $`\nu _\tau `$. Their expected lower limit on $`\tau _{\nu _s}`$ is around 0.1 sec. On the other hand, cosmology and astrophysics permit to obtain an upper bound on $`\tau _{\nu _s}`$ that may be complementary to direct experiments. If this happens to be the case, then the explanation of the KARMEN anomaly by a 33.9 MeV sterile neutrino would be ruled out. Recently two papers have appeared where the observation of SN 1987A were used to put bounds on the properties of the proposed sterile neutrino or light neutralino. In the present paper we derive constraints on the lifetime of a 33.9 MeV sterile neutrino from both more detailed consideration of SN 1987A and of primordial nucleosynthesis. With the present-day accuracy of the data on the primordial light-element abundances, the bounds that are found from SN 1987A tend to be stronger than those found from primordial nucleosynthesis. However, the latter remain interesting, first, because they may be competitive in the nearest future with an improved accuracy of the data on primordial nucleosynthesis, and second, in a hypothetical case where $`\nu _s`$ would have an anomalously strong (stronger than the usual weak) interaction with nucleons. In that case the SN 1987A bounds would not apply, while the nucleosynthesis bounds would survive. ## 2 Primordial Nucleosynthesis ### 2.1 General Features A heavy unstable sterile neutrino would influence big-bang nucleosynthesis (BBN) through its contribution to the cosmological energy density by speeding up the expansion and enlarging the frozen neutron-to-proton ratio, $`r_n=n/p`$, and less directly, though stronger, through its decay products, $`\nu _e,\nu _\mu `$, and $`\nu _\tau `$. The impact of $`\nu _\mu `$ and $`\nu _\tau `$ on BBN is rather straightforward: their energy density increases with respect to the standard case and this also results in an increase of $`r_n`$. This effect can be described by the increased number of effective neutrino species $`N_\nu `$ during BBN (in the standard case $`N_\nu =3`$). The increase of the energy density of $`\nu _e`$, due to decay of $`\nu _s`$ into $`\nu _e`$, has an opposite effect on $`r_n`$. Though a larger energy density results in faster cooling, the increased number of $`\nu _e`$ would preserve thermal equilibrium between neutrons and protons for a longer time and correspondingly the frozen $`n/p`$-ratio would become smaller. The second effect is stronger, so the net result is a smaller $`n/p`$-ratio. There is, however, another effect which is related to the distortion of the energy spectrum of $`\nu _e`$ from the decays of $`\nu _s`$. If the spectrum is distorted at the high-energy tail, as is the case, then proton formation in the reaction $`n+\nu _ep+e^{}`$ would be less efficient than neutron creation in the reaction $`\overline{\nu }_e+pn+e^+`$. We found that this effect is quite significant. The decays of $`\nu _s`$ into the $`e^+e^{}`$-channel will inject more energy into the electromagnetic part of the primeval plasma and this will diminish the relative contribution of the energy density of light neutrinos and diminish $`r_n`$. We could use the technique and numerical code from our earlier papers where the effects of non-equilibrium massless neutrinos (in the standard model) and possibly massive and/or unstable tau-neutrinos were precisely calculated. However, since we do not need the accuracy of a fraction of per cent achieved in these papers, we will use a simpler and considerably less time consuming approximate approach that we will now discuss. ### 2.2 The model If the KARMEN anomaly is explained by a heavy sterile neutrino with $`m_{\nu _s}=33.9`$ MeV then, as was mentioned in the Introduction, it may only mix with $`\nu _\tau `$, $`\nu _\tau `$ $`=`$ $`\nu _1\mathrm{cos}\mathrm{\Theta }+\nu _2\mathrm{sin}\mathrm{\Theta },`$ $`\nu _s`$ $`=`$ $`\nu _1\mathrm{sin}\mathrm{\Theta }+\nu _2\mathrm{cos}\mathrm{\Theta },`$ (6) where $`\mathrm{\Theta }`$ is the vacuum mixing angle and $`\nu _1`$ and $`\nu _2`$ are the mass eigenstates; the mass difference is positive: $`\delta m^2=m_2^2m_1^2(33.9\mathrm{MeV})^2>0`$. Through neutral-current interactions $`\nu _s`$ could decay into $`\nu _\tau `$ and a pair of other light leptons. The corresponding processes and their matrix elements are presented in table 1. The lifetime of $`\nu _s`$ is given by the expression $`\tau _{\nu _s}\mathrm{\Gamma }_{\nu _s}^1=\left[{\displaystyle \frac{(1+\stackrel{~}{g}_L^2+g_R^2)G_F^2m_{\nu _s}^5(\mathrm{sin}^22\mathrm{\Theta })/4}{192\pi ^3}}\right]^1{\displaystyle \frac{5.7\times 10^4\mathrm{sec}}{(\mathrm{sin}^22\mathrm{\Theta })/4}},`$ (7) where $`\stackrel{~}{g}_L=1/2+\mathrm{sin}^2\theta _W\mathrm{and}g_R=\mathrm{sin}^2\theta _W.`$ (8) According to the combined experimental data, the lifetime lies in the range $`0.1\mathrm{sec}<\tau _{\nu _s}<150\mathrm{sec}`$ (9) and correspondingly $`3.9\times 10^3<\mathrm{sin}2\mathrm{\Theta }<0.15.`$ (10) Even with a very small mixing, $`\nu _s`$ could be abundantly produced in the early universe when the temperature was higher than $`m_{\nu _s}`$. Their production rate can be estimated as $$\mathrm{\Gamma }_{\nu _s}=\frac{1}{2}\mathrm{sin}^22\mathrm{\Theta }_\mathrm{M}\mathrm{\Gamma }_W,$$ (11) where $`\mathrm{\Gamma }_W=2.5G_F^2T^5`$ is the averaged weak interaction rate and $`\mathrm{\Theta }_\mathrm{M}`$ is the mixing angle in the medium. According to the calculations of refs. $`\mathrm{sin}2\mathrm{\Theta }_\mathrm{M}{\displaystyle \frac{(\mathrm{sin}2\mathrm{\Theta })/2}{1+0.76\times 10^{19}T^6(\delta m^2)^1}}{\displaystyle \frac{(\mathrm{sin}2\mathrm{\Theta })/2}{1+6.6\times 10^{23}T^6}},`$ (12) where $`T`$ and $`\delta m^2`$ are taken in MeV. One sees that matter effects are not important for $`T<5`$ GeV. Comparing the production rate (11) with the Hubble expansion rate $$H=\sqrt{\frac{8\pi ^3}{90}g_{}(T)}\frac{T^2}{M_{\mathrm{Pl}}},$$ (13) we find that sterile neutrinos could be abundantly produced in the early universe if, roughly speaking, $`(\mathrm{sin}2\mathrm{\Theta })^2(T/3\mathrm{MeV})^3>1`$. Even for $`\mathrm{\Theta }`$ at the lower limit given by the relation (10), the equilibrium condition is fulfilled for $`T120`$ MeV. For this small value of $`\mathrm{\Theta }`$ the original equilibrium number density of $`\nu _s`$ would be diluted at smaller $`T`$ by annihilation of pions and muons. However, as we see in what follows, nucleosynthesis strongly disfavors large values of $`\tau _{\nu _s}`$. Correspondingly, for large mixing angles $`\nu _s`$ remains in equilibrium for much smaller $`T`$, and this dilution is not essential. The evolution of $`\nu _s`$ at lower temperatures, $`T<m_{\nu _s}`$, is considered below. We calculate the number density of the heavy $`\nu _s`$ assuming that initially they were in equilibrium and in the process of freeze-out they interacted with the thermal equilibrium bath of light particles. With the evident correction for the $`\nu _s`$ decay, these calculations are very similar to the usual freeze-out calculation of massive species. There is one important difference however. Normally massive particles disappear in the process of mutual annihilation, so that the rate of freezing is proportional to the number density of the particle in question, $`\dot{n}_m/n_m\sigma _{\mathrm{ann}}n_m`$, which becomes exponentially small when $`T<m`$. Sterile neutrinos may also disappear in collisions with massless leptons $`\nu _s+\mathrm{}_1\mathrm{}_2+\mathrm{}_3,`$ (14) so their extinction through this process would be more efficient than by mutual annihilation, $`\nu _s+\overline{\nu }_s\mathrm{all}`$. The decoupling temperature of the process (14) can be approximately found from the decoupling of the usual massless neutrinos at $`T_{\nu _\tau }2`$ MeV by rescaling it by the mixing parameter, $`(\mathrm{sin}2\mathrm{\Theta })^{2/3}`$. To obtain a better estimate we solve the corresponding kinetic equation for $`\nu _s`$ in sec. 2.3. This permits us to determine the distribution function, $`f_{\nu _s}(x,y)`$, where $`x=m_0a(t)\mathrm{and}y=pa(t),`$ (15) are convenient dimensionless variables with $`a(t)`$ the cosmic scale factor and $`m_0`$ the normalization mass that we choose as $`m_0=1`$ MeV. In what follows we will often skip $`m_0`$, keeping in mind that the relevant quantities are measured in MeV. In terms of these variables, the kinetic equations have the form $`(_tHp_p)f=Hx_xf=I_{\mathrm{coll}},`$ (16) where $`I_{\mathrm{coll}}`$ $`=`$ $`{\displaystyle \frac{1}{2E}}{\displaystyle \underset{i}{}\left(\frac{d^3p_i}{2E_i(2\pi )^3}\right)\underset{f}{}\left(\frac{d^3p_i}{2E_f(2\pi )^3}\right)}`$ (17) $`\times (2\pi )^4\delta ^{(4)}({\displaystyle \underset{i}{}}p_i{\displaystyle \underset{f}{}}p_f)\left\|A_{if}\right\|^2F(f_i,f_f),`$ is the collision integral and $`F(f_i,f_f)={\displaystyle \underset{i}{}}f_i{\displaystyle \underset{f}{}}(1f_f)+{\displaystyle \underset{f}{}}f_f{\displaystyle \underset{i}{}}(1f_i)`$ (18) with sub-$`i`$ and sub-$`f`$ meaning initial and final particles. After the distribution function of $`\nu _s`$ is found, the next step is to find the distribution functions of the light neutrinos and in particular their energy densities. The distributions of electrons and positrons are of course assumed to be very close to equilibrium because of their very fast thermalization due to the interaction with the photon bath. However, the evolution of the photon temperature, due to decay and annihilation of the massive $`\nu _s`$, becomes different from the standard one, $`T_\gamma 1/x`$, by an extra factor $`(1+\mathrm{\Delta })>1`$: $`T_\gamma =[1+\mathrm{\Delta }(x)]/x.`$ (19) At sufficiently high temperatures, $`T>T_W2`$ MeV, light neutrinos and electrons/positrons were in strong contact, so that the neutrino distributions were also very close to the equilibrium ones. If $`\nu _s`$ disappeared sufficiently early, while thermal equilibrium between $`e^\pm `$ and neutrinos remained, then $`\nu _s`$ would not have any observable effect on primordial abundances, because only the contribution of neutrino energy density relative to the energy density of $`e^\pm `$ and $`\gamma `$ is essential for nucleosynthesis. Hence a very short-lived $`\nu _s`$ has a negligible impact on primordial abundances, while with an increasing lifetime the effect becomes stronger. Indeed at $`T<T_W`$ the exchange of energy between neutrinos and electrons becomes very weak and the energy injected into the neutrino component is not immediately redistributed between all the particles. The branching ratio of the decay of $`\nu _s`$ into $`e^+e^{}`$ is approximately 1/9, so that the neutrino component is heated much more than the electromagnetic one. As we mentioned above, this leads to a faster cooling and to a larger $`n/p`$-ratio. So, for the numerical calculations we adopt the following procedure. We assume that at sufficiently high temperature, e.g. $`T_i=5`$ MeV, or equivalently at the initial value of the scale factor $`x_i=1/T_i=0.2`$, there is a complete thermal equilibrium between active neutrinos and electrons. Starting from that moment we calculate the corrections to the active neutrino distribution functions $`\delta f_\nu =f_\nu f_\nu ^{\mathrm{eq}},`$ (20) where the equilibrium distribution function is assumed to have the standard Fermi-Dirac form with a temperature that drops as $`1/x`$, $`f_\nu ^{\mathrm{eq}}=(e^y+1)^1.`$ (21) The evolution of the photon temperature, i.e. $`\mathrm{\Delta }(x)`$, is determined from the energy balance equation $`d\rho /dx=3(\rho +p)/x,`$ (22) where $`\rho `$ and $`p`$ are respectively the total energy and pressure densities in the cosmological plasma. It is worth noting that we normalized the scale factor in such a way that at the initial moment $`x_iT_i=1`$. If we change the initial moment for calculating $`\mathrm{\Delta }`$, it would result in a different definition of $`x`$, but this is unobservable. We have checked that for $`x_i=0.1`$–0.3 the results weakly depend upon the value of $`x_i`$, and that the corrected neutrino distribution $`f_\nu =f_\nu ^{\mathrm{eq}}+\delta f`$ quite accurately maintain the equilibrium shape with the same temperature as electrons, $`T_\nu =T_\gamma `$ at these early times. We made several assumptions that permitted to simplify calculations very much: $`\nu _s`$ was assumed non-relativistic, the equilibrium distribution functions “inside” the collision integral were taken in the Boltzmann approximation, while “outside” they were taken in Fermi-Dirac form, and the kinetic equations for $`\delta f_\nu `$ were taken in a simplified form, so that some fraction of neutrino energy was lost (see sec. 2.4). All these assumptions lead to a weaker impact of $`\nu _s`$ on nucleosynthesis, so the real bound should be somewhat stronger than what is presented below. One more comment is in order. We do not take into account oscillations between $`\nu _s`$ and $`\nu _\tau `$ for $`T<m_{\nu _s}`$. This is perfectly justified because in the interesting range of neutrino energies the oscillation frequency is so high and the velocities of $`\nu _s`$ and $`\nu _\tau `$ are so much different that the coherence is quickly lost and they can be considered as independent particles. Medium effects are also not important for the considered positive mass difference $`\delta m^210^3\mathrm{MeV}^2`$. ### 2.3 Evolution of Heavy Neutrinos The evolution of the occupation number of $`\nu _s`$ is determined by its decays and inverse decays, listed in table 1, and by the reactions (14) with all possible sets of light leptons permitted by quantum numbers, presented in table 2. Taking both contributions into the collision integral and assuming that the massless species are in thermal equilibrium with temperature $`T`$ and that both helicity states of $`\nu _s`$ are equally populated, we obtain $`_xf_{\nu _s}(x,y)`$ $`=`$ $`{\displaystyle \frac{1.48x}{\tau _{\nu _s}(\mathrm{sec})}}\left({\displaystyle \frac{10.75}{g_{}(T)}}\right)^{1/2}{\displaystyle \frac{f_{\nu _s}^{\mathrm{eq}}f_{\nu _s}}{(Tx)^2}}`$ (23) $`\left[{\displaystyle \frac{m_{\nu _s}}{E_{\nu _s}}}+{\displaystyle \frac{3\times 2^7T^3}{m_{\nu _s}^3}}\left({\displaystyle \frac{3\zeta (3)}{4}}+{\displaystyle \frac{7\pi ^4}{144}}\left({\displaystyle \frac{E_{\nu _s}T}{m_{\nu _s}^2}}+{\displaystyle \frac{p_{\nu _s}^2T}{3E_{\nu _s}m_{\nu _s}^2}}\right)\right)\right],`$ where $`E_{\nu _s}=\sqrt{m_{\nu _s}^2+(y/x)^2}`$ and $`p_{\nu _s}=y/x`$ are the energy and momentum of $`\nu _s`$ respectively and $`g_{}(T)`$ is the effective number of massless species in the plasma determined as the ratio of the total energy density to the equilibrium energy density of one bosonic species with temperature $`T`$, $`g_{}=\rho _{\mathrm{tot}}/(\pi ^2T^4/30)`$. The coefficient $`1.48`$ comes from expressing the lifetime in sec according to the relation $`1/\mathrm{MeV}=0.658\times 10^{21}`$ sec and from the value of the Hubble parameter $`H=\sqrt{8\pi \rho _{\mathrm{tot}}/3M_{\mathrm{Pl}}^2}`$. The first term in square brackets, $`m_{\nu _s}/E_{\nu _s}`$, comes from summing the squared matrix elements in table 1 (decay) while the rest is obtained by using that the sum of the squared matrix elements from table 2 (collisions) can be written as $`\|M\|^2\left(1+\stackrel{~}{g}_L^2+g_R^2\right)\left[(p_1p_2)(p_3p_4)+2(p_1p_4)(p_3p_2)\right].`$ (24) The assumption of Boltzmann statistics for massless leptons considerably simplifies the calculations, because the integral over their momenta can be taken explicitly. This approximation results in a larger value of the collision integral, i.e. in faster decay and reaction rates and correspondingly to a smaller abundance of $`\nu _s`$. Thus the restriction on $`\tau _{\nu _s}`$ obtained in this approximation is weaker than the real one. At high temperatures, $`T>5`$ MeV, (23) was integrated with the simplifying assumption $`Tx=1`$. This also leads to a weaker bound on $`\tau _{\nu _s}`$. Assuming entropy conservation one can check that $`Tx`$ may change by a factor 1.1 due to $`\nu _s`$ decays and annihilation. In fact the effect is stronger because the number density of $`\nu _s`$ is much larger than the equilibrium one. However, since we started at $`T=5`$ MeV $`<m_{\nu _s}`$, when the energy density of $`\nu _{\nu _s}`$ was already somewhat suppressed, the variation of $`Tx`$ from that moment would be weaker. Naively one would expect a variation of $`f_{\nu _s}`$ similar to the variation of $`Tx`$. However, it is argued below that the effect may be much stronger. Indeed, since $`\nu _s`$ disappears in the collisions with massless particles, its number density is much more sensitive to the variation of the coefficient in front of the collision integral. In the case of the standard freezing in two-body annihilation the frozen number density, $`n_f`$, is known to be inversely proportional to the annihilation cross-section. This is not true for the case of reactions (14). One can solve the kinetic equation explicitly and check that the result is exponentially sensitive to the coefficient in front of $`(f_{\nu _s}^{\mathrm{eq}}f_{\nu _s})`$. Therefore we took the variation of $`Tx`$ into account, starting from $`x_i=0.2`$, according to the energy balance law (22), see also (34). A simpler and more common method based on entropy conservation is not accurate enough because the sterile neutrinos strongly deviate from equilibrium in the essential range of temperatures and thus entropy is not conserved. In the non-relativistic limit we may integrate both sides of (23) over $`d^3y/(2\pi )^3`$ to obtain the number density of $`\nu _s`$ in the comoving volume $`\overline{n}_{\nu _s}=x^3n_{\nu _s}=\overline{n}_{\nu _s}^{\mathrm{eq}}{\displaystyle _0^x}𝑑x_1{\displaystyle \frac{d\overline{n}_{\nu _s}^{\mathrm{eq}}}{dx_1}}\mathrm{exp}\left[{\displaystyle _{x_1}^x}𝑑x_2B(x_2)\right],`$ (25) where $`B(x)={\displaystyle \frac{1.48x}{(Tx)^2(\tau _{\nu _s}/\mathrm{sec})}}\left({\displaystyle \frac{10.75}{g_{}(T)}}\right)^{1/2}\left[1+{\displaystyle \frac{3\times 2^7(Tx)^3}{x^3m_{\nu _s}^3}}\left({\displaystyle \frac{3\zeta (3)}{4}}+{\displaystyle \frac{\frac{7\pi ^4}{144}(Tx)}{xm_{\nu _s}}}\right)\right].`$ (26) Usually $`Tx`$ is assumed to be constant, quite often normalized as $`Tx=1`$. A small variation of this quantity usually is not very important for the value of $`\overline{n}`$. In our case, however, estimating the integral by the saddle point method reveals that the dependence on $`Tx`$ appears in the exponent and the effect may be significant even for a small variation of $`Tx`$. The dependence on $`Tx`$ for the scattering term results in an increase of $`\overline{n}_{\nu _s}`$, while that for the decay term results in a decrease of $`\overline{n}_{\nu _s}`$. The results of the numerical solution of (23) for different values of lifetimes are presented in figs. 1, 2, and 3. The first one shows the evolution of $`x^3n_{\nu _s}`$ for $`\tau _{\nu _s}=0.1,\mathrm{\hspace{0.17em}0.2},\mathrm{and}\mathrm{\hspace{0.17em}\hspace{0.17em}0.3}`$, the second is $`x^4\rho _{\nu _s}`$, while the third figure is a snap-shot of the distribution functions, $`f_{\nu _s}`$, at $`x=1`$. All the distributions are significantly higher than the equilibrium ones. For example, the equilibrium distribution is about 10 orders of magnitude below the calculated curve for $`\tau _{\nu _s}=0.1`$. ### 2.4 Light Neutrinos At the onset of nucleosynthesis, $`\nu _s`$ practically disappeared from the primeval plasma, at least for sufficiently small lifetimes, which are near the bound obtained below. However, the products of their decays and annihilation distorted the standard nucleosynthesis conditions: the energy density as well as the spectrum of neutrinos were different from the standard ones, and this changed the light element abundances. We will look for the solution of the kinetic equations governing the evolution of the distribution functions of light neutrinos in the form $`f_{\nu _a}(x,y)=\left(1+e^y\right)^1+\delta f_{\nu _a}(x,y),`$ (27) where $`a=e,\mu ,\mathrm{or}\tau `$. The first term is equal to the equilibrium distribution function for the case when the temperature drops as the inverse scale factor, i.e. $`Tx=1`$. In reality this is not the case and the function $`\delta f`$ takes into account both the variation of $`Tx`$ and the spectrum modification of the active neutrinos. The distribution of $`e^\pm `$ always has the equilibrium form, $`f_e=[1+\mathrm{exp}(E/T)]^1=[1+\mathrm{exp}(y/Tx)]^1`$ and the product $`Tx`$ is taken in the form $`Tx=1+\mathrm{\Delta }(x),`$ (28) where $`\mathrm{\Delta }`$ is assumed to be small, such that a perturbative expansion in $`\mathrm{\Delta }`$ can be made. A similar first-order perturbative expansion is made with respect to $`\delta f`$, so that the collision integral becomes linear in terms of $`\mathrm{\Delta }`$ and $`\delta f`$. It is assumed that initially both $`\delta f`$ and $`\mathrm{\Delta }`$ were zero. To this end one should find an appropriate value of the initial “time”, $`x_i`$, or temperature, $`T_i`$. The temperature should be sufficiently small, such that the number density of $`\nu _s`$ is already low and hence the function $`\mathrm{\Delta }`$ would not rise too much. On the other hand, the temperature should not be too low, otherwise equilibrium between neutrinos and electrons would not be established and the assumption of $`\delta f=0`$ could be grossly wrong. As we see in what follows, convenient initial values that satisfy both conditions are $`x_i0.2`$ and $`T_i5`$ MeV for $`\tau _{\nu _s}0.2`$ sec. Of course, the later rise of temperature by the well known factor 1.4 because of $`e^+e^{}`$-annihilation is taken into account explicitly. ### 2.5 Collision Integral and Source Term There are two kinds of terms in the collision integral, “easy” ones where the unknown function $`\delta f`$ depends upon the external variable $`y`$, which is not integrated upon, and “difficult” ones when $`\delta f`$ is under the integral. The terms of the first kind come, with negative sign, from external particles in the initial state. They give $`_x\delta f_{\nu _e}(x,y)=0.26\left({\displaystyle \frac{10.75}{g_{}}}\right)^{1/2}\delta f_{\nu _e}(1+g_L^2+g_R^2)(y/x^4),`$ (29) where $`g_L=\mathrm{sin}^2\theta _W+0.5`$, while $`g_R=\mathrm{sin}^2\theta _W`$. For $`\nu _\mu `$ and $`\nu _\tau `$ the functions $`\stackrel{~}{g}_L`$ and $`g_R`$ are given by expressions (8). Using the simple expression (29) one finds that the rate of approach of $`\nu _e`$ to equilibrium is given by $`\delta f_{\nu _e}\mathrm{exp}(0.13y/x^3).`$ (30) One sees that equilibrium would be efficiently restored for $`\nu _e`$ if $`x<0.5y^{1/3}`$ or $`T>\mathrm{\hspace{0.17em}2}y^{1/3}`$ MeV, while for $`\nu _{\mu ,\tau }`$ equilibrium would be restored if $`T>2.25y^{1/3}`$ MeV. These results are close to the standard estimates known in the literature (for the most recent ones see e.g. ). However, there are the following drawbacks in this derivation. First, it was done under the assumption that the source of distortion acted for a finite time and thus (30) could be valid only after the source has been switched off. If there is a constantly working source, $`S(x,y)`$, equilibrium is always distorted roughly by the factor $`S/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is the effective reaction rate. Another and more serious argument against the validity of (30) is that the “difficult” part of the collision integral was neglected. One can see from the general expressions (17,18) that $`\delta f`$ also appears under the integral over momenta, and it can be seen that it comes mostly with a positive sign. These terms counteract the smoothing action of expression (29) and shift equilibrium restoration to considerably higher temperatures. One can take them into account exactly, numerically solving the kinetic equations with the exact collision integral that can be reduced down to a 2-dimensional integration over particle energies (see e.g. ), but this is a very time consuming procedure. Instead we will here use a much simpler approach. We will approximately represent such “difficult” terms by the integrals over energy chosen in such a way that the kinetic equations satisfy the law of particle conservation in the comoving volume if only elastic scattering is taken into account, i.e. the kinetic equation should automatically give $`_x(x^3n)=0`$. The exact equations should simultaneously satisfy energy conservation law, but working with the exact equations is much more complicated. We checked the validity of our approximate procedure for a case of a more simple reaction amplitude, where we compared the approximate and exact results and found very good agreement. The approximation that we use breaks the energy conservation law in the kinetic equation, so that some small fraction of energy going into light particles from $`\nu _s`$ decays and/or annihilation is lost. This diminishes the effects that we are discussing so that the real constraints on $`\tau _{\nu _s}`$ should be stronger. The kinetic equations for light neutrinos in this approximation can be written as $`_x\delta f_{\nu _e}(x,y)`$ $`=`$ $`S_{\nu _e}(x,y)+0.26\left({\displaystyle \frac{10.75}{g_{}}}\right)^{1/2}(1+g_L^2+g_R^2)(y/x^4)`$ $`\times \{\delta f_{\nu _e}+{\displaystyle \frac{2}{15}}{\displaystyle \frac{e^y}{1+g_L^2+g_R^2}}[1+0.75(g_L^2+g_R^2)]`$ $`\times \left[{\displaystyle 𝑑y_2y_2^3\delta f_{\nu _e}(x,y_2)}+{\displaystyle \frac{1}{8}}{\displaystyle 𝑑y_2y_2^3\left(\delta f_{\nu _\mu }(x,y_2)+\delta f_{\nu _\tau }(x,y_2)\right)}\right]`$ $`+{\displaystyle \frac{3}{5}}\mathrm{\Delta }(x){\displaystyle \frac{g_L^2+g_R^2}{1+g_L^2+g_R^2}}e^y(11y/121)\}.`$ The coefficients $`g_{L,R}`$ for $`\nu _{\mu ,\tau }`$ are given by (8) and for $`\nu _e`$ are presented after (29). The source term $`S`$ describes injection of non-equilibrium neutrinos by $`\nu _s`$ decays or reactions (14) with light leptons. In what follows we include only decays. We have estimated the contributions of the reactions (14) and found that they slightly improve the restrictions that we have obtained. For $`\nu _e`$, $`\nu _\mu `$, and $`\nu _\tau `$ the contributions of the decay term are respectively $`S_{\nu _e,\nu _\mu }(x,y)`$ $`=`$ $`{\displaystyle \frac{0.012}{\tau _{\nu _s}x^2}}\left({\displaystyle \frac{10.75}{g_{}}}\right)^{1/2}\left(1{\displaystyle \frac{16y}{9m_{\nu _s}x}}\right)\left(n_{\nu _s}n_{\nu _s}^{\mathrm{eq}}\right)\theta \left(m_{\nu _s}x/2y\right),`$ (32) $`S_{\nu _\tau }(x,y)`$ $`=`$ $`{\displaystyle \frac{0.024}{\tau _{\nu _s}x^2}}\left({\displaystyle \frac{10.75}{g_{}}}\right)^{1/2}\left[1{\displaystyle \frac{16y}{9m_{\nu _s}x}}+{\displaystyle \frac{2}{3}}\left(1+\stackrel{~}{g}_L^2+g_R^2\right)\left(1{\displaystyle \frac{4y}{3m_{\nu _s}x}}\right)\right]`$ (33) $`\times \left(n_{\nu _s}n_{\nu _s}^{\mathrm{eq}}\right)\theta \left(m_{\nu _s}x/2y\right),`$ where $`n_{\nu _s}(x)`$ is the number density of $`\nu _s`$ and $`\theta (y)`$ is the step function which ensures energy conservation in the decay. The factor $`0.012`$ comes from the product of the branching ratio BR$`=96/4/(1+\stackrel{~}{g}_L^2+g_R^2)=21.3`$ with the factors $`3/m_{\nu _s}^3=7.7\times 10^5`$ ($`m_{\nu _s}`$ in MeV) from the normalization, we divide by $`2/\pi ^2=0.203`$ from the number density, and 6.582 from the relation between MeV and sec. Dividing by the Hubble parameter gives a factor, $`1.221/(1.88\sqrt{8\pi /3})=0.2244`$, and we find $`21.3\times (7.7\times 10^5)\times 0.2244\times 6.582/0.203=0.012`$. The coefficient in front of the collision integral is in reality momentum dependent, and hence is slightly different from 0.4 for $`\nu _e`$ (and 0.29 for $`\nu _{\mu ,\tau }`$), which are often used in the literature. We extracted the correct momentum dependent coefficients from our Standard Model code , and used this in the calculations. This slightly weakens the bound on the lifetime. The function $`\mathrm{\Delta }(x)`$ is determined from the energy balance condition (22) which in the present case reads $`{\displaystyle \frac{d\mathrm{\Delta }}{dx}}={\displaystyle \frac{1}{4x^4\rho _{EM}}}\left[{\displaystyle \frac{xd(x^3\rho _{\nu _s})}{dx}}+{\displaystyle \frac{d(x^4\delta \rho _\nu )}{dx}}\right],`$ (34) where $`\rho _{EM}`$ is the energy density in the electromagnetic sector, and $`\rho _{\nu _s}`$ and $`\rho _\nu `$ are the energy density of $`\nu _s`$ and the total energy density of all light neutrinos respectively. We also used another method for calculating $`Tx`$ starting from earlier times and obtained stronger results, so we believe that the limits that we obtain here are quite safe. ### 2.6 A Few Numerical Technicalities Let us summarize a few technicalities related to the numerical approach. We divide the time into 3 regions. First we integrate only (23) from very high temperatures, $`T=50`$ MeV ($`x=0.02`$) and until $`T=5`$ MeV ($`x=0.2`$). We assume that initially the sterile particles are in equilibrium (see sec. 2.2). In this way we can follow the freeze out and initial decay of the sterile particles, without worrying about the equilibrium active neutrinos and the electromagnetic plasma. In the next region, $`0.2<x<50`$, we solve (23, 2.5, 34) with the initial values $`\mathrm{\Delta }=0`$ and $`\delta f_{\nu _{\mathrm{active}}}=0`$. It is worth mentioning that from the very beginning we separately calculate and include the annihilation of the electrons, which increase the photon temperature with a standard factor $`1.4`$. Finally, for very high $`x>50`$ ($`T<0.03`$ MeV) when all sterile neutrinos have disappeared and the active neutrinos have long decoupled, we solved only the kinetic equations governing the $`n`$-$`p`$-reactions needed for the nucleosynthesis code. We use an 800 point grid in momentum in the region $`0<y<80`$, and checked that the results are insensitive to the doubling of the grid. For the BBN calculations we use $`\eta _{10}=5`$. ### 2.7 Results We have solved (23, 2.5, 34) numerically for different lifetimes. In fig. 1 we plot the evolution of the normalized number density, $`x^3n_{\nu _s}`$, as a function of $`x`$. One sees for $`\tau _{\nu _s}=0.3`$ how $`\nu _s`$ first freeze out, followed by the subsequent decay. In fig. 2 a similar plot of the normalized energy density, $`x^4\rho _{\nu _s}`$, is presented as a function of $`x`$, and the effect of the sterile particle being massive is evident. In fig. 3 we present a snap-shot of the distribution function at the time $`x=1`$. The equilibrium distribution function at this time is about 10 orders of magnitude smaller than the curve for $`\tau =0.1`$ sec. The calculated energy densities of all light neutrino species relative to the electromagnetic energy density, $`\rho _e+\rho _\gamma `$, are presented in fig. 4 as functions of $`x`$. One sees that this fraction is higher for longer lifetimes, especially around the time $`x=1`$, when the $`n/p`$-ratio freezes out, leading to an expected increase in the final helium abundance. In fig. 5 a snap-shot of the spectrum of $`\nu _e`$, namely $`y^2f_{\nu _e}`$ and the distortion $`y^2\delta f_{\nu _e}`$ are presented for $`x=1`$ for the lifetimes $`0.1,0.2,0.3`$. A distortion of the electronic neutrino spectrum has a strong impact on nucleosynthesis, while $`\nu _\mu `$ and $`\nu _\tau `$ act only by their total energy density. It is noteworthy that the increase in $`\rho _{\nu _e}`$ acts in the opposite direction to an increase in $`\rho _{\nu _\mu ,\nu _\tau }`$, since it reduces the effective number of light neutrinos, or in other words, it gives rise to a smaller mass fraction of primordial <sup>4</sup>He, while an increase in the high energy part of $`\nu _e`$ spectrum results in a larger mass fraction of <sup>4</sup>He. The results of the calculations have been imported into the modified Kawano nucleosynthesis code, and the abundances of all light elements have been calculated. At each time step $`x`$, we find the corresponding photon temperature and total energy density. Furthermore we integrate the kinetic equation governing the $`n/p`$ evolution taking into account the distorted spectrum of $`\nu _e`$. The final helium abundance is presented as a function of the $`\nu _s`$ lifetime in fig. 6. By translating these results into effective number of neutrinos one sees that if we allow for $`\mathrm{\Delta }N=0.2`$ (as suggested by ), then only lifetimes lower than $`\tau _{\nu _s}=0.17`$ sec are permitted. If one is more conservative and allows for one extra neutrino species, $`\mathrm{\Delta }N=1.0`$, then lifetimes longer than $`\tau _{\nu _s}=0.24`$ sec are excluded. A drop in the helium abundance around $`\tau _{\nu _s}=0.1`$ is related to the dominant role of the $`\nu _e`$ energy density, since the spectrum distortion is shifted to smaller energies. Finally in fig. 7 we compare the KARMEN experimental data with the bound obtained from BBN. One sees that the expected new data from the NOMAD Collaboration together with our BBN bound leave a small allowed window for a sterile neutrino with a lifetime around 0.1–0.2 sec. ## 3 Supernova Limits ### 3.1 Small Mixing Angle The mixing angles of sterile neutrinos with the standard active flavors is tightly constrained by standard arguments related to supernova (SN) physics and to the neutrino observations of SN 1987A. Some of these arguments have been sketched out in the context of the KARMEN anomaly in Ref. . The simplest limit arises from the “energy-loss argument.” The SN 1987A observations imply that a SN core may not emit too much energy in an “invisible channel” as this would unduly shorten the observed neutrino burst. Reasonably accurate limits are obtained by demanding that the “exotic” energy-loss rate should obey $$ϵ\begin{array}{c}<\hfill \\ \hfill \end{array}10^{19}\mathrm{erg}\mathrm{g}^1\mathrm{s}^1,$$ (35) where $`ϵ`$ is to be calculated at typical average conditions of a SN core ($`\rho =3\times 10^{14}\mathrm{g}\mathrm{cm}^3`$, $`T=30\mathrm{MeV}`$). Sterile neutrinos $`\nu _s`$ are produced because they mix with one of the standard ones. If the standard neutrino masses are all in the sub-eV range, the assumed sterile neutrino mass of $`m_{\nu _s}=33.9\mathrm{MeV}`$ assures a mass difference so large that medium effects on the oscillations can be neglected. The oscillation frequency is so large that a standard neutrino $`\nu _a`$, once produced, oscillates many times before collisions interrupt the coherent development of the flavor amplitude. The average probability of finding the original $`\nu _a`$ in the $`\nu _s`$ flavor state is $`\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })`$ where $`\mathrm{\Theta }`$ is the $`\nu _a`$-$`\nu _s`$-mixing angle. Oscillations are interrupted with the collision rate $`\mathrm{\Gamma }`$ of the $`\nu _a`$ flavor, leading to the standard sterile-neutrino production rate of $`\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })\mathrm{\Gamma }`$. We estimate $`\mathrm{\Gamma }`$ as the neutral-current collision rate on free nucleons, ignoring correlation and degeneracy effects, so that: $$\mathrm{\Gamma }=\frac{C_V^2+3C_A^2}{\pi }G_F^2n_BE_\nu ^2,$$ (36) with $`G_F`$ the Fermi constant and $`n_B`$ the number density of baryons (nucleons). For a mix of protons and neutrons we use an average neutral-current coupling constant of $`(C_V^2+3C_A^2)1`$. If we further assume that the trapped active neutrinos do not have a significant chemical potential (true for $`\nu _\mu `$ and $`\nu _\tau `$, but not for $`\nu _e`$), the energy-loss rate is $$ϵ_{\nu _s}=\frac{\mathrm{sin}^2(2\mathrm{\Theta })}{2}\frac{G_F^2}{\pi ^3m_N}_0^{\mathrm{}}𝑑E_\nu \frac{E_\nu ^5}{e^{E_\nu /T}+1},$$ (37) where $`m_N`$ is the nucleon mass. For simplicity we use Maxwell-Boltzmann statistics for the neutrinos (we ignore the $`+1`$ in the denominator under the integral) and find: $`ϵ_{\nu _s}`$ $`=`$ $`{\displaystyle \frac{\mathrm{sin}^2(2\mathrm{\Theta })}{2}}{\displaystyle \frac{120G_F^2T^6}{\pi ^3m_N}},`$ $`=`$ $`\mathrm{sin}^2(2\mathrm{\Theta })\mathrm{\hspace{0.17em}2.8}\times 10^{26}\mathrm{erg}\mathrm{g}^1\mathrm{s}^1T_{30}^6`$ where $`T_{30}=T/30\mathrm{MeV}`$. Comparing this result with (35) leads to a limit $$\mathrm{sin}^2(2\mathrm{\Theta })\begin{array}{c}<\hfill \\ \hfill \end{array}3\times 10^8T_{30}^6.$$ (39) The temperature $`T=30\mathrm{MeV}`$ is at the lower end of what is found in typical numerical calculations so that this limit is reasonably conservative. For the mixing with $`\nu _e`$ the limit is more restrictive because the electron neutrinos in a SN core are highly degenerate, leading to a larger conversion rate and a larger amount of energy liberated per collision. Several authors found $`\mathrm{sin}^2(2\mathrm{\Theta })\begin{array}{c}<\hfill \\ \hfill \end{array}10^{10}`$ for this case . In our calculation of the emission rate we have taken the $`\nu _s`$ to be effectively massless. The average energy of trapped standard neutrinos is about $`3T`$ which far exceeds $`m_{\nu _s}`$. The scattering cross section scales with $`E_\nu ^2`$, favoring the emission of high-energy neutrinos. The average $`\nu _s`$ energy emerging from a SN core is thus found to be about $`5T`$ so that neglecting $`m_{\nu _s}`$ is a good approximation—the $`\nu _s`$ are highly relativistic. The KARMEN experiment implies that the lifetime of $`\nu _s`$ exceeds about $`0.1\mathrm{s}`$ so that these particles escape from the SN before decaying. On the other hand, they must decay on the long way from SN 1987A to us. If the decay products include $`e^+e^{}`$ pairs one will also get $`\gamma `$ rays from inner bremsstrahlung in the decay. The non-observation of a $`\gamma `$ ray burst coinciding with SN 1987A then leads to further limits . The decay products likely would include standard neutrinos which would have shown up in the detectors. However, because of the high energies of the sterile neutrinos which are representative of the SN core temperature, such events would have much larger energies than those expected from thermal neutrino emission at the neutrino sphere. Therefore, the emission of sterile neutrinos from the core and their subsequent decay cannot mimic the standard SN neutrino signal. ### 3.2 Large Mixing Angle These upper limits on the mixing angle are only valid if the sterile neutrinos actually escape from the SN core after production, a condition that is not satisfied if one of the mixing angles is too large. Since the mixing angle between $`\nu _e`$ or $`\nu _\mu `$ and $`\nu _s`$ is already constrained from laboratory experiments to be small in this sense, we worry here only about the $`\nu _\tau `$-$`\nu _s`$ mixing angle. We continue to assume that the standard-neutrino mass eigenstates are small so that the mass difference and hence the oscillation frequency between $`\nu _\tau `$ and $`\nu _s`$ will be large compared to a typical collision rate in a SN core and that the mixing angle in the medium is identical to the vacuum mixing angle. Therefore, after production the chance of finding a $`\nu _s`$ in the active $`\nu _\tau `$ state will be given by the average value $`\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })`$. If the mixing angle is large this will mean that the sterile neutrino essentially acts as an active one. It will be trapped, and the energy-loss argument is not applicable. If the mixing angle is not quite maximal, sterile neutrinos will still be trapped, but their mean free path will be larger than that of an active flavor. Energy is transported out of a SN core by neutrino diffusion, a mechanism which is more effective if the mfp is larger. Simply put, a neutrino can transport energy over distances of order the mfp so that more distant regions are thermally coupled if the mfp becomes larger. Most of the transporting of energy is done by the most weakly coupled particles which are still trapped. For example, in a SN core the photon contribution to the energy transport is negligible because their mfp is very much smaller than that of standard neutrinos. Conversely, the contribution of sterile neutrinos will increase with an increasing mfp, i.e. with a decreasing $`\mathrm{sin}^2(2\mathrm{\Theta })`$. The effect on the SN 1987A signal will be identical to the effect of freely escaping sterile neutrinos, i.e. the signal will be shortened. We stress that the signal duration is determined by the diffusion time scale throughout the star. Therefore, increasing the mfp in the deep interior of the star shortens the cooling time scale. In a numerical study the neutrino opacities were artificially decreased. It was found that the efficiency of neutrino transfer in the star should not be more than about twice the standard value to remain consistent with the SN 1987A signal characteristics. Likewise, in a different study the number of standard neutrino flavors was artificially increased, again leading to an increased efficiency of energy transfer. Doubling the effective number of neutrino flavors appears excluded from the SN 1987A data. The interaction rate of sterile neutrinos is that of a standard $`\nu _\tau `$, times $`\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })`$. Conversely, the collision rate for $`\nu _\tau `$ is $`1\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })`$ times the standard rate because a $`\nu _\tau `$ has an average chance $`1\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })`$ of being measured as a $`\nu _s`$. Assuming that the standard transfer of energy is dominated by $`\nu _\mu `$, $`\nu _\tau `$ and their anti-particles (the mfp for $`\nu _e`$ and $`\overline{\nu }_e`$ is much shorter due to charged-current reactions), then adding the sterile neutrino will enhance the rate of energy transfer by a factor $$\frac{1}{1+1}\left(1+\frac{1}{1\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })}+\frac{1}{\frac{1}{2}\mathrm{sin}^2(2\mathrm{\Theta })}\right).$$ (40) Maximal mixing corresponds to $`\mathrm{sin}^2(2\mathrm{\Theta })=1`$, implying that both $`\nu _\tau `$ and $`\nu _s`$ each scatter with half the standard rate, so their mfp is each increased by a factor of 2. Moreover, the sterile neutrino contributes a second channel for the transfer of energy. Therefore, the energy flux carried by maximally mixed $`\nu _s`$ and $`\nu _\tau `$ is four times that carried by a standard $`\nu _\tau `$. This explains the limiting behavior of (40) for maximal mixing. The new experimental limit on the $`\nu _s`$-$`\nu _\tau `$ mixing angle is $`\mathrm{sin}^2(2\mathrm{\Theta })<10^2`$. Therefore, the efficiency of energy transfer out of a SN core would be enhanced by more than two orders of magnitude. Such an enhancement is certainly not compatible with the SN 1987A signal, implying that the mixing angle has to be very small, i.e. that it must obey (39). Therefore, the loop hole of a large $`\nu _s`$-$`\nu _\tau `$ mixing angle has been plugged by the new experimental constraints. ## 4 Conclusion Cosmological and astrophysical arguments seem to exclude the interpretation of the KARMEN anomaly by an unstable sterile neutrino mixed with $`\nu _\tau `$. The arguments based on SN 1987A are stronger than the nucleosynthesis bound. The supernova limit is given by (39), i.e. $`\tau _{\nu _s}>6\times 10^4`$ sec. This result, together with the direct experimental constraints on $`\nu _\tau `$, completely excludes a 33.9 MeV sterile neutrino. Primordial nucleosynthesis permits to exclude roughly $`\tau _{\nu _s}>0.2`$ sec. So with the existing direct experimental limits, some window, $`\tau _{\nu _s}=0.10.2`$ sec, remains open. The simplifications that we made in deriving the BBN bound typically lead to a weaker result, so the real bound may be somewhat stronger. However, if exact calculations confirm the decrease of <sup>4</sup>He around $`\tau _{\nu _s}=0.1`$ found in this paper, then BBN will never exclude lifetimes $`0.1`$ sec. The effect happens to be surprisingly sensitive to usually neglected phenomena, particularly the shape of the $`\nu _e`$ spectrum and to the non-adiabatic variation of temperature. We hope to do exact (but rather long) calculations later. Together with a possible improvement of the observational data on light element abundances and a better understanding of the subsequent changes of light elements in the course of cosmological evolution, the BBN bound may become competitive with the supernova one. If so, the very attractive hypothesis of a 33.9 MeV sterile $`\nu _s`$ would be killed by two independent arguments. Moreover, one may imagine a case when the supernova arguments are not applicable, while the nucleosynthesis ones still operate. For example, if $`\nu _s`$ possesses an anomalously strong interaction with nucleons (stronger than the usual weak one), then it would not noticeably change the energetics of supernovae that we described here, but would affect nucleosynthesis practically at the same level as discussed in section 2. This new interaction might be related to the anomalously high mass of $`\nu _s`$. If this is the case then a small window for a 33.9 MeV $`\nu _s`$ with lifetime $`\tau _{\nu _s}=0.1`$–0.2 sec, still exists at the present time, and stronger experimental bounds, as well as more accurate calculations of the impact of $`\nu _s`$ on primordial nucleosynthesis, are needed. ## Acknowledgment We are grateful to N. Krasnikov who brought this subject to our attention. In Munich, this work was partly supported by the Deutsche Forschungsgemeinschaft under grant No. SFB 375. DS thanks the NOMAD collaboration, in particular L. DiLella, L. Camilleri, S. Gninenko and A. Kovzelev, for the invitation and hospitality during their December-1999 collaboration meeting where this work was presented.
warning/0002/nlin0002007.html	ar5iv	text	# A family of multi-value cellular automaton model for traffic flow ## 1 Introduction Traffic phenomena attract much attention of physicists in recent years. It shows a complex phase transition from free to congested state, and many theoretical models have been proposed so far. Among them we will focus on deterministic cellular automaton (CA) models. CA models are simple, flexible, and suitable for computer simulations of discrete phenomena. The rule–184 CA has been widely used as a prototype of deterministic model of traffic flow. In the model, lane is single and cars can move by one site at most every time step. Its several variations have been proposed recently: First, Fukui and Ishibashi (FI) proposed a high speed extension of the rule–184 CA. In the model, cars can move more than one site per unit time if there are successive vacant spaces ahead. Second, Fukś and Boccara proposed a ‘monitored traffic model’, which is a kind of quick start (QS) model. In the model, drivers can prospect vacant spaces due to car motion in the next site and can move more quickly compared with the rule–184 CA. Third, Takayasu and Takayasu proposed a slow start (SlS) model, in which cars can not move just after they stop and wait for a unit time. This represents an asymmetry of stopping and starting behavior. We call all these variants a family of rule–184 CA in this paper. A desirable condition for CA as a traffic model is that it can show a phase transition observed in real data. Let us see an example of real data of Tomei expressway taken by the Japan Highway Public Corporation. Fig.1 shows flow–density diagrams often called ‘fundamental diagram’, of both lanes in January 1996 on up line at 170.64 km point from Tokyo. Fig.1 (a) and (b) are diagrams of driving and acceleration lane respectively. We see that a phase transition from free to congested state around a density of 25 vehicles$`/`$km in both lanes, and there is a clearer discontinuity in the acceleration lane. The multiple states of flow at the same density around the critical density are also observed particularly in the acceleration lane, and a skeleton curve of plot points shows a shape of ‘inverse $`\lambda `$’. These phenomena are often observed in real data. There is another observed fact that over a certain critical density, a perturbation to a uniform traffic causes a formation of jam, and ‘stop-and-go’ wave propagates backward. Cars accelerate and decelerate alternately in the wave. We adopt these two experimental facts, multiple states and stop-and-go wave, as criteria to judge whether a CA is suitable for traffic model or not. Among the rule–184 CA family, only SlS model shows multiple states and also shows a stop-and-go wave. Recently, a multi-value generalization of the rule–184 CA has been proposed by using an ultradiscrete method. Its evolution equation is $$U_j^{t+1}=U_j^t+\mathrm{min}(U_{j1}^t,LU_j^t)\mathrm{min}(U_j^t,LU_{j+1}^t),$$ (1) where $`U_j^t`$ represents the number of cars at site $`j`$ and time $`t`$, and $`L`$ is an integer constant. Each site is assumed to hold $`L`$ cars at most. We can prove that if $`L>0`$ and $`0U_j^tL`$ for any $`j`$, then $`0U_j^{t+1}L`$ holds for any $`j`$. Thus (1) is considered to be $`(L+1)`$–value CA. Since (1) is obtained from an ultradiscretization of the Burgers equation, we call it Burgers CA (BCA). Note that BCA is equivalent to the rule–184 CA in a special case of $`L=1`$. The positive integer $`L`$ can be interpreted physically in the following three ways: First, the road is an $`L`$–lane freeway in a coarse sense, and effect of lane-changing rule is not expressed explicitly. In Appendix, we show BCA with $`L=2`$ can be interpreted as a two lane model of the rule–184 CA with an explicit lane–changing rule. Second, $`U_j^t/L`$ represents a probability distribution of cars in a single-lane freeway. In this case, $`U_j^t`$ itself no longer represents a real number of cars at site $`j`$. Third, also in a single-lane freeway, length of each site is assumed to be long enough to contain $`L`$ cars. By introducing the free parameter $`L`$ into the rule–184 CA family, we can generalize them to multi-value CA and obtain rich algebraic properties and wide applications to various transport phenomena. Though BCA does not show multiple states at a density in the fundamental diagram, we have shown that its high speed extension shows multiple states around a critical density. Metastable states of flow in the model exist on a prominent part of a line in a free phase over a critical density. The model is also considered to be a multi-value generalization of FI model. This model is mentioned in Sec.2.3 to compare other models. Other models in the rule–184 CA family, SlS model and QS model, are two-value CA like FI model. We generalize them to multi-value ones and investigate their properties as a traffic model in detail. We use a max–plus representation to express time evolution rule of those generalized models. This representation comes from the ultradiscrete method which has a close relation with the max–plus algebra, as shown in our previous papers. Using max–plus representation, we can express evolution rule of the above CA’s in a conservation from like (1), which automatically derives a conservation of total number of cars. ## 2 Generalization to multi-value CA In this section, we will present four CA models which are multi-value generalization of the rule–184 CA family. Before explaining models, we shortly review a rule of car movement of BCA. Note that we assume $`L`$ is a lane number of the road in the description of all models. In BCA, cars at site $`j`$ move to vacant spaces at their next site $`j+1`$ as many as possible. Therefore, car flow from $`j`$ to $`j+1`$ at time $`t`$ becomes $`\mathrm{min}(U_j^t,LU_{j+1}^t)`$. Thus adding $`\mathrm{min}(U_{j1}^t,LU_j^t)`$ (flow from $`j1`$ to $`j`$) and $`\mathrm{min}(U_j^t,LU_{j+1}^t)`$ (flow from $`j`$ to $`j+1`$) to $`U_j^t`$ (present car number), we obtain $`U_j^{t+1}`$ and derive (1). ### 2.1 Multi-value QS model First, we generalize QS model to a multi-value one. In the multi-value model, cars at site $`j`$ move to vacant spaces at site $`j+1`$ per unit time as many as possible. This movement rule is similar to that of BCA. However, a difference is that drivers at $`j`$ estimate vacant spaces at $`j+1`$ by predicting how many cars move from $`j+1`$ to $`j+2`$. Thus, they estimate vacant spaces at $`j+1`$ at time $`t`$ to be $`LU_{j+1}^t\text{(present spaces)}+\mathrm{min}(U_{j+1}^t,LU_{j+2}^t)\text{(predicted spaces due to BCA movement)}`$. Therefore considering the number of cars coming into and escaping from site $`j`$, evolution equation on $`U_j^t`$ is given by $`U_j^{t+1}`$ $`=`$ $`U_j^t+\mathrm{min}(U_{j1}^t,LU_j^t+\mathrm{min}(U_j^t,LU_{j+1}^t))`$ (2) $`\mathrm{min}(U_j^t,LU_{j+1}^t+\mathrm{min}(U_{j+1}^t,LU_{j+2}^t))`$ $`=`$ $`U_j^t+\mathrm{min}(U_{j1}^t,2LU_j^tU_{j+1}^t)\mathrm{min}(U_j^t,2LU_{j+1}^tU_{j+2}^t).`$ The final form of (2) means that the number of movable cars ($`U_j^t`$) is limited by vacant spaces ($`2LU_{j+1}^tU_{j+2}^t`$) in their next two sites. In the case of $`L=1`$, this model is identical to QS model proposed by Fukś and Boccara, and it is the rule–3212885888 CA of a neighborhood size ‘5’ after Wolfram’s terminology. ### 2.2 Multi-value SlS model In the original SlS model, lane is single and cars can not move just after they stop and wait for a unit time. Movable cars move to vacant spaces in their next site like BCA. In the multi-value case, we should distinguish standing cars and moving cars in each site because each site can hold plural cars. The number of cars at site $`j`$ blocked by cars at $`j+1`$ at time $`t1`$ is represented by $`U_j^{t1}\mathrm{min}(U_j^{t1},LU_{j+1}^{t1})`$. These cars cannot move at $`t`$, then the maximum number of cars movable to site $`j`$ at $`t`$ is given by $`U_j^t\{U_j^{t1}\mathrm{min}(U_j^{t1},LU_{j+1}^{t1})\}`$. Therefore, the multi-value SlS CA is given by $`U_j^{t+1}=U_j^t`$ $`+`$ $`\mathrm{min}[U_{j1}^t\{U_{j1}^{t1}\mathrm{min}(U_{j1}^{t1},LU_j^{t1})\},LU_j^t]`$ (3) $``$ $`\mathrm{min}[U_j^t\{U_j^{t1}\mathrm{min}(U_j^{t1},LU_{j+1}^{t1})\},LU_{j+1}^t)]`$ We note that this model includes the original SlS model in the case $`L=1`$. Since $`U^{t+1}`$ is determined by $`U^t`$ and $`U^{t1}`$, the evolution equation is second order in time and an effect of inertia of cars is included in this model. ### 2.3 EBCA2 model In our previous paper, we propose an extended BCA (EBCA) model in which cars can move two site forward at a unit time if the successive two sites are not fully occupied. In this paper, we call the model ‘EBCA2’ and another model with a similar extension ‘EBCA1’ described in the next subsection. A main difference between two models is a priority of movement of fast and slow cars to vacant spaces. Fast cars with speed 2 move prior to slow ones with speed 1 in EBCA2 and otherwise in EBCA1. Evolution equation of EBCA2 is given by $$U_j^{t+1}=U_j^t+\mathrm{min}(b_{j1}^t+a_{j2}^t,LU_j^t+a_{j1}^t)\mathrm{min}(b_j^t+a_{j1}^t,LU_{j+1}^t+a_j^t),$$ (4) where $`a_j^t\mathrm{min}(U_j^t,LU_{j+1}^t,LU_{j+2}^t)`$ which is the number of cars moving by two sites and $`b_j^t\mathrm{min}(U_j^t,LU_{j+1}^t)`$. In reference , we only study the cases of $`L=1`$ and 2. In this paper, we show some properties of flow for $`L`$ in the next section. Note that (4) in the case of $`L=1`$ is FI model which is equivalent to the rule–3436170432 CA. ### 2.4 EBCA1 model There is another possibility when we extend BCA to a high-speed model and we call this extended new model EBCA1. In contrast to EBCA2, slow cars with speed 1 move prior to fast ones with speed 2. Then car movement from $`t`$ to $`t+1`$ consists of the following two successive procedures, * Cars move to vacant spaces in their next site as many as possible. * Only cars moved in procedure a) can move more one site and they move to their next vacant spaces as many as possible. The number of moving cars at site $`j`$ and time $`t`$ in procedure a) is given by $`b_j^t\mathrm{min}(U_j^t,LU_{j+1}^t)`$. In procedure b), the number of moving cars at site $`j+1`$ becomes $`\mathrm{min}(b_j^t,LU_{j+2}^tb_{j+1}^t+b_{j+2}^t)`$, where the second term in $`\mathrm{min}`$ represents vacant spaces at site $`j+2`$ after the first procedure a). Therefore, considering a total number of cars entering into and escaping from site $`j`$, evolution equation of EBCA1 is given by $`U_j^{t+1}`$ $`=`$ $`U_j^t+b_{j1}^tb_j^t`$ (5) $`+\mathrm{min}(b_{j2}^t,LU_j^tb_{j1}^t+b_j^t)\mathrm{min}(b_{j1}^t,LU_{j+1}^tb_j^t+b_{j+1}^t)`$ $`=`$ $`U_j^t+\mathrm{min}(b_{j1}^t+b_{j2}^t,LU_j^t+b_j^t)\mathrm{min}(b_j^t+b_{j1}^t,LU_{j+1}^t+b_{j+1}^t).`$ Note that (5) with $`L=1`$ differs from FI model and it is equivalent to rule–3372206272 CA. We can consider that this rule is an extension of SlS model to high-speed motion since cars which stop at procedure a) cannot move at procedure b). Let us consider a relation between SlS and EBCA1 models in detail. If we write each procedure of EBCA1 explicitly using an intermediate time step, we obtain $`U_j^{t+1/2}`$ $`=`$ $`U_j^t+b_{j1}^tb_j^t`$ (6) $`U_j^{t+1}`$ $`=`$ $`U_j^{t+1/2}+\mathrm{min}(U_{j1}^{t+1/2}\{U_{j1}^tb_{j1}^t\},LU_j^{t+1/2})`$ (7) $`\mathrm{min}(U_j^{t+1/2}\{U_j^tb_j^t\},LU_{j+1}^{t+1/2}).`$ Equations (6) and (7) represent above procedure a) and b) respectively, and $`U_j^{t+1/2}`$ denotes car number at site $`j`$ just after procedure a). If we replace $`t+1`$ by $`t+1/2`$ in (1), we obtain (6). Moreover, if we replace $`t`$ by $`t+1/2`$ and $`t1`$ by $`t`$ in (3), we obtain (6). Therefore we can consider EBCA1 to be a ‘combination’ of BCA and SlS rules. ## 3 Fundamental diagram and multiple states We discuss fundamental diagrams of new CA models described in the previous section. In the followings, we will consider a periodic road or a circuit. All models in Sec.2 can be expressed in a conservation form such as $$\mathrm{\Delta }_tU_j^t+\mathrm{\Delta }_jq_j^t=0,$$ (8) where $`\mathrm{\Delta }_t`$ and $`\mathrm{\Delta }_j`$ are forward difference operator with respect to indicated variable, and $`q_j^t`$ represents a traffic flow. Average density $`\rho `$ and average flow $`Q^t`$ over all sites are defined by $$\rho \frac{1}{KL}\underset{j=1}{\overset{K}{}}U_j^t,Q^t\frac{1}{KL}\underset{j=1}{\overset{K}{}}q_j^t,$$ (9) where $`K`$ is number of sites in a period. Since all models are in a conservation form, average density does not depend on time and we can use $`\rho `$ without a script $`t`$. Figures 2 (a)–(d) are density–flow diagrams of previous models with $`L=2`$ and $`K=30`$. Figures 2 (a)–(d) corresponds to multi-value QS, multi-value SlS, EBCA2 and EBCA1 models respectively. If we set a number $`N`$ of cars, we obtain a unique density $`\rho =N/KL`$ but there are many initial distributions of cars. Each initial distribution does not always reach a steady flow and often make a periodic state. $`Q^t`$ can also changes periodically in time. Therefore, we plot every points in figures by averaging $`Q^t`$ from $`t=2K`$ to $`t=4K`$. (We use $`Q`$ as an average flow.) Note that $`t=2K`$ is long enough from initial time to obtain a periodic state and $`2K`$ time steps is much longer than its period. Moreover, we obtain different values of $`Q`$ from different initial distributions for a given $`\rho `$ in Figs.4 (b) $``$ (d). In Ref., this type of state giving different $`Q`$ values for the same density are called ‘multiple states’. We use this terminology for other models in this paper. Multiple states exist around a critical density. Especially for multi-value SlS model (Fig.2 (b)) and EBCA1 model (Fig.2 (d)), we see a thick branch other than lines forming an inverse $`\lambda `$ shape and the distribution of data in the branch look random. This fact is interesting because evolution rules are completely deterministic. Next, we focus on EBCA1 model and examine its properties in detail. Figure 3 is a skeleton diagram of EBCA1 model. Branches in Fig.3 are drawn as straight lines and we obtain them using specific initial conditions. For example, states ‘$`\mathrm{}11102021110202\mathrm{}`$’, ‘$`\mathrm{}121212\mathrm{}`$’, ‘$`\mathrm{}002002\mathrm{}`$’ and ‘$`\mathrm{}222222\mathrm{}`$’ are all steady states in the model, and flows of these states are plotted on B, C, D and E respectively in Fig.3. There are two branches in congested phase of higher density, that is, B–C and D–E. We call the branch D–E and B–C congested phase zero and phase one respectively. It is noted that there are many states not on the skeleton in Fig.3. For example, state $`\mathrm{}211211\mathrm{}`$ is a steady state of the model, and it has density 2/3 and flow 1/2. Fig.4 shows a time evolution of flow $`Q^t`$ in EBCA1 model with a density of phase transition region. Number of total sites $`K`$ is 60 and 240 in Fig.4 (a) and (b) respectively. In the figures, we can see periodic oscillation of flow appearing soon after an initial time. We observe that the maximum period of oscillation becomes longer as system size becomes larger. Moreover, oscillation in a period is not simple as shown in the figures. Fig.4 (c) shows a power spectrum $`I`$ versus frequency $`f`$ defined by $`I=\|_{t=0}^{T1}Q^t\mathrm{exp}(2\pi ift/T)\|/T`$. We set $`K=240`$ and $`T=1000`$. This indicates that the irregular oscillation is similar to white noise. Fig.5 is a fundamental diagram of EBCA1 with $`L=1`$. It is interesting that there exist multiple states even in the case $`L=1`$, which is equivalent to the rule–3372206272 CA. Among deterministic two-value CA models, only SlS model is known to show metastable state so far. Thus, EBCA1 with $`L=1`$ becomes the second example of such kind of two-value CA. At the end of this section, we give a comment on a case of $`L>3`$. Fundamental diagram of multi-value QS model is the same regardless of $`L`$. In other models, small branches increases around the critical density as $`L`$ becomes larger. Let us show this phenomenon using EBCA2 which has clear branches. The fundamental diagram of EBCA2 with $`L=7`$ is given in Fig.6. There are many branches around the critical density. We can give a partial explanation of these branches. Let us assume initial values on all sites are restricted to $`n`$ or $`Ln`$ where $`0n<L/2`$. Then, if we use a transformation $`V=(Un)/(L2n)`$ from $`U`$ to $`V`$, $`V=0`$ corresponds to $`U=n`$ and $`V=1`$ to $`U=Ln`$. We can easily show that evolution equation on $`V`$ is obtained by replacing $`U`$ by $`V`$ and $`L`$ by 1 in (4). Equation (4) with $`L=1`$ is FI model itself and it is two-value CA. Therefore, if all initial values of $`U`$ is restricted to $`n`$ or $`Ln`$, then $`U`$ at any site is always $`n`$ or $`Ln`$ and evolutional state is the same as that of FI model by replacing 0 by $`n`$ and 1 by $`Ln`$. Above fact implies a shape of the fundamental diagram. Consider an arbitrary state of FI model, that is, (4) with $`L=1`$. Let us assume density is $`\rho `$ and average flow is $`Q`$ for that state. We can obtain a corresponding state of (4) with $`L>0`$ by replacing 0 by $`n`$ and 1 by $`Ln`$. Then, using (4), we can easily show that density of that state is $`(1\frac{2n}{L})\rho +\frac{n}{L}`$ and flow is $`(1\frac{2n}{L})Q+\frac{2n}{L}`$. Moreover, we can also show that a fundamental diagram of FI model is exactly given by two segments of which end points are $`(0,0)`$, $`(\frac{1}{3},\frac{2}{3})`$ and $`(\frac{1}{3},\frac{2}{3})`$, $`(1,0)`$ respectively. Therefore, a diagram of EBCA2 with $`L>0`$ has at least a segment with end points $`(\frac{n}{L},\frac{2n}{L})`$ and $`(\frac{1}{3}+\frac{n}{3L},\frac{2}{3}+\frac{2n}{3L})`$ and that with $`(\frac{1}{3}+\frac{n}{3L},\frac{2}{3}+\frac{2n}{3L})`$ and $`(1\frac{n}{L},\frac{2n}{L})`$ for any $`n`$ ($`0n<L/2`$). Figure 6 is a $`L=7`$ case and the shape of diagram is strictly a superposition of all these segments. ## 4 Stability of flow In this section, we investigate stability of the branches in fundamental diagram obtained in the previous section. First, we study a uniform state $`\mathrm{}11111\mathrm{}`$, which gives the maximum flow in each model. Let us define a weak perturbation by a perturbation changing a local state $`11`$ to $`20`$, and a strong perturbation by a perturbation changing $`1111`$ to $`2200`$. Both perturbations clearly do not change the density. These perturbations mean that a car or two cars in the uniform flow suddenly decrease their speed and consequently fully-occupied sites ‘2’ appear. Figures 7 (a)–(c) show an instability of uniform flow by a weak perturbation in multi-value SlS, EBCA2 and EBCA1 models respectively. We can observe a stop-and-go wave propagating backward from the perturbed site in Fig.7 (a) (multi-value SlS) and (c) (EBCA1), while the wave does not appear in (b) (EBCA2). In Fig.7 (c), flow $`Q`$ decreases by the weak perturbation and transits from A to F in Fig.3. Moreover, we see that the final steady state contains locally two states corresponding to B and C in Fig.3. Flow of state B is 5/6 and that of C is 1/2, both states are in the phase one, and they give the maximum and minimum flow in that phase. Numerical results from various initial states with a weak perturbation shows that a state of higher density tends to give this type of extreme local states on a branch after some time steps. Moreover, if we give a strong perturbation of uniform state $`\mathrm{}111\mathrm{}`$ corresponding to A, it goes down directly to G(1/2,1/2) in the phase zero. Fig.7 (d) shows an instability due to a strong perturbation in EBCA1 model. In this case, stop-and-go wave does not occur. Moreover, we see that the final state consists of two local states $`\mathrm{}200200\mathrm{}`$ and $`\mathrm{}222222\mathrm{}`$ corresponding to D and E in Fig.3, which also give the maximum and minimum flow in the phase zero respectively. Since a perturbed uniform state becomes a congested one which consists of states giving the maximum and the minimum flow in a phase, we can predict a length of the final congested bunch. Assume $`\alpha `$ denotes a ratio of length of final congested bunch to length of total sites. Density of B, F and C are 5/12, 1/2 and 3/4 respectively. Thus we obtain $$\frac{5}{12}(1\alpha )+\frac{3}{4}\alpha =\frac{1}{2},$$ and $`\alpha `$ becomes $`1/4`$ which coincides with that of Fig.7 (c). In the case of the phase zero, we also obtain $`\alpha =1/4`$(Fig.7(d)). Contrary to above facts, there also exist some states on branches stable against a weak perturbation. Let us consider states on branch A–D in Fig.3. There exists a state $`\mathrm{}011011011\mathrm{}`$ on the branch and we obtain $`\mathrm{}011020011\mathrm{}`$ by a weak perturbation. Against the perturbation, perturbed state is stable and it is still on the branch A–D. We call this type of state metastable state. Moreover, let us consider a state $`\mathrm{}111120111120\mathrm{}`$ corresponding to F in Fig.3. We can obtain $`\mathrm{}201120111120\mathrm{}`$ by a weak perturbation, and its flow does not change. Therefore, metastable states exist not only on branch A–D but also on the phase one branch. In these cases, the effect of the site “2” does not spread over the entire system during the time evolution. Thus some part of the multiple branches are metastable states, which can be observed in long time simulations starting from random initial conditions. ## 5 Concluding discussions In this paper, some multi-value generalizations of rule–184 family are studied by using max–plus representation. We obtain evolution equations in a conserved form giving $`(L+1)`$-value CA. We have proved that EBCA1 model is a generalization of SlS model, and it gives multiple states and stop-and-go waves. We consider that a suitable traffic model gives those two phenomena at least. If we observe a real traffic data like Fig.1, existence of multiple states is clear. In a real traffic, cars are considered to be always perturbed by some traffic effects. Therefore, we consider existence of stable branches against a weak perturbation is also desirable for a model. Moreover, stop-and-go waves are often observed in a real congested traffic. Considering these points, EBCA1 is the most suitable model among the models of this paper. Next, we give a presumption on a fundamental diagram of real data. In Fig.1, we see the multiple state clearly in an acceleration lane. In a driving lane, since the average speed of cars is slower than the acceleration lane, flow tends to be stable and the branch is not clear. In the acceleration lane, drivers move faster than cars in the driving lane, then over-dense free flow will be likely to occur and we can see the multiple state around a critical density in the diagram. Moreover, from the stability analysis we find that there are metastable states even in the congested state in the EBCA1 model. Although we cannot clearly see the fact due to fluctuation of data in Fig.1, we expect that there exists both a metastable “weak jam”(line B–C) and a “strong jam”(line D–E) in a real highway traffic. Finally, we point out future problems. We have investigated mainly about models with $`L=2`$. About multiple branches due to larger $`L`$, only EBCA2 gives clear branches and we can easily discuss their stability. Though other models can also make multiple branches, transformation of perturbed states is complicated and its analysis is difficult. We should solve this difficulty because we consider EBCA1 model is more suitable for a traffic model. Moreover, two-dimensional extension of these models and its application to a real traffic network is one of the important future problems. Acknowledgment The authors are grateful to Professors, Shinji Takesue, Hisao Hayakawa, Makoto Kikuchi, Shinichi Tadaki, Yuki Sugiyama for fruitful discussions and helpful comments. This work is partially supported by Grant-in-Aid from the Ministry of Education, Science and Culture. Appendix In this appendix, we show that BCA model with $`L=2`$ can be interpreted as a two-lane model of rule–184 CA. Let us call one of the two lanes A-lane and the other B-lane. Both lanes are divided into discrete sites as shown in Fig.A1. Assume $`A_j^t`$ and $`B_j^t`$ denote number of cars at site $`j`$ and time $`t`$ in A-lane and B-lane respectively. Since all sites can hold only one car at most, $`A_j^t`$ and $`B_j^t`$ are always 0 or 1. A car at site $`j`$ in A-lane moves according to the following rule: * If site $`j+1`$ in A-lane is empty, the car moves to that site. * If site $`j+1`$ in A-lane is not empty and site $`j`$ and $`j+1`$ in B-lane are both empty, the car can move to site $`j+1`$ in B-lane. * Otherwise, the car stays at site $`j`$ in A-lane. As for a car in B-lane, symmetrical rule with respect to lane symbol is applied. The above rule can be interpreted as follows: Every car move in its own lane prior to the other lane (rule (a)). However, if a car can not move in its own lane and only if there are no traffic in the other lane, it changes a lane and move forward in the other lane (rule (b)). Since cars move independently in each lane according to rule–184 if rule (b) is omitted, the above rule can be interpreted as a two-lane model based on rule–184 with lane-changing effect. We can express the rule by a couple of evolution equations as follows: $`A_j^{t+1}`$ $`=`$ $`A_j^t+\mathrm{min}(A_{j1}^t,1A_j^t)\mathrm{min}(A_j^t,1A_{j+1}^t)`$ (10) $`+\mathrm{min}(1A_{j1}^t,1A_j^t,B_{j1}^t,B_j^t)\mathrm{min}(A_j^t,A_{j+1}^t,1B_j^t,1B_{j+1}^t),`$ $`B_j^{t+1}`$ $`=`$ $`B_j^t+\mathrm{min}(B_{j1}^t,1B_j^t)\mathrm{min}(B_j^t,1B_{j+1}^t)`$ (11) $`+\mathrm{min}(1B_{j1}^t,1B_j^t,A_{j1}^t,A_j^t)\mathrm{min}(B_j^t,B_{j+1}^t,1A_j^t,1A_{j+1}^t).`$ The last two terms in both equations express lane-changing effect. If those terms are omitted, both equations become independent each other and are equivalent to rule–184 CA (BCA model with $`L=1`$). If $`U_j^t`$ is defined by $$U_j^t=A_j^t+B_j^t,$$ (12) it denotes a sum of cars of both lanes. Since $`A_j^t`$ and $`B_j^t`$ is always 0 or 1, we can derive an evolution equation on $`U_j^t`$ from (10) and (11) as follows: $$U_j^{t+1}=U_j^t+\mathrm{min}(U_{j1}^t,\mathrm{\hspace{0.17em}2}U_j^t)\mathrm{min}(U_j^t,\mathrm{\hspace{0.17em}2}U_{j+1}^t).$$ (13) This equation is equivalent to BCA with $`L=2`$. Figure Captions 1. Observed data of flow (vehicles$`/`$5min.) versus density (vehicles$`/`$Km) on Tomei expressway. This diagram was taken by Japan Highway Public Corporation. (a) driving lane, (b) acceleration lane. 2. Fundamental diagrams of multi-value CA models with $`L=2`$ and $`K=30`$: (a) multi-value QS, (b) multi-value SlS, (c) EBCA2, (d)EBCA1 models. 3. Schematic fundamental diagram of EBCA1 model. 4. Time evolution of flow in EBCA1 model. (a) $`K=60`$, (b) $`K=240`$, (c) power spectrum $`I`$ versus frequency $`f`$ of traffic flow for $`K=240`$. 5. Fundamental diagram of EBCA1 model with $`L=1`$. 6. Fundamental diagram of EBCA2 model with $`L=7`$. 7. Instability of uniform flow due to a weak perturbation. (a) multi-value SlS, (b) EBCA2, (c) EBCA1 models. (d) shows an instability due to a strong perturbation in EBCA1 model. In all figures, black, dark gray, light gray squares denote a value 2, 1, 0 respectively. 8. Interpretation of BCA with $`L=2`$ as a model of coupled single lanes.
warning/0002/cond-mat0002274.html	ar5iv	text	# Effect of Magnetic field on the Pseudogap Phenomena in High-𝑇_c Cuprates ## 1 Introduction Since the discovery of high-temperature (High-$`T_\mathrm{c}`$) superconductivity by Bednortz and M$`\ddot{\mathrm{u}}`$ller, $`^{\text{?}\text{)}}`$ the anomalous normal state properties have been studied for many years from the various points of view. In particular, the pseudogap phenomena in under-doped cuprates have been recognized as one of the most important issues. There are enormous studies for the issue from both experimental and theoretical points of view. However, the complete understanding still remains to be obtained. The pseudogap phenomena mean the suppression of the spectral weight near the Fermi energy without any long range order. They are universal phenomena observed in various compounds of under-doped cuprates. Various experiments such as nuclear magnetic resonance (NMR), $`^{\text{?}\text{)}}`$ optical conductivity, $`^{\text{?}\text{)}}`$ transport, $`^{\text{?}\text{)}}`$ angle-resolved photo-emission spectroscopy (ARPES), $`^{\text{?}\text{)}}`$ tunneling spectroscopy, $`^{\text{?}\text{)}}`$ electronic specific heat, $`^{\text{?}\text{)}}`$ and so on have indicated the existence of the pseudogap in the normal state High-$`T_\mathrm{c}`$ cuprates from optimally-doped to under-doped region. In particular, NMR measurements of $`1/T_1T`$ have shown the existence of the pseudogap in the spin excitation channel from early years. $`^{\text{?)}}`$ In the previous paper, we have explained the pseudogap phenomena as a precursor of the strong coupling superconductivity. $`^{\text{?}\text{)}}`$ Since the effective Fermi energy $`\epsilon _\mathrm{F}`$ is renormalized by the electron-electron correlation, the ratio $`T_\mathrm{c}/\epsilon _\mathrm{F}`$ increases in the strongly correlated electron systems. The ratio indicates the strength of the superconducting coupling. Therefore, the strong coupling superconductivity has a general importance for the superconductivity in the strongly correlated electron systems. Moreover, it is natural to consider the strong coupling superconductivity in High-$`T_\mathrm{c}`$ cuprates because of the high critical temperature $`T_\mathrm{c}`$ itself. The strong coupling superconductivity necessarily leads to the strong thermal superconducting fluctuations. Such strong fluctuations in the quasi-two dimensional systems have serious effects on the electronic state and give rise to the pseudogap phenomena. $`^{\text{?)}}`$ Actually, various experiments have indicated the close relationship between the pseudogap phenomena and the superconductivity. In particular, ARPES have directly shown the pseudogap in the one-particle spectral weight $`^{\text{?)}}`$ and suggested its close relevance and continuity to the superconducting gap. $`^{\text{?}\text{)}}`$ Other scenarios have been theoretically proposed for the pseudogap phenomena. In the resonating valence bond (RVB) theory, there are two distinct excitations, spinon and holon. The pseudogap is described as a spinon pairing (so-called ’spin gap’). $`^{\text{?}\text{)}}`$ The magnetic scenarios based on the anti-ferromagnetic or SDW gap formation or their precursor have been proposed by various authors. $`^{\text{?}\text{)}}`$ Furthermore, the pairing scenarios as a precursor of the superconductivity are classified into several types. The phase fluctuation scenarios have been proposed by Emery and Kivelson $`^{\text{?}\text{)}}`$ and calculated by various authors. $`^{\text{?}\text{)}}`$ The scenario based on the strong coupling superconductivity has been proposed $`^{\text{?}\text{)}}`$ on the basis of the famous Nozi$`\stackrel{`}{\mathrm{e}}`$res and Schmitt-Rink formalism. $`^{\text{?}\text{)}}`$ The Nozi$`\stackrel{`}{\mathrm{e}}`$res and Schmitt-Rink formalism is justified in the low density limit. However, the nearly half-filled lattice system should be regarded as a rather high density case. Therefore, the Nozi$`\stackrel{`}{\mathrm{e}}`$res and Schmitt-Rink formalism cannot be applied to the pseudogap phenomena in High-$`T_\mathrm{c}`$ cuprates. Our scenario is based on the strong coupling superconductivity, but is different from the Nozi$`\stackrel{`}{\mathrm{e}}`$res and Schmitt-Rink formalism. We think of the pseudogap as the gap brought about by the resonance scattering $`^{\text{?}\text{}\text{?}\text{)}}`$ with the strong superconducting fluctuations. The strong superconducting fluctuations necessarily exist in case of the strong coupling superconductivity in the quasi-two dimensional systems. We have shown that the pseudogap phenomena are naturally understood on the basis of the resonance scattering scenario. $`^{\text{?)}}`$ Recently, the magnetic field effects on the NMR spin-lattice relaxation rate $`1/T_1`$ have been measured and discussed by several groups to determine the correct scenarios for the pseudogap phenomena .$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ The experimental results are interpreted as follows. The magnetic field effects cannot be observed in under-doped cuprates in which the strong pseudogap phenomena occurs in the wide temperature region. $`^{\text{?, ?)}}`$ In particular, the onset temperature $`T^{}`$ does not vary. On the other hand, the magnetic field effects are visible from optimally-doped to slightly over-doped cuprates in which only the weak pseudogap phenomena are observed in the narrow temperature region. $`^{\text{?, }\text{?}\text{)}}`$ The observed magnetic field dependences are explained by the conventional weak coupling theory. $`^{\text{?}\text{, ?)}}`$ However, for the under-doped cuprates, we have no theoretical explanation of the magnetic field effects on the pseudogap phenomena. In this paper, we point out that the magnetic field effects are naturally and continuously understood from under-doped to over-doped cuprates on the basis of our resonance scattering scenario. In particular, there is an interpretation that regards the experimental results for under-doped systems as a negative evidence for the pairing scenario. $`^{\text{?)}}`$ Our results conclude that this interpretation is inappropriate. It is generally considered that the superconducting fluctuations are remarkably influenced by the magnetic field, while the effects of the magnetic field on the spin-fluctuations are considered to be small. Therefore, the experimental results may be interpreted as an evidence for the magnetic scenario for the pseudogap. The misinterpretation is caused by the loss of the understanding for the strong coupling superconductivity. Therefore, we give an explanation for the magnetic field effects on the pseudogap phenomena on the basis of the strong coupling superconductivity. Actually, the experimental results including their doping dependence rather support our scenario for the pseudogap phenomena. $`^{\text{?)}}`$ This paper is constructed as follows. In §2, we give a model Hamiltonian and explain the theoretical framework adopted in this paper. In §3, we explicitly calculate the single particle self-energy $`\Sigma ^\mathrm{R}(𝒌,\omega )`$, density of states $`\rho (\epsilon )`$, NMR spin-lattice relaxation rate $`1/T_1T`$ and their magnetic field dependences. In §4, we discuss the transport phenomena in the pseudogap phase. In §5, we summarize the obtained results and give discussions. ## 2 Theoretical Framework In this section, we describe the theoretical framework in this paper. We calculate the magnetic field effects on the pseudogap phenomena by using the same formalism as is used in our previous paper. $`^{\text{?)}}`$ Therefore, in the first subsection §2.1, we briefly explain the formalism and show the outline of the obtained results in ref.8. In the second subsection §2.2, we introduce the magnetic field effects thorough the Landau quantization and give a rough estimate for the effects. Hereafter, we adopt the unit $`\mathrm{}=c=k_\mathrm{B}=1`$. ### 2.1 Pseudogap phenomena under zero magnetic field We adopt the following two-dimensional model Hamiltonian which has a $`d_{x^2y^2}`$-wave superconducting ground state, with High-$`T_\mathrm{c}`$ cuprates in mind. $`H={\displaystyle \underset{𝒌,s}{}}\epsilon _𝒌c_{𝒌,s}^{}c_{𝒌,s}+{\displaystyle \underset{𝒌,𝒌^{\mathbf{}},𝒒}{}}V_{𝒌𝒒/2,𝒌^{\mathbf{}}𝒒/2}c_{𝒒𝒌^{\mathbf{}},}^{}c_{𝒌^{\mathbf{}},}^{}c_{𝒌,}c_{𝒒𝒌,},`$ (2.1) where $`V_{𝒌,𝒌^{\mathbf{}}}`$ is the $`d_{x^2y^2}`$-wave separable pairing interaction, $`V_{𝒌,𝒌^{\mathbf{}}}=g\phi _𝒌\phi _𝒌^{\mathbf{}},`$ (2.2) $`\phi _𝒌=\mathrm{cos}k_x\mathrm{cos}k_y.`$ (2.3) Here, $`g`$ is negative. $`\phi _𝒌`$ is the $`d_{x^2y^2}`$-wave form factor. We consider the dispersion $`\epsilon _𝒌`$ given by the tight-binding model for a square lattice including the nearest- and next-nearest-neighbor hopping $`t`$, $`t^{}`$, respectively, $`\epsilon _𝒌=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)+4t^{}\mathrm{cos}k_x\mathrm{cos}k_y\mu .`$ (2.4) We fix the lattice constant $`a=1`$. We adopt $`t=0.5\mathrm{eV}`$ and $`t^{}=0.45t`$. These parameters well reproduce the Fermi surface of the typical High-$`T_\mathrm{c}`$ cuprates, $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_{6+\delta }`$ and $`\mathrm{Bi}_2\mathrm{Sr}_2\mathrm{CaCu}_2\mathrm{O}_{8+\delta }`$. We choose the chemical potential $`\mu `$ so that the filling $`n=0.9`$. This filling corresponds to the hole doping $`\delta =0.1`$. The Fermi surface is shown in Fig.1. In reality, the origin of the pairing interaction should be considered to be the anti-ferromagnetic spin fluctuations. $`^{\text{?}\text{}\text{?}\text{)}}`$ The spin fluctuations not only cause the pairing interaction but also affect the electronic state. $`^{\text{?, ?, }\text{?}\text{}\text{?}\text{)}}`$ There are studies dealing with the pairing correlation arising from the spin fluctuations on the basis of the fluctuation exchange (FLEX) approximation. $`^{\text{?}\text{}\text{?}\text{)}}`$ However, we do not introduce these effects because these details do not seriously affect the pseudogap phenomena as a precursor of the $`d_{x^2y^2}`$-wave superconductivity. There is a feedback effect on the pairing interaction arising from the pseudogap. The pseudogap affects the low frequency component of the spin fluctuations. However, the pairing interaction is mainly caused by the high frequency component of the spin fluctuations. Therefore, we can neglect the feedback effect on the pairing interaction and start from the model with an attractive interaction. The superconducting fluctuations are expressed by the T-matrix (Fig.2), $`t(𝒒,\mathrm{i}\mathrm{\Omega }_n)^1`$ $`=`$ $`g^1+\chi _0(𝒒,\mathrm{i}\mathrm{\Omega }_n),`$ (2.5) $`\chi _0(𝒒,\mathrm{i}\mathrm{\Omega }_n)`$ $`=`$ $`T{\displaystyle \underset{𝒌^{\mathbf{}},\omega _m}{}}G(𝒌^{\mathbf{}},\mathrm{i}\omega _m)G(𝒒𝒌^{\mathbf{}},\mathrm{i}\mathrm{\Omega }_n\mathrm{i}\omega _m)\phi _{𝒌^{\mathbf{}}𝒒/2}^2.`$ (2.6) Here, $`\omega _m=2\pi (m+\frac{1}{2})T`$ and $`\mathrm{\Omega }_n=2\pi nT`$ are the fermionic and bosonic Matsubara frequencies, respectively. Here, the scattering vertex arising from the superconducting fluctuations, $`\mathrm{\Gamma }(𝒌,𝒒𝒌:𝒌^{\mathbf{}},𝒒𝒌^{\mathbf{}}:\mathrm{i}\mathrm{\Omega }_n)`$ is factorized into $`\mathrm{\Gamma }(𝒌,𝒒𝒌:𝒌^{\mathbf{}},𝒒𝒌^{\mathbf{}}:\mathrm{i}\mathrm{\Omega }_n)=\phi _{𝒌𝒒/2}t(𝒒,\mathrm{i}\mathrm{\Omega }_n)\phi _{𝒌^{\mathbf{}}𝒒/2}`$. The form factor $`\phi _𝒌`$ in the above expression gives rise to the $`d_{x^2y^2}`$-wave shape of the pseudogap. When $`1+g\chi _0(\mathrm{𝟎},0)=0`$, $`t(\mathrm{𝟎},0)`$ diverges and the superconductivity occurs. This is the famous Thouless criterion which is equivalent to that of the BCS theory in the weak coupling limit. $`^{\text{?)}}`$ Analytically continued T-matrix $`t(𝒒,\mathrm{\Omega })`$ can be regarded as a propagator of the fluctuating Cooper pairs. Here, we are interested in the normal state near the superconducting critical point, where the superconducting fluctuations are enhanced. There, $`1+g\chi _0(\mathrm{𝟎},0)`$ is small and $`t(𝒒,\mathrm{\Omega })`$ is strongly enhanced around $`𝒒=\mathrm{\Omega }=0`$. Even in the weak coupling limit, there are corrections on the various quantities due to the superconducting fluctuations. They are well known as the Aslamazov-Larkin term (AL term) $`^{\text{?}\text{)}}`$ and the Maki-Thompson term (MT term). $`^{\text{?}\text{)}}`$ These terms are the corrections on the two-body correlation function. On the other hand, the superconducting fluctuations more seriously affect the one-particle electronic states in the strong or intermediate coupling region. The superconducting fluctuations give rise to the pseudogap phenomena. The weak correction on the density of states (DOS correction term) has been discussed for High-$`T_\mathrm{c}`$ cuprates within the weak coupling theory. $`^{\text{?, }\text{?}\text{)}}`$ Our calculation corresponds to an extension of these weak coupling theories to the strong coupling ones. Because $`t(𝒒,\mathrm{\Omega })`$ is strongly enhanced around $`𝒒=\mathrm{\Omega }=0`$, its main contribution to the single particle self-energy $`\Sigma (𝒌,\omega )`$ originates from the vicinity of $`𝒒=\mathrm{\Omega }=0`$. Therefore, we expand $`t^1(𝒒,\mathrm{\Omega })`$ in the vicinity of $`𝒒=\mathrm{\Omega }=0`$. This expansion corresponds to the time-dependent-Ginzburg-Landau (TDGL) expansion. $`gt^1(𝒒,\mathrm{\Omega })=t_0+b𝒒^2(a_1+\mathrm{i}a_2)\mathrm{\Omega }.`$ (2.7) The properties of the TDGL parameters are discussed in detail in our previous paper. $`^{\text{?)}}`$ The outline is the following. As is described above, $`t_0=1+g\chi _0(\mathrm{𝟎},0)`$ is $`0`$ at the critical point and is sufficiently small in the vicinity of the critical point. The parameter $`b`$ is generally related to the coherence length $`\xi _0`$, $`b\xi _0^2`$. The parameter $`a_2`$ express the time scale of the fluctuations. Roughly speaking, the parameters $`a_2`$ and $`b`$ are described as $`a_2\rho _\mathrm{d}(0)/T_\mathrm{c}`$ $`b\rho _\mathrm{d}(0)/T_\mathrm{c}^2.`$ (2.8) Here, we have defined the effective density of states for the $`d_{x^2y^2}`$-wave symmetry, $`\rho _\mathrm{d}(\epsilon )=_𝒌\rho _𝒌(\epsilon )\phi _𝒌^2`$, where, $`\rho _𝒌(\epsilon )`$ is the one-particle spectral weight $`\rho _𝒌(\epsilon )=A_𝒌(\epsilon )=\frac{1}{\pi }\mathrm{Im}G^\mathrm{R}(𝒌,\epsilon )`$. It should be noticed that $`\rho _\mathrm{d}(\epsilon )`$ is more sensitive to the pseudogap formation rather than the usual density of states $`\rho (\epsilon )=_𝒌\rho _𝒌(\epsilon )`$. Because of the high critical temperature $`T_\mathrm{c}`$ and the renormalization effect by the pseudogap, both $`a_2`$ and $`b`$ are strongly reduced in the strong coupling superconductivity. These features indicate that the scattering vertex due to the superconducting fluctuations is strongly enhanced. Although the T-matrix calculation used in this paper does not include the renormalization effect, these behaviors are obtained qualitatively. On the other hand, $`a_1`$ is not so reduced by the strong coupling superconductivity. Especially, High-$`T_\mathrm{c}`$ cuprates have a comparatively large value of $`a_1`$ because of their strong particle-hole asymmetry. Therefore, we cannot neglect $`a_1`$, although it is usually neglected in the weak coupling theories. $`^{\text{?}\text{)}}`$ In the T-matrix calculation, we estimate the TDGL parameters by using the non-interacting Green function $`G^{(0)\mathrm{R}}(𝒌,\omega )=(\omega \epsilon _𝒌+\mathrm{i}\delta )^1`$ (Fig.2(a)). This estimation corresponds to the Gauss approximation for the superconducting fluctuations. The critical temperature is determined in the mean field level, $`T_\mathrm{c}=T_{\mathrm{MF}}`$. As will be discussed in §5, the self-consistent T-matrix calculation includes the somewhat critical fluctuations (Figs.2(b) and 11). However, fundamental features do not change. As is minutely described in the previous paper ,$`^{\text{?)}}`$ the resonance scattering by the strong superconducting fluctuations gives rise to the pseudogap on both the one-particle spectral weight and the density of states. In this paper we describe these features by using the T-matrix calculation (Fig.3(a)). Although we pointed out the several important points in the self-consistent T-matrix calculation (Fig.3(b)), the pseudogap is properly described within the T-matrix calculation. $`^{\text{?)}}`$ In the T-matrix calculation, the self-energy is given by $`\Sigma (𝒌,\mathrm{i}\omega _n)=T{\displaystyle \underset{𝒒,\mathrm{i}\mathrm{\Omega }_m}{}}t(𝒒,\mathrm{i}\mathrm{\Omega }_m)G^{(0)}(𝒒𝒌,\mathrm{i}\mathrm{\Omega }_m\mathrm{i}\omega _n)\phi _{𝒌𝒒/2}^2.`$ (2.9) After the analytic continuation, we obtain $`\Sigma ^\mathrm{R}(𝒌,\omega )`$ $`=`$ $`{\displaystyle \underset{𝒒}{}}{\displaystyle }{\displaystyle \frac{\mathrm{d}\mathrm{\Omega }}{\pi }}[b(\mathrm{\Omega })\mathrm{Im}t(𝒒,\mathrm{\Omega })G^{(0)\mathrm{A}}(𝒒𝒌,\mathrm{\Omega }\omega )`$ (2.10) $`f(\mathrm{\Omega })t(𝒒,\mathrm{\Omega }+\omega )\mathrm{Im}G^{(0)\mathrm{R}}(𝒒𝒌,\mathrm{\Omega })]\phi _{𝒌𝒒/2}^2,`$ $`\mathrm{Im}\Sigma ^\mathrm{R}(𝒌,\omega )`$ $`=`$ $`{\displaystyle \underset{𝒒}{}}{\displaystyle \frac{\mathrm{d}\mathrm{\Omega }}{\pi }[b(\mathrm{\Omega }+\omega )+f(\mathrm{\Omega })]\mathrm{Im}t(𝒒,\mathrm{\Omega }+\omega )\mathrm{Im}G^{(0)\mathrm{R}}(𝒒𝒌,\mathrm{\Omega })\phi _{𝒌𝒒/2}^2}.`$ (2.11) Here,$`f(\mathrm{\Omega })`$ and $`b(\mathrm{\Omega })`$ are the Fermi and Bose distribution functions, respectively. We explicitly estimate the TDGL expansion parameters and numerically calculate the single particle self-energy $`\Sigma ^\mathrm{R}(𝒌,\omega )`$. Typical features of the single-particle self-energy are shown in Fig.4. Here, we exclude the trivial Hartree-Fock term shown in Fig.3(c). It is notable that the real part of the self-energy has a positive slope at the low frequency, and the imaginary part has a sharp peak in its absolute value there. Both features are anomalous compared with the conventional Fermi liquid theory. These anomalous features of the single particle self-energy should be regarded as the effects of the resonance scattering. Of course, the resonance scattering becomes strong as the superconductivity becomes strong coupling one and the systems approach the critical point. These features lead to the pseudogap. The corresponding one-particle spectral weight $`A(𝒌,\omega )`$ and density of states $`\rho (\omega )`$ are shown in Fig.5. Both show the gap structure above $`T_\mathrm{c}`$. It should be noticed that the pseudogap is the characteristics of the strong coupling superconductivity and does not occur in the weak coupling limit. In Fig.4 and 5, we have included the magnetic field effect described in the next subsection. The pseudogap reduces the critical temperature $`T_\mathrm{c}`$. The reduction becomes more remarkable as the coupling constant increases. Therefore, although the mean field critical temperature $`T_{\mathrm{MF}}`$ remarkably increases with the coupling constant $`\|g\|`$, $`T_\mathrm{c}`$ does not vary so much. $`^{\text{?)}}`$ Here, $`T_\mathrm{c}`$ is scaled by the effective Fermi energy $`\epsilon _\mathrm{F}`$. By considering the fact, it is naturally understood that $`T_\mathrm{c}`$ decreases with the doping quantity in the under-doped region. As the doping quantity decreases, the system approaches to the Mott insulator. Therefore, it should be considered that the renormalization effect for the effective Fermi energy $`\epsilon _\mathrm{F}`$ is enhanced with decreasing the doping quantity. Since $`T_\mathrm{c}/\epsilon _\mathrm{F}`$ is almost independent of the coupling constant $`\|g\|`$ in the strong coupling region, $`T_\mathrm{c}`$ decreases with $`\epsilon _\mathrm{F}`$ in the under-doped region. Thus, our theory naturally and appropriately explains the pseudogap phenomena in High-$`T_\mathrm{c}`$ cuprates. ### 2.2 Magnetic field effects on the pseudogap phenomena In this subsection, we introduce the magnetic field effect. In this paper, we consider the magnetic field applied along the c-axis, $`B\stackrel{}{c}`$. The main effect of the magnetic fields is the Landau level quantization for the superconducting fluctuations. It corresponds to the quantization of the orbital motion of the fluctuating Cooper pairs. The quantization is expressed by the replacement of the quadratic term of the momentum as $`𝒒^24\mathrm{e}B(n+\frac{1}{2})`$$`^{\text{?)}}`$ The Landau quantization has the following two important effects. One is the Landau degeneracy which generally enhances the fluctuations. The Landau degeneracy reduces the dimensionality of the fluctuations. The other is the suppression of the superconductivity which weakens the pseudogap. The distance to the critical point increases as $`t_0t_0+2b\mathrm{e}B`$. This corresponds to the energy level of the Lowest Landau level. When considering at the fixed temperature, the dominant effect is the latter. We can see that the characteristic magnetic field $`B_{\mathrm{ch}}`$ for the pseudogap phenomena is scaled by the quantity $`t_0/b`$, that is $`B_{\mathrm{ch}}t_0/b`$$`^{\text{?}\text{)}}`$ The ratio $`b/t_0`$ corresponds to the square of the GL correlation length $`\xi _{\mathrm{GL}}`$ for the superconducting fluctuations, that is, $`b/t_0=\xi _{\mathrm{GL}}^2`$. The magnetic field effects are scaled by the quantity $`B\xi _{\mathrm{GL}}^2`$. Thus, the pseudogap is affected by the magnetic field according to the magnetic flux penetrating the correlated area $`\xi _{\mathrm{GL}}^2`$. As we mentioned above, the parameter $`b`$ is small in the strong coupling case. Moreover, the fact that the pseudogap phenomena take place in the wide temperature region means that $`t_0`$ is large near the pseudogap onset temperature $`T^{}`$. As a result, the characteristic magnetic field $`B_{\mathrm{ch}}`$ is large, especially near $`T^{}`$. In other words, the magnetic field effects are remarkably small in case of the strong coupling superconductivity. Especially, the onset temperature $`T^{}`$ does not vary. Since $`\xi _{\mathrm{GL}}`$ diverges at the critical temperature $`T_\mathrm{c}`$, the magnetic field effects are sure to appear near $`T_\mathrm{c}`$. However, the region in which the effects appear is remarkably small. On the other hand, in the relatively weak coupling case, the magnetic field dependence is large and $`T^{}`$ may vary. These features well explain the results of the high field NMR measurements including their doping dependence. $`^{\text{?, ?, ?, ?)}}`$ Here, we have neglected the Zeeman coupling term. Although the Zeeman coupling term plays an important role at the low temperature in superconducting state, $`^{\text{?}\text{)}}`$ it has only higher order correction in the fluctuating region. This fact can be simply understood as follows. The lowest order correction of the Zeeman coupling term on the superconducting fluctuations is the second order and described as $`4\mathrm{e}B(n+\frac{1}{2})4\mathrm{e}B(n+\frac{1}{2})+\frac{8}{v^2}(\mu B)^2`$. Here, $`v`$ is a mean value of the quasi-particle velocity on the Fermi surface. $`\mu =g\mu _\mathrm{B}/2`$ is the magnetic moment of the electrons. Here, $`g`$ is the $`g`$-value and $`\mu _\mathrm{B}`$ is the Bohr magneton. Thus, the Zeeman coupling term slightly weaken the superconducting fluctuations. However, it has only higher order effect with respect to the magnetic field compared with the Landau quantization. Therefore, the effect of the Zeeman coupling term is extraordinary small in the weak coupling limit since the typical magnetic field is small. Also for High-$`T_\mathrm{c}`$ cuprates, it is higher order and remarkably small compared with the effect of the Landau quantization in the magnetic field of the experimentally relevant order. Actually, the magnetic field adopted in this paper is the order of $`\mathrm{e}B10^2`$ in our unit. That corresponds to $`B10\mathrm{T}\mathrm{e}\mathrm{s}\mathrm{l}\mathrm{a}`$. In this case, the effect of the Zeeman coupling term is higher order than that of Landau quantization as $`10^2`$. Thus, it is justified to neglect the Zeeman coupling term. Of course, we cannot neglect the Zeeman coupling term under the extraordinary high magnetic field in the strong coupling limit. However, such an extreme situation is not realistic. When the magnetic field is applied perpendicular to the c-axis $`B\stackrel{}{c}`$, effect of the Zeeman coupling term is relatively important because the coherence length along the c-axis $`\xi _\mathrm{c}`$ is small. ## 3 Magnetic Field Dependence of the NMR Spin-Lattice Relaxation Rate, $`1/T_1T`$ In this section, we actually calculate the NMR spin-lattice relaxation rate $`1/T_1T`$ under the magnetic field with the recent high field NMR measurements in mind. We calculate $`1/T_1T`$ by using the general expression, $`1/T_1T={\displaystyle \underset{𝒒}{}}\|A(𝒒)\|^2[{\displaystyle \frac{1}{\omega }}\mathrm{Im}\chi _\mathrm{s}^\mathrm{R}(𝒒,\omega )_{\omega 0}].`$ (3.1) Here, we neglect the momentum dependence of the hyperfine coupling $`A(𝒒)`$ for simplicity, which does not affect the magnetic field dependence of $`1/T_1T`$. We calculate the spin susceptibility $`\chi _\mathrm{s}^\mathrm{R}(𝒒,\omega )`$ which corresponds to the two-body correlation function shown in Fig.6. Here, the solid lines are the renormalized Green function $`G^\mathrm{R}(𝒌,\omega )=(\omega \epsilon _𝒌\Sigma ^\mathrm{R}(𝒌,\omega ))^1`$. The self-energy $`\Sigma ^\mathrm{R}(𝒌,\omega )`$ is calculated by using the T-matrix approximation as we described before. The effects of the superconducting fluctuations are included in the self-energy. In calculating the self-energy $`\Sigma ^\mathrm{R}(𝒌,\omega )`$, we linearize the dispersion relation as $`\epsilon _{𝒌𝒒}=\epsilon _𝒌𝒗_𝒌𝒒`$. This linearization is justified because the only small region around $`𝒒=0`$ contributes to the self-energy. We replace the quadratic term as $`𝒒^24\mathrm{e}B(n+\frac{1}{2})`$. This process corresponds to the Landau level quantization for the superconducting fluctuations. From eq.(3.1), $`1/T_1T`$ is expressed as, $`1/T_1T`$ $`=`$ $`{\displaystyle \underset{𝒌,𝒒}{}}{\displaystyle \frac{\mathrm{d}\omega }{\pi }f^{}(\omega )\mathrm{Im}G^\mathrm{R}(𝒌,\omega )\mathrm{Im}G^\mathrm{R}(𝒌\mathbf{+}𝒒,\omega )}`$ (3.2) $`=`$ $`\pi {\displaystyle d\omega f^{}(\omega )\rho (\omega )^2}.`$ Here, $`f^{}(\omega )`$ is the differential of the Fermi distribution function. This expression is reduced to the well-known expression $`1/T_1T=\pi \rho (0)^2`$ at $`T=0`$. After all, we calculate the decrease of $`1/T_1T`$ by the suppression of the density of states. Generally speaking, we can consider the Aslamazov-Larkin term (AL term) and the Maki-Thompson term (MT term) as corrections by the fluctuations on the two-body correlation function. $`^{\text{?, ?)}}`$ However, the AL term dose not exist in calculating the spin susceptibility $`\chi _\mathrm{s}^\mathrm{R}(𝒒,\omega )`$. We can understand this fact by considering the spin index for the spin singlet pairing. $`^{\text{?, ?)}}`$ The contribution from the MT term is small in case of the d-wave pairing, and suppressed by the slight elastic scattering. $`^{\text{?)}}`$ Therefore, we have only to calculate the decrease of $`1/T_1T`$ by the pseudogap as the effect of the superconducting fluctuations. Of course, we have to take account of the anti-ferromagnetic spin-fluctuations in order to describe the whole temperature dependence of $`1/T_1T`$. $`1/T_1T`$ increases owing to the anti-ferromagnetic spin fluctuations in the normal phase ($`T>T^{}`$), and decreases owing to the superconducting fluctuations in the pseudogap phase ($`T^{}>T>T_\mathrm{c}`$). Generally speaking, the magnetic field is considered to have a great effect on the superconducting fluctuations, while the effect on the spin-fluctuations is comparatively small. Because we pay attention to the magnetic field dependence in this paper, we have only to calculate the decrease of $`1/T_1T`$ due to the superconducting fluctuations and its magnetic field dependence. Actually, the misinterpretation for the experimental results is caused by the loss of the understanding of the magnetic field effect on the superconducting fluctuations in case of the strong coupling superconductivity. Our calculation gives a clear understanding about the magnetic field dependence of the pseudogap phenomena. Even if the effect of the exchange enhancement is taken into account, the results for the magnetic field effect do not change, qualitatively. At the last of this section, we definitely calculate the effect of the exchange enhancement within the random phase approximation (RPA). Qualitatively, the same results are given there. The calculated results are shown in Fig.7, 8 and 9. In all figures, the magnetic field is varied as $`4\mathrm{e}B=0.01,0.02,0.05`$ and $`0.1`$ in our unit. The horizontal axis is the temperature scaled by the zero-field critical temperature $`T_{\mathrm{c0}}`$. All results can be understood by considering the characteristic magnetic field, as we have mentioned in the previous section, $`B_{\mathrm{ch}}t_0/b=\xi _{\mathrm{GL}}^2.`$ (3.3) The results for the relatively weak coupling case $`g=0.5`$ are shown in Fig.7. In this case, only the weak pseudogap is observed in the narrow temperature region. We consider that this case corresponds to the slightly over-doped or optimally-doped cuprates. The magnetic field dependence of $`1/T_1T`$ is clearly observed and $`T^{}`$ varies. These behaviors are consistent with the NMR experiments in the slightly over-doped cuprates. $`^{\text{?)}}`$ In case of $`g=0.8`$, the magnetic field dependence is small since the parameter $`b`$ decreases(Fig.8). In particular, $`1/T_1T`$ is almost independent of the magnetic field near the onset temperature $`T^{}`$ where the parameter $`t_0`$ is large. On the other hand, the magnetic field dependence can be observed in the vicinity of the critical temperature $`T_\mathrm{c}`$, since $`t_0`$ is small there. We can see the different magnetic field dependences of the density of states according to the distance to the critical point (Fig.9). The magnetic field effect is visible in the density of states just above $`T_\mathrm{c}`$ (Fig.9(a)). The density of states at the low energy are recovered with increasing the magnetic field. On the other hand, the effect is almost invisible when the temperature is apart form $`T_\mathrm{c}`$ (Fig.9(b)). The results for the considerably strong coupling case $`g=1.0`$ is shown in Fig.10. In this case, the strong pseudogap anomaly exists in the wide temperature region. The magnetic field effects become still smaller. The magnetic field dependence is narrowly observed in the vicinity of $`T_\mathrm{c}`$. We consider that these strong coupling cases correspond to the under-doped cuprates. These behaviors well explain the experimental results in the under-doped cuprates. $`^{\text{?)}}`$ The weak effect in the vicinity of $`T_\mathrm{c}`$ is also observed in the experimental results. Thus, the interpretation of the experimental results as a negative evidence for the pairing scenario is inappropriate. The strength of the superconducting coupling is indicated by the ratio $`T_\mathrm{c}/\epsilon _\mathrm{F}`$. The ratio increases due to the mass renormalization by the electron-electron correlation. It should be considered that the mass-renormalization is enhanced with decreasing the doping quantity. The attractive interaction becomes strong at the same time, since the anti-ferromagnetic spin fluctuations are enhanced. Therefore, it is expected that the superconductivity becomes the strong coupling as the doping quantity decreases. Thus, the strength of the superconducting coupling naturally changes with the doping in accordance with our expectation. It should be noticed that the change of the magnetic field effects is continuous from weak to strong coupling. In other words, the calculated results explain the NMR measurements continuously and entirely from over-doped to under-doped cuprates. Therefore, the recent high field NMR measurements including their doping dependence are regarded as an affirmative evidence for the pairing scenario. In order to confirm the effect of the Landau degeneracy to enhance the fluctuations, we show the Fig.11. In Fig.11, the horizontal axis is scaled by the critical temperature under the magnetic field $`T_{\mathrm{cH}}`$. By keeping the distance to the critical point, we can remove the effect of the suppression of the superconductivity. Therefore, we can see the effect of the Landau degeneracy. The results show that $`1/T_1T`$ decreases with increasing the magnetic field. It is because of the Landau degeneracy. The Landau degeneracy enhances the superconducting fluctuations and make the pseudogap stronger. Then, $`1/T_1T`$ is still more reduced. Therefore, even in the rather weak pseudogap case, the pseudogap may be observed clearly under the high magnetic field. In other words, the magnetic field makes the pseudogap visible in more over-doped cuprates. At the last of this section, we consider the effect of the exchange enhancement. The exchange enhancement is taken into account within the random phase approximation (RPA). The basic results about the magnetic field effects are not changed. However, it is definitely shown that the peak of $`1/T_1T`$ ($`T=T^{}`$) does not change in the strong coupling case, while the peak changes in the weak coupling case. The dynamical spin susceptibility $`\chi _{\mathrm{RPA}}(𝒌,\omega )`$ calculated by RPA is expressed as follows. $`\chi _{\mathrm{RPA}}^\mathrm{R}(𝒒,\omega )`$ $`=`$ $`{\displaystyle \frac{\chi _0^\mathrm{R}(𝒒,\omega )}{1U\chi _0^\mathrm{R}(𝒒,\omega )}},`$ (3.4) $`\chi _0(𝒒,\mathrm{i}\omega _n)`$ $`=`$ $`T{\displaystyle \underset{𝒌,\omega _m}{}}G(𝒌,\mathrm{i}\omega _m)G(𝒌+𝒒,\mathrm{i}\mathrm{\Omega }_m+\mathrm{i}\omega _n).`$ We fix $`U=1.5`$ afterward. $`1/T_1T`$ is calculated by eq.(3.1). Here, we take into account the momentum dependence of the hyperfine coupling $`\|A(𝒒)\|^2=\frac{1}{2}[\{A_1+2B(\mathrm{cos}(q_x)+\mathrm{cos}(q_y))\}^2+\{A_2+2B(\mathrm{cos}(q_x)+\mathrm{cos}(q_y))\}^2]`$. The hyperfine coupling constants $`A_1,A_2`$ and $`B`$ is evaluated as $`A_1=0.84B`$ and $`A_2=4B`$$`^{\text{?}\text{)}}`$ The following results are not affected by the choice of the parameters, qualitatively. The calculated results are shown in Figs.12 and 13. In the high temperature region, $`1/T_1T`$ is enhanced owing to the exchange enhancement. Near the critical temperature, $`1/T_1T`$ is reduced owing to the superconducting fluctuations. As a result, $`1/T_1T`$ shows its peak at $`T=T^{}`$ above $`T_\mathrm{c}`$. It is a well-known pseudogap phenomenon in NMR measurements. In the weak coupling case $`g=0.5`$ (Fig.12), the magnetic field effect is clearly observed. $`T^{}`$ is lowered by the magnetic field. On the other hand, in the relatively strong coupling case $`g=0.8`$ (Fig.13), the magnetic filed effect is remarkably small. $`T^{}`$ is not changed by the magnetic field. $`1/T_1T`$ shows the magnetic filed dependence only in the vicinity of the critical temperature $`T_\mathrm{c}`$. These features are the same as those derived by the calculation without the effect of the exchange enhancement. The results for the spin-echo decay rate $`1/T_{2G}`$ are shown in the inset of Figs.12 and 13. $`1/T_{2G}`$ is calculated by the following expression. $`1/T_{2G}^2={\displaystyle \underset{𝒒}{}}[\|A_{}(𝒒)\|^2\mathrm{Re}\chi _\mathrm{s}^\mathrm{R}(𝒒,0)]^2[{\displaystyle \underset{𝒒}{}}\|A_{}(𝒒)\|^2\mathrm{Re}\chi _\mathrm{s}^\mathrm{R}(𝒒,0)]^2.`$ (3.6) Here, the dynamical spin susceptibility is calculated by RPA, and $`\|A_{}(𝒒)\|^2=\{A_2+2B(\mathrm{cos}(q_x)+\mathrm{cos}(q_y))\}^2`$$`^{\text{?)}}`$ $`1/T_{2G}`$ shows the pseudogap phenomena. However, the effect of the pseudogap on $`1/T_{2G}`$ is weaker than that on $`1/T_1T`$. The pseudogap appears in the narrower temperature region. $`1/T_{2G}`$ shows its peak below the pseudogap onset temperature $`T^{}`$ in $`1/T_1T`$. These results are consistent with the experimental results. $`^{\text{?}\text{)}}`$ These results indicate that the effects of the pseudogap are weak on the real part of the spin susceptibility rather than on the imaginary part at the low frequency. The dissipation (imaginary part) directly reflects the low energy density of state. However, the static property (real part) does not necessarily so. In other wards, the pseudogap suppresses the weight of the spin susceptibility at low frequency. However, the effect on the total weight is rather small. In particular, the $`d`$-wave pseudogap only weakly affects the real part near the anti-ferromagnetic wave vector $`𝒒=(\pi ,\pi )`$. The momentum dependence of the hyperfine coupling $`A_{}(𝒒)`$ still more weaken the effect of the pseudogap on $`1/T_{2G}`$. The above features are in common with the superconducting state. $`^{\text{?}\text{)}}`$ That is natural because the pseudogap and the superconducting gap have the same $`d_{x^2y^2}`$-wave form. The magnetic field dependence of $`1/T_{2G}`$ has the same features as those of $`1/T_1T`$. ## 4 Transport in the Pseudogap Phase In the previous sections, we have paid attention to the NMR spin-lattice relaxation rate $`1/T_1T`$ and its magnetic field dependence in the pseudogap phase. Besides that, the pseudogap affects many other measurements. These effects may be understood by considering the suppression of the one-particle spectral weight and that of the low frequency anti-ferromagnetic spin fluctuations. $`^{\text{?)}}`$ In particular, the transport phenomena are enough interesting to be discussed here, because they reflect the characteristic momentum dependences of High-$`T_\mathrm{c}`$ cuprates and the relationship between the spin fluctuations and the superconducting fluctuations. The following qualitative discussion deserves to be described here, because there is no explicit calculation based on our understanding. First, we describe how the transport phenomena in under-doped cuprates are understood in the normal phase ($`T>T^{}`$). They are anomalous at a glance. However, we can understand them by considering the magnetic interaction caused by the anti-ferromagnetic spin fluctuations. $`^{\text{?, }\text{?}\text{)}}`$ The momentum dependence of the lifetime of quasi-particles is important to understand the transport properties. The momentum dependent lifetime is due to the scattering by the anti-ferromagnetic spin fluctuations. $`^{\text{?, ?)}}`$ ’Hot spot’ means the part of the Fermi surface in which $`\epsilon _𝒌=\epsilon _{𝒌+𝑸}`$. Here, $`𝑸`$ is a anti-ferromagnetic wave vector $`𝑸=(\pi ,\pi )`$. At ’hot spot’, quasi-particles suffer an immediate scattering by the anti-ferromagnetic spin fluctuations at $`𝒒=𝑸`$. ’Cold spot’ is the area on the Fermi surface far from ’hot spot’. There, quasi-particles do not suffer the immediate scattering. Therefore, the lifetime is long at ’cold spot’ and short at ’hot spot’. This momentum dependent lifetime is a general property of the systems with anti-ferromagnetic spin fluctuations. $`^{\text{?}\text{)}}`$ For High-$`T_\mathrm{c}`$ cuprates, ’hot spot’ is located near $`(\pi ,0)`$ and its symmetric points. ’Cold spot’ is located near $`(\pi /2,\pi /2)`$. At the same time, the pseudogap is large at ’hot spot’ and small at ’cold spot’ because of its $`d_{x^2y^2}`$-wave shape. ’Hot spot’ does not contribute to the in-plane conductivity, because the conductivity is almost determined by the most easily flowing quasiparticles. The in-plane conductivity is mainly determined by ’cold spot’. The quasiparticles at ’cold spot’ are sure to have the $`T^2`$-damping rate at the low temperature limit which is consistent with the conventional Fermi liquid theory. However, they have the $`T`$-linear damping rate above the relatively low crossover temperature ($`T>T_{\mathrm{cr}}`$). It is because of the low energy magnetic excitations. The transformation of the Fermi surface which leads to a form more appropriate to the nesting reduces the crossover temperature $`T_{\mathrm{cr}}`$. The transformation itself is due to the anti-ferromagnetic spin fluctuations. $`^{\text{?)}}`$ As a result, the in-plane resistivity shows a $`T`$-linear law in the normal phase ($`T>T^{}>T_{\mathrm{cr}}`$). It should be noticed that $`T`$-linear resistivity is not due to the Curie-Weiss law $`\chi _\mathrm{s}(𝑸)1/(T+\theta )`$, or $`1/T_1T1/(T+\theta )`$. The calculations inappropriately treating the momentum dependent lifetime attribute the $`T`$-linear resistivity to the Curie-Weiss law. $`^{\text{?, }\text{?}\text{)}}`$ For example, the approximate relation between the in-plane resistivity $`\rho _{\mathrm{ab}}`$ and $`1/T_1T`$, $`\rho _{\mathrm{ab}}T^2/(T_1T)`$ is derived. $`^{\text{?)}}`$ If appropriately considering ’hot spot’ and ’cold spot’, the $`T`$-linear resistivity is realized more generally, but in more high temperature region. $`^{\text{?)}}`$ This generality is important to understand the $`T`$-linear in-plane resistivity in the pseudogap phase. The other important character of High-$`T_\mathrm{c}`$ cuprates is a momentum dependence of the interlayer hopping matrix element $`t_{}(𝒌)`$. The band calculation has shown that the dispersion along the c-axis is large at ’hot spot’ and is nearly $`0`$ at ’cold spot’. $`^{\text{?}\text{)}}`$ $`t_{}(𝒌)`$ is approximately expressed as $`^{\text{?, }\text{?}\text{)}}`$ $`t_{}(𝒌)({\displaystyle \frac{\mathrm{cos}k_x\mathrm{cos}k_y}{2}})^2.`$ (4.1) Since the quasiparticle velocity along the c-axis is nearly $`0`$ at ’cold spot’, ’cold spot’ does not contribute to the c-axis conductivity. On the other hand, the contribution from ’hot spot’ is reduced by the short lifetime, in spite of the large velocity along the c-axis. As a result, the c-axis transport becomes incoherent. Thus, we can understand the coherent in-plane conductance and the incoherent c-axis conductance at the same time. $`^{\text{?)}}`$ The momentum dependent lifetime enhances the Hall coefficient $`R=\sigma _{xy}/\sigma _{xx}^2H`$$`^{\text{?)}}`$ However, the vertex correction plays a more important role for the Hall coefficient. $`^{\text{?)}}`$ It is because of the momentum derivative of the total current $`J_\nu `$ in the general formula given by Kohno and Yamada. $`^{\text{?}\text{)}}`$ The Hall coefficient is strongly enhanced by the vertex correction. In the conventional metals, the vertex correction gives only an constant factor arising from the Umklapp scattering and has no significant effect. $`^{\text{?}\text{)}}`$ The significance of the vertex correction is also due to the anti-ferromagnetic spin fluctuations. The vertex correction is not so important for the longitudinal conductivity even in the systems with spin fluctuations. $`^{\text{?)}}`$ Here, we consider the transport phenomena in the pseudogap phase. The main effects of the pseudogap on the transport phenomena are the following two points. One is the pseudogap itself. The other is the suppression of the anti-ferromagnetic spin fluctuations. Because of the singlet pairing correlation, the low frequency part of the anti-ferromagnetic spin fluctuations is expected to be suppressed. Indeed, NMR measurements show the suppression of $`(T_1T)^1`$ and $`(T_{2G})^1`$$`^{\text{?, ?, }\text{?}\text{)}}`$ The low frequency part of the anti-ferromagnetic spin fluctuations causes the quasiparticle damping. Therefore, the quasiparticle damping due to the anti-ferromagnetic spin fluctuations is immediately affected by the pseudogap. As is shown in §2, the imaginary part of the self-energy due to the superconducting fluctuations, $`\mathrm{Im}\Sigma ^\mathrm{R}(𝒌,\omega )`$ is remarkably large in the pseudogap phase. The large imaginary part leads to the pseudogap near $`(\pi ,0)`$. Therefore, the contribution to the conductivity from the quasiparticles near $`(\pi ,0)`$ is remarkably suppressed by the pseudogap itself. However, quasiparticles at ’hot spot’ do not contribute to the in-plane conductivity from the beginning. The in-plane conductivity is determined by the contribution from ’cold spot’. Therefore, the pseudogap itself is not important for the in-plane conductivity. The suppression of the anti-ferromagnetic spin fluctuations slightly reduces the scattering at ’cold spot’. Quasiparticles at ’cold spot’ are not affected by the strong scattering due to the spin fluctuations at $`𝒒=𝑸`$. Therefore, the effect of the suppression of the anti-ferromagnetic spin fluctuations is small at ’cold spot’ rather than at ’hot spot’. As a result, the in-plane resistivity slightly deviates downward and keep the $`T`$-linear law. $`^{\text{?)}}`$ This behavior is observed in many under-doped compounds and the downward deviation coincides with the pseudogap. $`^{\text{?)}}`$ Thus, the $`T`$-linear resistivity generally appears owing to the low frequency magnetic excitations. The $`T`$-linear law in the pseudogap phase can not be understood by the phenomenological relation, $`\rho _{\mathrm{ab}}T^2\chi _\mathrm{s}(𝑸)`$ or $`\rho _{\mathrm{ab}}T^2/(T_1T)`$. On the other hand, the pseudogap itself has more drastic effect on the c-axis conductivity. The c-axis conductance is determined by the contribution from the vicinity of $`(\pi ,0)`$. Quasiparticles near $`(\pi ,0)`$ decrease the contribution to the c-axis conductivity owing to the pseudogap. The pseudogap is large there. Therefore, the c-axis conductivity is remarkably suppressed by the pseudogap. This effect is confirmed within the formalism in this paper. We calculate the c-axis resistivity by using the Kubo formula and neglecting the vertex correction. The c-axis conductivity is expressed as $`\sigma _\mathrm{c}(T)=de^2{\displaystyle \underset{𝒌}{}}t_{}^2(𝒌){\displaystyle \frac{\mathrm{d}\omega }{\pi }(f^{}(\omega ))\mathrm{Im}G^\mathrm{R}(𝒌,\omega )\mathrm{Im}G^\mathrm{R}(𝒌,\omega )}.`$ (4.2) Here, $`d`$ is the interlayer distance. We normalize the conductivity by the constant factor $`de^2`$. The calculated result is shown in Fig.14. The c-axis resistivity shows the semi-conductive behavior near the critical point $`T_\mathrm{c}`$. It is because the scattering due to the superconducting fluctuations becomes remarkable with approaching the critical temperature. Thus, we can understand the drastic increase of the c-axis resistivity in the pseudogap phase, while the in-plane resistivity changes only a little. Because of the momentum dependence of the hopping matrix element, the c-axis transport reflects the electronic state near $`(\pi ,0)`$. Therefore, we can see the pseudogap by observing the c-axis optical conductivity, $`^{\text{?)}}`$ while the in-plane optical conductivity does not clearly indicate the pseudogap. For the Hall conductivity, both two points play an important role. Because of the suppression of the spin fluctuations, the enhancement of the Hall coefficient due to the spin fluctuations is reduced. Moreover, since the vertex correction is large around ’hot spot’, $`^{\text{?)}}`$ the pseudogap itself affects the vertex correction. The pseudogap opens at ’hot spot’ and reduces the effects of the vertex correction. Due to the two effects, the Hall coefficient is reduced and shows its peak in the pseudogap phase. These results well explain the observed transport phenomena in the pseudogap phase. $`^{\text{?)}}`$ Thus, the transport phenomena in the pseudogap phase is naturally understood by considering the $`d_{x^2y^2}`$-wave pseudogap. ## 5 Summary and Discussion In this paper, we have shown that the pairing scenario based on the strong coupling superconductivity well explains the effects of the magnetic field on the pseudogap phenomena in High-$`T_\mathrm{c}`$ cuprates. We have shown in the previous paper $`^{\text{?)}}`$ that the pseudogap phenomena are properly described as a precursor of the superconductivity under the reasonable conditions. In this paper, we have used the same formalism for calculating the single-particle self-energy and introduce the magnetic field effects thorough the Landau level quantization. We explicitly calculated the NMR spin-lattice relaxation rate $`1/T_1T`$ to compare the obtained results with the results of the recent high field NMR measurements. The dominant effect of the Landau quantization is the suppression of the superconductivity, while the Landau degeneracy itself enhances the superconducting fluctuations. From the simple discussion, we can see that the characteristic magnetic field is scaled as $`B_{\mathrm{ch}}t_0/b=\xi _{\mathrm{GL}}^2`$. Actually, the calculated results support this behavior. In the relatively weak coupling case, the weak pseudogap is observed in the narrow temperature region. Then, the characteristic magnetic field is small and the magnetic field effects are visible. On the other hand, in the strong coupling case where the pseudogap is observed in the wide temperature region, the characteristic magnetic field is large. In particular, it is remarkably large near the onset temperature $`T^{}`$. Therefore $`T^{}`$ is almost independent of the magnetic field. The magnetic field effects are visible only in the vicinity of the critical temperature $`T_\mathrm{c}`$. It should be noticed that the characteristic magnetic field $`B_{\mathrm{ch}}`$ near $`T^{}`$ is different from that near $`T_\mathrm{c}`$. When the pseudogap phenomena take place at $`T^{}`$, the superconducting correlation length $`\xi _{\mathrm{GL}}`$ is still short. Therefore, the pseudogap is not so affected by the magnetic field near $`T^{}`$. By considering that the effective Fermi energy $`\epsilon _\mathrm{F}`$ decreases and the attractive interaction increases with decreasing the doping quantity, the calculated results well explain the high field NMR measurements including their doping dependence. The explanation is continuous from over-doped to under-doped cuprates. There is an interpretation that the magnetic field independence of the pseudogap phenomena in under-doped cuprates is an evidence denying the pairing scenarios for the pseudogap. $`^{\text{?)}}`$ However, the pairing scenario based on the strong coupling superconductivity naturally explains the experiments including their doping dependence. Moreover, the continuous understanding in the phase diagram rather support the pairing scenario. In the pseudogap phase, the self-energy correction due to the superconducting fluctuations is a common mechanism in reducing the density of states and $`1/T_1T`$. Because the pseudogap phenomena continuously take place from slightly over-doped to under-doped cuprates, their magnetic field dependences should be continuously understood. The pseudogap becomes strong as the doping quantity decreases. The magnetic field dependence of the weak pseudogap case can be understood within the conventional weak coupling theory for the superconducting fluctuations. $`^{\text{?, ?)}}`$ Our theory is an extension of the theory. This fact indicates the correctness of our description for the pseudogap phenomena in under-doped cuprates based on the strong coupling superconductivity. On the other hand, it is not clear whether the magnetic origin may be consistent with the magnetic field dependence, especially in the weak pseudogap case. It is because the magnetic exchange coupling $`J`$ is the order of $`J1000\mathrm{K}`$ and the applied magnetic field is the order of $`\mu _\mathrm{B}B10\mathrm{K}`$. Moreover, we have discussed the transport phenomena in the pseudogap phase. Generally speaking, the transport phenomena in the normal phase are explained by the effects of the anti-ferromagnetic spin fluctuations. We have shown that the transport coefficients in the pseudogap phase are naturally understood by considering the characteristic momentum dependences of both the spin- and the superconducting fluctuations in addition to the momentum dependence of the c-axis transfer matrix. The c-axis conductivity is mainly determined by the region near $`(\pi ,0)`$. Therefore, the c-axis transport directly reflects the pseudogap. We have definitely shown the remarkable increase of the c-axis resistivity in the pseudogap phase. Here, we give a brief discussion on the self-consistent calculation. In the self-consistent calculation the pseudogap is described in a similar way. The fundamental picture does not change also in the self-consistent calculation, although the renormalization effects on the TDGL parameters exist. However, the self-consistent T-matrix calculation is one of the methods introducing the criticality of the superconducting fluctuations. This effect corresponds to the forth order term in the Ginzburg-Landau description. The forth order term in the Ginzburg-Landau action is expressed by the diagram shown in Fig.15. This term indicates the repulsive interaction between the fluctuating Cooper pairs (that is the mode coupling term). The effect of this term is included in the self-consistent calculation at least in the level of the Hartree-Fock approximation. Thus, the criticality of the superconducting fluctuations is introduced. The criticality makes the magnetic field dependence still smaller. To put it in detail, in the self-consistent calculation, $`t_0`$ depends on the magnetic field. As the magnetic field suppresses the pseudogap, $`t_0(B)`$ is reduced. Therefore, the distance to the superconductivity $`t_0(B)+2b\mathrm{e}B`$ varies more slowly than $`t_0(0)+2b\mathrm{e}B`$. Thus, the magnetic field dependence is reduced by the criticality. Anyway, the strong coupling superconductivity is the essential factor for the magnetic field independence, as we described in this paper. The existence of the wide critical region is a result of the strong coupling superconductivity. The more systematic measurements of the magnetic field dependences in the various doping rate will be an important verification to determine the origin of the pseudogap in High-$`T_\mathrm{c}`$ cuprates. ## Acknowledgements The authors are grateful to Mr. Koikegami for fruitful discussions. The authors are grateful to Dr. G-q. Zheng for teaching us the experimental results. Numerical computation in this work was partly carried out at the Yukawa Institute Computer Facility. The present work was partly supported by a Grant-In-Aid for Scientific Research from the Ministry of Education, Science, Sports and Culture, Japan. One of the authors (Y.Y) has been supported by a Research Fellowships of the Japan Society for the Promotion of Science for Young Scientists.
warning/0002/hep-ph0002245.html	ar5iv	text	# IMPACT OF PHYSICAL PRINCIPLES AT VERY HIGH ENERGY SCALES ON THE SUPERPARTICLE MASS SPECTRUM ## I Introduction The Minimal Supersymmetric Standard Model (MSSM) is a well-motivated extension of the Standard Model (SM) that includes broken supersymmetry (SUSY) at the weak scale. To construct the MSSM, one postulates: * the gauge group and the matter content of the SM, where the various fields are replaced by superfields: $`\widehat{Q}_i=\left(\begin{array}{c}\widehat{u}_i\\ \widehat{d}_i\end{array}\right),\widehat{L}_i=\left(\begin{array}{c}\widehat{\nu }_i\\ \widehat{e}_i\end{array}\right),\widehat{U}_i^c,\widehat{D}_i^c,\widehat{E}_i^c,`$ where $`i=1,2,3`$ corresponding to the various generations; * an extended Higgs sector that includes two different $`SU(2)`$ doublet Higgs superfields $`\widehat{H}_u(\mathrm{𝟐})=\left(\begin{array}{c}\widehat{h}_u^+\\ \widehat{h}_u^0\end{array}\right),\mathrm{and}\widehat{H}_d(\overline{\mathrm{𝟐}})=\left(\begin{array}{c}\widehat{h}_d^{}\\ \widehat{h}_d^0\end{array}\right);`$ to allow superpotential Yukawa couplings (and hence, masses) for both up and down type fermions. The introduction of two doublets of Higgsinos is also just right to cancel the chiral anomaly from the gauginos. * an $`R`$-parity conserving renormalizable superpotential, <sup></sup><sup></sup>Our sign convention for the $`\mu `$-term is defined by the chargino and neutralino mass matrices given in Eqs. (33) and (34) of the review by X. Tata, Ref. . $`\widehat{f}=\mu \widehat{H}_u^a\widehat{H}_{da}+f_uϵ_{ab}\widehat{Q}^a\widehat{H}_u^b\widehat{U}^c+f_d\widehat{Q}^a\widehat{H}_{da}\widehat{D}^c+f_e\widehat{L}^a\widehat{H}_{da}\widehat{E}^c+\mathrm{},`$ where $`ϵ_{ab}`$ is the completely antisymmetric $`SU(2)`$ tensor with $`ϵ_{12}=1`$, and the ellipses refer to Yukawa couplings for the second and third generations; soft supersymmetry breaking (SSB) terms consistent with Lorentz invariance and SM gauge invariance, $`_{soft}={\displaystyle \underset{r}{}}m_r^2\|\varphi _r\|^2`$ $``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\lambda }{}}M_\lambda \overline{\lambda }_\alpha \lambda _\alpha +[B\mu \stackrel{~}{H}_d\stackrel{~}{H}_u+h.c]`$ $`+`$ $`[A_uf_uϵ\stackrel{~}{Q}\stackrel{~}{H}_u\stackrel{~}{u}_R^{}+A_df_d\stackrel{~}{Q}\stackrel{~}{H}_d\stackrel{~}{d}_R^{}+A_ef_e\stackrel{~}{L}\stackrel{~}{H}_d\stackrel{~}{e}_R^{}+\mathrm{}+h.c.],`$ where contraction over the $`SU(2)`$ indices is understood, and the ellipses again refer to terms of the second and third generation trilinear scalar couplings. In practice, because only third generation Yukawa couplings are sizeable, the $`A`$-parameters of just the third family are frequently relevant. Although we have not shown this explicitly, the Yukawa couplings and the $`A`$-parameters are, in general, (complex) matrices in generation space. The resulting framework then requires $`100`$ parameters beyond those of the SM, and hence is not very predictive. Since, the phenomenology that we consider is generally insensitive to inter-generation mixing of quarks and squarks, we assume that these matrices are diagonal. Furthermore, since we do not discuss $`CP`$ violating effects, we take the superpotential and soft SUSY breaking parameters to be real. Even so, a large number of additional parameters remains. Most of these occur in the SSB sector of the model, which simply reflects our ignorance of the mechanism of supersymmetry breaking. To gain predictivity, despite the lack of a compelling model of SUSY breaking, we must make additional simplifying assumptions about symmetries of interactions at energy scales not directly accessible to experiments, or postulate other physical principles that determine the origin of the soft SUSY breaking terms. The most popular model in which to embed the MSSM is the minimal supergravity model (mSUGRA) . In this model, supersymmetry is broken in a “hidden sector” which consists of fields which couple to the fields of the visible sector (the MSSM fields) only gravitationally. Within the framework of supergravity grand unification, the additional assumption that the vacuum expectation value ($`vev`$) of the gauge kinetic function does not break the unifying gauge symmetry leads to a common mass $`m_{1/2}`$ for all gauginos. In addition, it is usually assumed that there exists a common mass $`m_0`$ for all scalars and a common trilinear term $`A_0`$ for all soft SUSY breaking trilinear interactions. Universal soft SUSY breaking scalar masses are not, however, a consequence of the supergravity framework but an additional assumption. The (universal) soft parameters are assumed to be renormalized at some high scale $`M_XM_{GUT}M_{Planck}`$. These are assumed to have values comparable to the weak scale, $`M_{weak}`$, resulting in an elegant solution to the fine-tuning problem associated with the Higgs sector. Motivated by the apparently successful gauge coupling unification in the MSSM, the scale $`M_{GUT}2\times 10^{16}`$ GeV is usually adopted for the scale choice $`M_X`$. The resulting effective theory, valid at energy scales $`E<M_{GUT}`$, is then just the MSSM with soft SUSY breaking terms that unify at $`M_{GUT}`$. The soft SUSY breaking scalar and gaugino masses, the trilinear $`A`$ terms and in addition a bilinear soft term $`B`$, the gauge and Yukawa couplings and the supersymmetric $`\mu `$ term are all then evolved from $`M_{GUT}`$ to some scale $`MM_{weak}`$ using renormalization group equations (RGE). The large top quark Yukawa coupling causes the squared mass of $`H_u`$ to be driven to negative values, which signals the radiative breakdown of electroweak symmetry (REWSB); this then allows one to determine the value of $`\mu ^2`$ in terms of $`M_Z^2`$, possibly at the expense of some fine-tuning. Finally, it is customary to trade the parameter $`B`$ for $`\mathrm{tan}\beta `$, the ratio of Higgs field vacuum expectation values. The resulting weak scale spectrum of superpartners and their couplings can then be derived in terms of four continuous parameters plus one sign $$m_0,m_{1/2},A_0,\mathrm{tan}\beta \mathrm{and}sign(\mu ),$$ (3) in addition to the usual parameters of the standard model. This calculational procedure has been embedded into the event generator ISAJET thereby allowing detailed predictions for the collider events within this framework. The mSUGRA model, while highly predictive, rests upon a number of simplifying assumptions that are invalid in specific models of physics at energy scales $`M_{GUT}M_{Planck}`$. Thus, in the search for weak scale supersymmetry, the mSUGRA model may give misleading guidance as to the possible event signatures expected at high energy collider experiments. Indeed the literature is replete with models with non-universal soft SUSY breaking mass terms at the high scales. In this paper, we survey a variety of these models (as well as others that lead to universality) and comment on possible phenomenological implications, especially for high energy collider experiments. For the most part, we restrict our attention to models which reduce to the $`R`$-parity conserving MSSM at scales $`Q<M_{GUT}`$. The event generator ISAJET (versions $`>7.37`$) has recently been upgraded to accomodate supersymmetric models with non-universal soft SUSY breaking masses at the $`GUT`$ scale. To generate such models, the user must input the usual mSUGRA parameter set Eq. 3, but may in addition select one or several of the following options: $`NUSUG1:`$ $`M_1,M_2,M_3`$ $`NUSUG2:`$ $`A_t,A_b,A_\tau `$ $`NUSUG3:`$ $`m_{H_d},m_{H_u}`$ $`NUSUG4:`$ $`m_{Q_1},m_{D_1},m_{U_1},m_{L_1},m_{E_1}`$ $`NUSUG5:`$ $`m_{Q_3},m_{D_3},m_{U_3},m_{L_3},m_{E_3}.`$ If one or more of the $`NUSUGi`$ ($`i=15`$) inputs are selected, then the $`GUT`$ scale universal soft breaking masses are overwritten and a weak-scale MSSM mass spectrum is generated. ISAJET then computes the corresponding branching fractions and sparticle cross sections, so that specific theoretical predictions for $`GUT`$ scale SSB masses can be mapped onto explicit predictions for the high energy collider events expected to arise from these models. In addition, the ISAJET keyword $`SSBCSC`$ has been introduced in ISAJET versions $`7.50`$. Using $`SSBCSC`$, the user may choose any scale between the weak scale and the Planck scale at which to impose the above SSB boundary conditions. We illustrate its use in Sec. XI C where it is necessary to introduce boundary conditions at the string scale rather than at $`M_{GUT}`$. To facilitate the examination of these models by our experimental colleagues, we present here a survey of a number of well-motivated models which usually lead to non-universality of SSB parameters. Our survey is far from exhaustive, but is meant to present a flavor of the range of possibilities available for such models. For each model, we 1. present a short description of the physics, 2. delineate the parameter space, 3. indicate how, within the model framework, collider events may be generated using ISAJET, and 4. comment upon some of the general features of SUSY events expected at collider experiments. The models selected include the following: * $`SU(5)`$ grand unified models with universal soft SUSY breaking masses at scales higher than $`Q=M_{GUT}`$, * $`SU(5)`$ models where supersymmetry breaking occurs via non-singlet hidden sector superfields, * the MSSM plus an intermediate-scale right-handed neutrino which leads to see-saw neutrino masses, * models with extra $`D`$-term contributions to scalar masses that are generically present if the rank of the unifying gauge group exceeds 4, * minimal and general $`SO(10)`$ grand unified models with universal soft SUSY breaking masses at scales higher than $`Q=M_{GUT}`$, * grand unified models with group structure $`G_{GUT}\times G_H`$, where $`G_H`$ contains a hypercolor interaction used to solve the doublet-triplet splitting problem, * effective supersymmetry models which lead to multi-TeV range scalar masses for the first two generations, but sub-TeV masses for third generation scalars and gauginos, * anomaly-mediated SUSY breaking models (AMSB), where the hidden sector resides in different spacetime dimensions from the visible sector, * the minimal gaugino mediation model, * 4-dimensional string models with Calabi-Yao or orbifold compactifications, and * models inspired by $`M`$-theory with SUSY breaking by one or several moduli fields. Space limitations preclude us from detailed discussions of these models. Here, we sketch the physics behind each model, and provide the reader with selected references where further details may be found. While much, but by no means all, of the material presented may be found in the literature, our hope is that the form in which we have presented it will facilitate, or even spur, closer examination of alternatives to the mSUGRA and gauge-mediated SUSY breaking models. ## II $`SU(5)`$ grand unified model with the SSB universality scale higher than $`M_{GUT}`$ As a working assumption, the scale at which all the SSB parameters are generated, is usually taken to be $`M_{GUT}`$. If this scale is substantially higher than this (but smaller than the Planck scale), renormalization group (RG) evolution induces a non-universality at the GUT scale. The effect can be significant if large representations are present. Here, we assume that supersymmetric $`SU(5)`$ grand unification is valid at mass scales $`Q>M_{GUT}2\times 10^{16}`$ GeV, extending at most to the reduced Planck scale $`M_P2.4\times 10^{18}`$ GeV. Below $`Q=M_{GUT}`$, the $`SU(5)`$ model breaks down to the MSSM with the usual $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ gauge symmetry. This framework is well described in, for instance, the work of Polonsky and Pomarol. In the $`SU(5)`$ model, the $`\widehat{D}^c`$ and $`\widehat{L}`$ superfields are elements of a $`\overline{\mathrm{𝟓}}`$ superfield $`\widehat{\varphi }`$, while the $`\widehat{Q}`$, $`\widehat{U}^c`$ and $`\widehat{E}^c`$ superfields occur in the $`\mathrm{𝟏𝟎}`$ representation $`\widehat{\psi }`$. The Higgs sector is comprised of three super-multiplets: $`\widehat{\mathrm{\Sigma }}(\mathrm{𝟐𝟒})`$ which is responsible for breaking $`SU(5)`$, plus $`\widehat{}_1(\overline{\mathrm{𝟓}})`$ and $`\widehat{}_2(\mathrm{𝟓})`$ which contain the usual Higgs doublet superfields $`\widehat{H_d}`$ and $`\widehat{H_u}`$ respectively, which occur in the MSSM. The superpotential is given by, $`\widehat{f}=\mu _\mathrm{\Sigma }tr\widehat{\mathrm{\Sigma }}^2`$ $`+`$ $`{\displaystyle \frac{1}{6}}\lambda ^{}tr\widehat{\mathrm{\Sigma }}^3+\mu _H\widehat{}_1\widehat{}_2+\lambda \widehat{}_1\widehat{\mathrm{\Sigma }}\widehat{}_2`$ (4) $`+`$ $`{\displaystyle \frac{1}{4}}f_tϵ_{ijklm}\widehat{\psi }^{ij}\widehat{\psi }^{kl}\widehat{}_2^m+\sqrt{2}f_b\widehat{\psi }^{ij}\widehat{\varphi }_i\widehat{}_{1j},`$ (5) where a sum over families is understood. $`f_t`$ and $`f_b`$ are the top and bottom quark Yukawa couplings, $`\lambda `$ and $`\lambda ^{}`$ are GUT Higgs sector self couplings, and $`\mu _\mathrm{\Sigma }`$ and $`\mu _H`$ are superpotential Higgs mass terms. Supersymmetry breaking is parametrized by the soft supersymmetry breaking terms: $`_{soft}`$ $`=`$ $`m__1^2\|_1\|^2m__2^2\|_2\|^2m_\mathrm{\Sigma }^2tr\{\mathrm{\Sigma }^{}\mathrm{\Sigma }\}m_5^2\|\varphi \|^2m_{10}^2tr\{\psi ^{}\psi \}{\displaystyle \frac{1}{2}}M_5\overline{\lambda }_\alpha \lambda _\alpha `$ (6) $`+`$ $`[B_\mathrm{\Sigma }\mu _\mathrm{\Sigma }tr\mathrm{\Sigma }^2+{\displaystyle \frac{1}{6}}A_\lambda ^{}\lambda ^{}tr\mathrm{\Sigma }^3+B_H\mu _H_1_2+A_\lambda \lambda _1\mathrm{\Sigma }_2`$ (7) $`+`$ $`{\displaystyle \frac{1}{4}}A_tf_tϵ_{ijklm}\psi ^{ij}\psi ^{kl}_2^m+\sqrt{2}A_bf_b\psi ^{ij}\varphi _i_{1j}+h.c.]`$ (8) The various soft masses and gauge and Yukawa couplings evolve with energy according to the 15 renormalization group equations given in Appendix A of Ref. . Here, we modify them to correspond with the sign conventions in ISAJET : $`{\displaystyle \frac{dm_{10}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[3f_t^2(m__2^2+2m_{10}^2+A_t^2)+2f_b^2(m__1^2+m_{10}^2+m_5^2+A_b^2){\displaystyle \frac{72}{5}}g_G^2M_5^2\right],`$ (9) $`{\displaystyle \frac{dm_5^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[4f_b^2(m__1^2+m_{10}^2+m_5^2+A_b^2){\displaystyle \frac{48}{5}}g_G^2M_5^2\right],`$ (10) $`{\displaystyle \frac{dm__1^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[4f_b^2(m__1^2+m_{10}^2+m_5^2+A_b^2)+{\displaystyle \frac{24}{5}}\lambda ^2(m__1^2+m__2^2+m_\mathrm{\Sigma }^2+A_\lambda ^2){\displaystyle \frac{48}{5}}g_G^2M_5^2\right],`$ (11) $`{\displaystyle \frac{dm__2^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[3f_t^2(m__2^2+2m_{10}^2+A_t^2)+{\displaystyle \frac{24}{5}}\lambda ^2(m__1^2+m__2^2+m_\mathrm{\Sigma }^2+A_\lambda ^2){\displaystyle \frac{48}{5}}g_G^2M_5^2\right],`$ (12) $`{\displaystyle \frac{dm_\mathrm{\Sigma }^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[{\displaystyle \frac{21}{20}}\lambda ^2(3m_\mathrm{\Sigma }^2+A_\lambda ^{}^2)+\lambda ^2(m__1^2+m__2^2+m_\mathrm{\Sigma }^2+A_\lambda ^2)20g_G^2M_5^2\right],`$ (13) $`{\displaystyle \frac{dA_t}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[9A_tf_t^2+4A_bf_b^2+{\displaystyle \frac{24}{5}}A_\lambda \lambda ^2+{\displaystyle \frac{96}{5}}g_G^2M_5\right],`$ (14) $`{\displaystyle \frac{dA_b}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[10A_bf_b^2+3A_tf_t^2+{\displaystyle \frac{24}{5}}A_\lambda \lambda ^2+{\displaystyle \frac{84}{5}}g_G^2M_5\right],`$ (15) $`{\displaystyle \frac{dA_\lambda }{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[{\displaystyle \frac{21}{20}}A_\lambda ^{}\lambda ^2+3A_tf_t^2+4A_bf_b^2+{\displaystyle \frac{53}{5}}A_\lambda \lambda ^2+{\displaystyle \frac{98}{5}}g_G^2M_5\right],`$ (16) $`{\displaystyle \frac{dA_\lambda ^{}}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{8\pi ^2}}\left[{\displaystyle \frac{63}{20}}A_\lambda ^{}\lambda ^2+3A_\lambda \lambda ^2+30g_G^2M_5\right],`$ (17) $`{\displaystyle \frac{df_t}{dt}}`$ $`=`$ $`{\displaystyle \frac{f_t}{16\pi ^2}}\left[9f_t^2+4f_b^2+{\displaystyle \frac{24}{5}}\lambda ^2{\displaystyle \frac{96}{5}}g_G^2\right],`$ (18) $`{\displaystyle \frac{df_b}{dt}}`$ $`=`$ $`{\displaystyle \frac{f_b}{16\pi ^2}}\left[10f_b^2+3f_t^2+{\displaystyle \frac{24}{5}}\lambda ^2{\displaystyle \frac{84}{5}}g_G^2\right],`$ (19) $`{\displaystyle \frac{d\lambda }{dt}}`$ $`=`$ $`{\displaystyle \frac{\lambda }{16\pi ^2}}\left[{\displaystyle \frac{21}{20}}\lambda ^2+3f_t^2+4f_b^2+{\displaystyle \frac{53}{5}}\lambda ^2{\displaystyle \frac{98}{5}}g_G^2\right],`$ (20) $`{\displaystyle \frac{d\lambda ^{}}{dt}}`$ $`=`$ $`{\displaystyle \frac{\lambda ^{}}{16\pi ^2}}\left[{\displaystyle \frac{63}{20}}\lambda ^2+3\lambda ^230g_G^2\right],`$ (21) $`{\displaystyle \frac{d\alpha _G}{dt}}`$ $`=`$ $`3\alpha _G^2/2\pi ,`$ (22) $`{\displaystyle \frac{dM_5}{dt}}`$ $`=`$ $`3\alpha _GM_5/2\pi ,`$ (23) with $`t=\mathrm{log}Q`$. To generate the weak scale MSSM mass spectrum, one begins with the input parameters $$\alpha _{GUT},f_t,f_b,\lambda ,\lambda ^{}$$ (24) stipulated at $`Q=M_{GUT}`$, where $`f_b=f_\tau `$ is obtained from the corresponding mSUGRA model. The first three of these can be extracted, for instance, from ISASUGRA, versions $`7.44`$. The couplings $`\lambda (M_{GUT})`$ and $`\lambda ^{}(M_{GUT})`$ are additional inputs, where $`\lambda (M_{GUT})0.7`$ to make the triplet Higgsinos heavy enough to satisfy experimental bounds on the proton lifetime. The gauge and Yukawa couplings can be evolved via the RGEs to determine their values at $`Q=M_P`$. Assuming universality at $`M_P`$ (this maximizes the effects of non-universality at the GUT scale), one imposes $`m_{10}`$ $`=`$ $`m_5=m__1=m__2=m_\mathrm{\Sigma }m_0`$ (25) $`A_t`$ $`=`$ $`A_b=A_\lambda =A_\lambda ^{}A_0,`$ (26) and evolves all the soft masses from $`M_P`$ to $`M_{GUT}`$. The MSSM soft breaking masses at $`M_{GUT}`$ are specified via $`m_Q^2=m_U^2=m_E^2m_{10}^2,`$ (27) $`m_D^2=m_L^2m_5^2,`$ (28) $`m_{H_1}^2=m__1^2,m_{H_2}^2=m__2^2,`$ (29) which can serve as input to ISAJET via the $`NUSUGi`$ keywords. Since there is no splitting amongst the gaugino masses, the gaugino masses may be taken to be $`M_1=M_2=M_3m_{1/2}`$ where $`m_{1/2}`$ is stipulated most conveniently at the $`GUT`$ scale. To obtain correct Yukawa unification, it is crucial to start with the correct weak scale Yukawa couplings. To calculate the values of the Yukawa couplings at scale $`Q=M_Z`$, one begins with the pole masses $`m_b=4.9`$ GeV and $`m_\tau =1.784`$ GeV. One may calculate the corresponding running masses in the $`\overline{MS}`$ scheme, and evolve $`m_b`$ and $`m_\tau `$ up to $`M_Z`$ using 2-loop SM RGEs. At $`Q=M_Z`$, the SUSY loop corrections to $`m_b`$ and $`m_\tau `$ must be included; ISAJET versions $`>7.44`$ uses the approximate formulae of Pierce et al.. A similar procedure is used to calculate the top quark Yukawa coupling at scale $`Q=m_t`$. SUSY particle mass spectra consistent with constraints from collider searches and with unified $`b`$ and $`\tau `$ Yukawa couplings (to 5%) are then obtained (assuming universality of scalar masses at the scale $`M_P`$), but only for $`\mu <0`$ and $`30\mathrm{tan}\beta 50`$, where the allowed range is weakly sensitive to $`\alpha _s`$. To illustrate the extent of non-universality due to $`SU(5)`$ running of SSB masses between $`M_P`$ and $`M_{GUT}`$, we explicitly examine a typical case. The corresponding input parameters as well as the values of SSB parameters at $`M_{GUT}`$ are listed in Table LABEL:tsu5\_1. The GUT scale input parameters extracted from ISAJET for $`\mathrm{tan}\beta =35`$ are $`f_t=0.534`$ and $`f_b=f_\tau =0.271`$ for the top, bottom and tau Yukawa couplings. We also adopt $`\lambda =1.0`$ and $`\lambda ^{}=0.1`$ for the $`SU(5)`$ Higgs couplings and $`g_{GUT}=0.717`$ for the unified $`SU(5)`$ gauge coupling. At the Planck scale, we then take $`m_0=150`$ GeV and $`A_0=0`$ GeV, parameters that are analogous to $`m_0`$ and $`A_0`$ at the GUT scale in the mSUGRA model. We take $`m_{1/2}(M_{GUT})=200`$ GeV for the universal gaugino masses. The evolution of SUSY mass parameters in the minimal $`SU(5)`$ model between $`M_P`$ and $`M_{GUT}`$ is shown in Fig. 1, assuming universality at $`M_P`$. We see that the rather high value of $`\lambda `$ induces a large splitting $`m_5^2m_{10}^2>m__1^2,m__2^2`$. Likewise, the large value of $`f_t`$ is responsible for the splitting $`m__1^2>m__2^2`$ at $`M_{GUT}`$. The large $`t`$ and $`b`$ Yukawa couplings are also responsible for the split between third generation and the first two generation values of $`m_{10}`$ and $`m_5`$. It is interesting to notice that reasonable values of the free parameters can give $`100`$% deviations from universality at $`M_{GUT}`$. In the cases that we checked, it was typically the Higgs scalars that are split by the large amount from the other scalars, primarily because $`\lambda `$ is large: for acceptable solutions, the corresponding non-universality between matter scalar masses was typically $`1020`$%. In Table LABEL:tsu5\_2, we list the corresponding values of selected weak scale sparticle masses for both the $`SU(5)`$ case and mSUGRA. The shift in scalar masses in this case can be up to $`20\%`$, with the biggest shift occuring in the $`\stackrel{~}{\mathrm{}}_R`$ and $`\stackrel{~}{\tau }_1`$ masses. ## III $`SU(5)`$ models with non-universal gaugino masses. Since supergravity is not a renormalizable theory, in general we may expect a non-trivial gauge kinetic function, and hence the possiblity of non-vanishing gaugino masses if SUSY is broken. Expanding the gauge kinetic function as $`f_{ab}=\delta _{ab}+\widehat{\mathrm{\Phi }}_{ab}/M_{Planck}+\mathrm{}`$, where the fields $`\widehat{\mathrm{\Phi }}_{ab}`$ transform as left handed chiral superfields under supersymmetry transformations, and as the symmetric product of two adjoints under gauge symmetries, we parametrize the lowest order contribution to gaugino masses by, $$d^2\theta \widehat{f}^a\widehat{f}^b\frac{\widehat{\mathrm{\Phi }}_{ab}}{M_{\mathrm{Planck}}}+h.c.\frac{F_\mathrm{\Phi }_{ab}}{M_{\mathrm{Planck}}}\lambda ^a\lambda ^b+\mathrm{},$$ (30) where the $`\lambda ^a`$ are the gaugino fields, and $`F_\mathrm{\Phi }`$ is the auxillary field component of $`\widehat{\mathrm{\Phi }}`$ that acquires a SUSY breaking $`vev`$. If the fields $`F_\mathrm{\Phi }`$ which break supersymmetry are gauge singlets, universal gaugino masses result. However, in principle, the chiral superfield which communicates supersymmetry breaking to the gaugino fields can lie in any representation in the symmetric product of two adjoints, and so can lead to gaugino mass termsThe results of this section are not new, but in the interest of completeness we thought it fit to include a review of these models in this section. that (spontaneously) break the underlying gauge symmetry. We require, of course, that SM gauge symmetry is preserved. Non-universal gaugino masses have been previously considered by other authors. In the context of $`SU(5)`$ grand unification, $`F_\mathrm{\Phi }`$ belongs to an $`SU(5)`$ irreducible representation which appears in the symmetric product of two adjoints: $$(\mathrm{𝟐𝟒}\times \mathrm{𝟐𝟒})_{\mathrm{symmetric}}=\mathrm{𝟏}\mathrm{𝟐𝟒}\mathrm{𝟕𝟓}\mathrm{𝟐𝟎𝟎},$$ (31) where only $`\mathrm{𝟏}`$ yields universal masses. The relations amongst the various GUT scale gaugino masses have been worked out e.g. in Ref. . The relative $`GUT`$ scale $`SU(3)`$, $`SU(2)`$ and $`U(1)`$ gaugino masses $`M_3`$, $`M_2`$ and $`M_1`$ are listed in Table III along with the approximate masses after RGE evolution to $`QM_Z`$. Here, motivated by the measured values of the gauge couplings at LEP, we assume that the $`vev`$ of the SUSY-preserving scalar component of $`\widehat{\mathrm{\Phi }}`$ is neglible. Each of the three non-singlet models is as predictive as the canonical singlet case, and all are compatible with the unification of gauge couplings. These scenarios represent the predictive subset of the more general (and less predictive) case of an arbitrary superposition of these representations. The model parameters may be chosen to be, $$m_0,M_3^0,A_0,\mathrm{tan}\beta \mathrm{and}sign(\mu ),$$ (32) where $`M_i^0`$ is the $`SU(i)`$ gaugino mass at scale $`Q=M_{GUT}`$. $`M_2^0`$ and $`M_1^0`$ can then be calculated in terms of $`M_3^0`$ according to Table III. Sample spectra for each case are exhibited in Table LABEL:tnusug. The phenomenology of these models has recently been examined in Ref. , and the SUSY reach presented for Fermilab Tevatron upgrade options for a variety of discovery channels. The results were found to be model-dependent. In particular, in the 24 model, a large splitting between weak scale values of $`m_{\stackrel{~}{Z}_2},m_{\stackrel{~}{W}_1}`$ and $`m_{\stackrel{~}{Z}_1}`$ gave rise to large rates for events with isolated leptons, so that SUSY discovery should be easier in this case than in the mSUGRA model. A special feature of this model is the sizeable cross section for $`(Z\mathrm{}\overline{\mathrm{}})+jets+\overline{)}E_T`$ events. Indeed, for certain ranges of model parameters, SUSY discovery seemed to be possible only via this channel. In contrast, for the 75 and 200 models, $`m_{\stackrel{~}{Z}_2}`$, $`m_{\stackrel{~}{W}_1}`$ and $`m_{\stackrel{~}{Z}_1}`$ were all nearly degenerate, so that leptons arising from –ino decays were very soft and difficult to detect. Consequently, there was hardly any reach for SUSY in these models at the Tevatron via leptonic channels, and the best reach occurred typically in the $`\overline{)}E_T+jets`$ channels. ## IV The MSSM with a right handed neutrino Experimental evidence strongly indicates the existence of neutrino oscillations, and almost certainly neutrino mass. The favoured interpretation is $`\nu _\mu \nu _\tau `$ oscillations, with $`\mathrm{\Delta }m^210^2`$ eV<sup>2</sup> and near-maximal mixing. An attractive method for introducing neutrino mass into the MSSM is via the see-saw mechanism. In this case, one can introduce an additional chiral superfieldOur purpose here is to illustrate the effect of introducing singlet neutrino superfields on the SSB parameters and the SUSY spectrum. An explanation of the atmospheric neutrino data would, of course, require us to introduce more than one such superfield and also interactions that violate lepton flavour conservation, but as long as these have only small Yukawa couplings, their effect on the spectrum should be negligible. ($`\widehat{N}^c`$) which transforms as a gauge singlet (whose fermionic component is the left-handed anti-neutrino and scalar component is $`\stackrel{~}{\nu }_R^{}`$). A Majorana mass term for the right-handed neutrino is allowed and, because $`\nu _R`$ is a SM singlet, its mass may be large: $`M_N10^{10}10^{16}`$ GeV. When electroweak symmetry is broken, a Dirac neutrino mass $`m_Dm_{\mathrm{}}`$ is also induced via the usual Higgs mechanism. The resulting neutrino mass matrix must be diagonalized, and one obtains a light physical neutrino mass $`m_\nu m_D^2/M_N`$ plus a dominantly singlet neutrino of mass $`MM_N`$. The superpotential for the MSSM with a singlet neutrino superfield $`\widehat{N}^c`$ (for just a single generation), is given by $$\widehat{f}=\widehat{f}_{MSSM}+f_\nu ϵ_{ij}\widehat{L}^i\widehat{H}_u^j\widehat{N}^c+\frac{1}{2}M_N\widehat{N}^c\widehat{N}^c$$ (33) while the soft SUSY breaking terms now include $$=_{MSSM}m_{\stackrel{~}{\nu }_R}^2\|\stackrel{~}{\nu }_R\|^2+[A_\nu f_\nu ϵ_{ij}\stackrel{~}{L}^i\stackrel{~}{H}_u^j\stackrel{~}{\nu }_R^{}+\frac{1}{2}B_\nu M_N\stackrel{~}{\nu }_R^2+h.c.].$$ (34) The parameters $`A_\nu `$, $`B_\nu `$ and $`m_{\stackrel{~}{\nu }_R}`$ are assumed to be comparable to the weak scale. Many of the relevant RGEs have been presented in Ref.. Here we present the complete set needed for determining the sparticle spectrum at the weak scale. The one-loop RGEs for the gauge couplings and gaugino masses are unchanged from the MSSM case, since the $`\widehat{N}^c`$ superfield is a gauge singlet. The Yukawa coupling RGEs are $`{\displaystyle \frac{df_t}{dt}}`$ $`=`$ $`{\displaystyle \frac{f_t}{16\pi ^2}}\left[6f_t^2+f_b^2+f_\nu ^2{\displaystyle \frac{16}{3}}g_3^23g_2^2{\displaystyle \frac{13}{15}}g_1^2\right]`$ (35) $`{\displaystyle \frac{df_b}{dt}}`$ $`=`$ $`{\displaystyle \frac{f_b}{16\pi ^2}}\left[f_t^2+6f_b^2+f_\tau ^2{\displaystyle \frac{16}{3}}g_3^23g_2^2{\displaystyle \frac{7}{15}}g_1^2\right]`$ (36) $`{\displaystyle \frac{df_\tau }{dt}}`$ $`=`$ $`{\displaystyle \frac{f_\tau }{16\pi ^2}}\left[3f_b^2+4f_\tau ^2+f_\nu ^23g_2^2{\displaystyle \frac{9}{5}}g_1^2\right]`$ (37) $`{\displaystyle \frac{df_\nu }{dt}}`$ $`=`$ $`{\displaystyle \frac{f_\nu }{16\pi ^2}}\left[3f_t^2+f_\tau ^2+4f_\nu ^23g_2^2{\displaystyle \frac{3}{5}}g_1^2\right].`$ (38) The RGEs for $`m_Q^2`$, $`m_U^2`$, $`m_D^2`$, $`m_E^2`$ and $`m_{H_d}^2`$ are all unchanged from the MSSM. However, for $`m_L^2`$, $`m_{\stackrel{~}{\nu }_R}^2`$ and $`m_{H_u}^2`$, we have $`{\displaystyle \frac{dm_L^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{16\pi ^2}}\left[{\displaystyle \frac{3}{5}}g_1^2M_1^23g_2^2M_2^2+f_\tau ^2X_\tau +f_\nu ^2X_n\right]`$ (39) $`{\displaystyle \frac{dm_{\stackrel{~}{\nu }_R}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{4}{16\pi ^2}}\left[f_\nu ^2X_n\right]`$ (40) $`{\displaystyle \frac{dm_{H_u}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{16\pi ^2}}\left[{\displaystyle \frac{3}{5}}g_1^2M_1^23g_2^2M_2^2+3f_t^2X_t+f_\nu ^2X_n\right]`$ (41) where we have defined $`X_n=m_L^2+m_{\stackrel{~}{\nu }_R}^2+m_{H_u}^2+A_\nu ^2`$ and $`X_t`$ and $`X_\tau `$ are given in Ref. . Finally, the RGEs for the $`A_i`$ parameters are given by $`{\displaystyle \frac{dA_t}{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{16\pi ^2}}\left[\mathrm{\Sigma }c_ig_i^2M_i+6f_t^2A_t+f_b^2A_b+f_\nu ^2A_\nu \right]`$ (42) $`{\displaystyle \frac{dA_b}{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{16\pi ^2}}\left[\mathrm{\Sigma }c_i^{}g_i^2M_i+6f_b^2A_b+f_t^2A_t+f_\tau ^2A_\tau \right]`$ (43) $`{\displaystyle \frac{dA_\tau }{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{16\pi ^2}}\left[\mathrm{\Sigma }c_i^{\prime \prime }g_i^2M_i+3f_b^2A_b+4f_\tau ^2A_\tau +f_\nu ^2A_\nu \right]`$ (44) $`{\displaystyle \frac{dA_\nu }{dt}}`$ $`=`$ $`{\displaystyle \frac{2}{16\pi ^2}}\left[\mathrm{\Sigma }c_i^{\prime \prime \prime }g_i^2M_i+3f_t^2A_t+4f_\nu ^2A_\nu +f_\tau ^2A_\tau \right],`$ (45) where the $`c_i`$, $`c_i^{}`$ and $`c_i^{\prime \prime }`$ are given in Ref. , and $`c_i^{\prime \prime \prime }=\{\frac{3}{5},3,0\}`$. These RGEs apply for scales $`Q>M_N`$, while the MSSM RGEs are used below $`Q=M_N`$. Below the scale $`M_N`$ the effective theory does not contain the right handed neutrino or sneutrino, so that the running of the corresponding parameters is frozen at their values at this scale. The RGE for the parameter $`B_\nu `$ is irrelevant for our analysis. This model has been explicitly included in ISAJET version $`7.48`$, via the keyword $`SUGRHN`$, which allows, in addition to $`mSUGRA`$ and/or $`NUSUGi`$ inputs, the following: $$m_{\nu _\tau },M_N,m_{\stackrel{~}{\nu }_{\tau R}},A_\nu ,$$ (46) where all masses are entered in GeV units. Then the neutrino Yukawa coupling is calculated, and the MSSM+RHN RGEs are used at scales $`Q>M_N`$, while MSSM RGEs are used below $`Q=M_N`$. A sample spectrum of masses is shown in Table LABEL:trhn, assuming $`m_{\nu _\tau }=10^9`$ GeV, $`M_N=10^{13}`$ GeV, $`m_{\stackrel{~}{\nu }_{\tau R}}=200`$ GeV and $`A_\nu =0`$. The main effect is that the additional Yukawa coupling drives the third generation slepton masses to somewhat lower values than the massless neutrino case. An upper limit on the parameter $`\mathrm{tan}\beta `$ occurs in mSUGRA for $`\mu <0`$ due to a breakdown in the REWSB mechanism, where the $`H_u`$ mass is not driven sufficiently negative by RG running. For the MSSM+RHN model, the additional Yukawa coupling $`f_\nu `$ aids somewhat in driving $`m_{H_u}^2`$ negative. It is natural to ask how much the additional Yukawa coupling $`f_\nu `$ would help to increase the allowed range of $`\mathrm{tan}\beta `$ while still satisfying the REWSB constraint. As an example, we checked that for the case $`m_0=m_{1/2}=200`$ GeV, $`A_0=0`$, and $`\mu <0`$, for which $`\mathrm{tan}\beta 45`$ in the mSUGRA framework, the inclusion of a right-handed neutrino with $`m_N=10^{13}`$, $`(10^{10})`$ $`((10^7))`$ GeV, only increases this range to 45.3 (45.7) ((46)), assuming $`f_\nu =f_t`$ at the $`GUT`$ scale. ## V Unifying gauge groups with rank $`5`$: D-terms In general, if the MSSM is embedded in a $`GUT`$ gauge group with rank $`5`$, and the $`GUT`$ gauge group is spontaneously broken to a gauge group of lower rank, there are additional $`D`$-term contributions to scalar masses. The important thing is that these contributions affect TeV scale physics even if the scale at which the $`GUT`$ symmetry is broken is very large: since symmetry breaking is arranged to occur in a nearly $`D`$-flat direction, these $`D`$-term contributions to scalar masses are still of order the weak scale, even though the extra particles have masses $`M_{GUT}`$. The $`D`$-terms must be added to the various SUSY scalar mass squared parameters at the high scale at which the breaking occurs, so that these effectively lead to non-universal boundary conditions for scalar masses. Kolda and Martin have analysed these contributions for gauge groups which are subgroups of $`E_6`$, which encompasses a wide range of well-motivated $`GUT`$ group choices. $`E_6`$ contains in addition to the SM $`SU(3)_C\times SU(2)_L\times U(1)_Y`$ gauge symmetry two additional $`U(1)`$ symmetries labelled as $`U(1)_X`$ and $`U(1)_S`$. The $`D`$-term contributions to scalar masses can then be parametrized as, $`\mathrm{\Delta }m_Q^2`$ $`=`$ $`{\displaystyle \frac{1}{6}}D_Y{\displaystyle \frac{1}{3}}D_X{\displaystyle \frac{1}{3}}D_S,`$ (47) $`\mathrm{\Delta }m_D^2`$ $`=`$ $`{\displaystyle \frac{1}{3}}D_Y+D_X{\displaystyle \frac{2}{3}}D_S,`$ (48) $`\mathrm{\Delta }m_U^2`$ $`=`$ $`{\displaystyle \frac{2}{3}}D_Y{\displaystyle \frac{1}{3}}D_X{\displaystyle \frac{1}{3}}D_S,`$ (49) $`\mathrm{\Delta }m_L^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}D_Y+D_X{\displaystyle \frac{2}{3}}D_S`$ (50) $`\mathrm{\Delta }m_E^2`$ $`=`$ $`D_Y{\displaystyle \frac{1}{3}}D_X{\displaystyle \frac{1}{3}}D_S,`$ (51) $`\mathrm{\Delta }m_{H_d}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}D_Y{\displaystyle \frac{2}{3}}D_X+D_S,`$ (52) $`\mathrm{\Delta }m_{H_u}^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}D_Y+{\displaystyle \frac{2}{3}}D_X+{\displaystyle \frac{2}{3}}D_S,`$ (53) where $`D_Y`$ is the usual $`D`$-term associated with weak hypercharge breaking. In light of our ignorance of the mechanism of gauge symmetry breaking, the contributions $`D_X`$ and $`D_S`$ can be treated as additional dimensionful parameters, that can range over positive as well as negative values. ### A Minimal $`SO(10)`$ model with gauge symmetry breaking at $`Q=M_{GUT}`$ A simple special case of the above arises if the $`GUT`$ gauge group $`SO(10)`$ is assumed to directly break to the SM gauge group at $`Q=M_{GUT}`$ so the theory below this scale is the MSSM, possibly together with a right-handed neutrino and sneutrino. In this case, the three generations of matter superfields plus an additional SM gauge singlet right handed neutrino superfield for each generation are each elements of the 16 dimensional spinor representation of $`SO(10)`$, and so are taken to have a common mass $`m_{16}`$ above $`M_{GUT}`$. The Higgs superfields of the MSSM belong to a single 10 dimensional fundamental representation of $`SO(10)`$, and acquire a mass $`m_{10}`$. At $`Q=M_{GUT}`$, the gauge symmmetry breaking induces $$D_X0;D_Y=D_S=0$$ (54) so that at this scale the scalar masses are broken according to (50). The MSSM masses at $`M_{GUT}`$ may then be written as $`m_Q^2=m_E^2=m_U^2=m_{16}^2+M_D^2`$ (55) $`m_D^2=m_L^2=m_{16}^23M_D^2`$ (56) $`m_{H_{u,d}}^2=m_{10}^22M_D^2,`$ (57) where we have reparametrized $`D_X=3M_D^2`$. If the right-handed neutrino mass is substantially below the GUT scale, the soft breaking sneutrino mass would evolve as in Eq. (40); at the $`GUT`$ scale it would then be given by, $$m_{\stackrel{~}{\nu }_R}^2=m_{16}^2+5M_D^2.$$ (58) In minimal $`SO(10)`$, the superpotential above $`M_{GUT}`$ has the form, $`\widehat{f}=f\widehat{\psi }\widehat{\psi }\widehat{\varphi }+\mathrm{}`$ (59) with just a single Yukawa coupling per generation, where $`\widehat{\psi }`$ and $`\widehat{\varphi }`$ are the $`\mathrm{𝟏𝟔}`$ dimensional spinor and $`\mathrm{𝟏𝟎}`$ dimensional Higgs superfields, respectively. We neglect possible inter-generational mixing and also assume that the right-handed neutrino has a mass $`M_{GUT}`$. The dots represent terms including for instance higher dimensional Higgs representations and interactions responsible for the breaking of $`SO(10)`$. We assume here for simplicity that these couplings are suppressed relative to the usual Yukawa couplings. In minimal $`SO(10)`$, all the Yukawa couplings are unified above $`M_{GUT}`$, which forces one into a region of very large $`\mathrm{tan}\beta 50`$ which is actually excluded assuming universality of scalars if the constraint of radiative electroweak symmetry breaking is included. It has been suggested, and recently shown, that $`D`$-term contributions have the correct form to allow for Yukawa unified solutions to the SUSY particle mass spectrum consistent with radiative electroweak symmetry breaking. The parameter space of this model can be taken as $$m_{16},m_{10},M_D^2,m_{1/2},A_0,sign(\mu ),$$ (60) where $`M_D^2`$ can be either positive or negative. Yukawa coupling unification forces $`tan\beta `$ to be in the range 45-52 – for many purposes its exact value is irrelevant. The parameter space of minimal $`SO(10)`$ SUSY GUT models was explored in Ref. . It was found that, requiring Yukawa unification good to 5%, no solutions could be found for values of $`\mu >0`$, while many solutions could be obtained for $`\mu <0`$, but only for positive values of $`M_D^2`$. The $`D`$-term forces $`m_{H_u}<m_{H_d}`$ at $`Q=M_{GUT}`$: this is necessary to drive $`m_{H_u}^2`$ negative before $`m_{H_d}^2`$, as is required for REWSB with $`\mathrm{tan}\beta >1`$. Implications of this model for the dark matter relic density, $`bs\gamma `$ decay rate, and collider searches, are presented in Ref. A sample spectrum from the mSUGRA model and a corresponding case in Yukawa-unified $`SO(10)`$ are shown in Table LABEL:tso10\_1. The $`D`$-term splitting that ameliorates REWSB also leaves a distinct imprint on the masses of the matter scalars: the left- sleptons and right- down-type squarks have smaller GUT scale squared masses than their counterparts. This can be reflected in the weak scale spectrum where left- sleptons can be lighter than right- sleptons, and the right bottom squark can be by far the lightest of all the squarks – perhaps, even within the kinematic reach of the Main Injector upgrade of the Tevatron, though its detection may be complicated. Note also the smaller absolute value of the $`\mu `$ parameter in the $`SO(10)`$ case: this results in lighter charginos and neutralinos with substantial, or even dominant, higgsino components and a smaller $`\stackrel{~}{Z}_2\stackrel{~}{Z}_1`$ mass difference. Finally, we remark that for the case shown, the lighter $`\stackrel{~}{\tau }`$ is dominantly $`\stackrel{~}{\tau }_L`$. It is well known that SUSY models with $`\mu <0`$ and large $`\mathrm{tan}\beta `$ yield a large rate for the decay $`bs\gamma `$. Indeed, this class of models is already severely constrained by experimental results on radiative $`b`$-decays. However, additional non-universality between generations is possible in this framework, which could alter the gluino loop contributions, and hence the final branching fraction for $`bs\gamma `$ decay. ## VI Mass splittings in $`SO(10)`$ above $`Q=M_{GUT}`$ ### A Minimal $`SO(10)`$ As discussed above, the minimal $`SO(10)`$ model contains three generations of matter superfields each in a 16 dimensional representation, and a single Higgs superfield in the 10 dimensional representation. The superpotential is as given in Eq. (59) with $`f`$ the common Yukawa coupling for the third generation. Other terms will also be present, including Yukawa couplings for the first two generations, as well as more complicated Higgs representations necessary for $`SO(10)`$ breaking. We will assume the Yukawa couplings involving these fields are all small, so the dominant contribution to RGE running comes from just the superpotential (59). We also assume associated $`SO(10)`$ soft SUSY breaking parameters: $`m_{16}`$, $`m_{10}`$, $`m_{1/2}`$ and $`A`$. Then the RGEs in the minimal $`SO(10)`$ model are calculable. For the gauge coupling we have, $$\frac{dg}{dt}=\frac{g^3}{16\pi ^2}(S24),$$ (61) where $`S`$ is the sum of Dynkin indices of the various chiral superfields in the model. With the above minimal field content, $`S=7`$. However, additional fields associated for instance with $`SO(10)`$ breaking ought to be present, and will increase the value of $`S`$. The Yukawa coupling RGE is, $$\frac{df}{dt}=\frac{1}{16\pi ^2}f\left(14f^2\frac{63}{2}g^2\right).$$ (62) For the gaugino mass we have the following RGE: $$\frac{dm_{1/2}}{dt}=\frac{1}{16\pi ^2}2(S24)g^2m_{1/2}$$ (63) For the scalar masses we have: $`{\displaystyle \frac{dm_{16}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left[10f^2\left(2m_{16}^2+m_{10}^2+A^2\right)45g^2m_{1/2}^2\right]`$ (64) $`{\displaystyle \frac{dm_{10}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}\left[8f^2\left(2m_{16}^2+m_{10}^2+A^2\right)36g^2m_{1/2}^2\right].`$ (65) Finally, the RGE for the trilinear mass parameter is $$\frac{dA}{dt}=\frac{1}{16\pi ^2}\left(28f^2A+63g^2m_{1/2}\right).$$ (66) As an illustration, we adopt the minimal $`SO(10)`$ case 5 spectra from Ref. for which Yukawa couplings unify at $`M_{GUT}`$. The model parameters and mass spectrum is listed in the “$`M_{GUT}`$ Unification” column of Table LABEL:msoten. We begin by using $`f(M_{GUT})=0.553`$ and $`g_{GUT}=0.706`$ (as given by the minimal $`SO(10)`$ model). We then evolve (using $`S=7`$) from $`M_{GUT}`$ to $`M_P`$ to find the corresponding Planck scale gauge and Yukawa couplings. At $`M_P`$, we assume universality of the three generations with $`m_{16}=629.8`$ GeV, while $`m_{10}=836.2`$ GeV. At $`M_{GUT}`$, we take $`m_{1/2}=348.8`$ GeV and $`A=186.5`$ GeV, with a $`D`$-term $`M_D=135.6`$ GeV. A Yukawa unified solution is obtained for $`\mathrm{tan}\beta =52.1`$ and the corresponding spectrum is shown in the last column titled “$`M_P`$ Unification”. In Fig. 2, we show by the solid lines the effect of running of SSB parameters between $`M_P`$ and $`M_{GUT}`$ for the minimal $`SO(10)`$ model, for parameters as in Table LABEL:msoten. The dashed lines show the corresponding situation for $`S=15`$, i.e. with one additional adjoint included; in this case, the running of the gauge coupling between $`M_{GUT}`$ and $`M_P`$ (see Eq. 61) is somewhat slower. We see that the splitting $`\delta m_{16}^2`$ between the GUT scale mass parameters of the first (or second) and third generations is reduced, albeit by a small amount. <sup>§</sup><sup>§</sup>§Since the right hand side of Eq. (62) is more negative when $`S=15`$ as compared to the $`S=7`$ case, the corresponding $`f`$ runs to smaller values in the former case. If we now consider the evolution of $`\delta m_{16}^2`$, for which the term depending on $`g`$ drops out, we see that this difference runs the most for $`S=7`$ for which $`f`$ is largest. In this sense, the difference shown by the solid lines may be regarded as a bound. The effect of $`SO(10)`$ running is that the first two generations of matter scalars run to higher masses, while the Higgs masses and third generation masses decrease somewhat. The corresponding weak scale sparticle masses are listed in Table LABEL:msoten, without and with the effect of Planck to $`GUT`$ scale running. The main effect is a $`27`$% change in the mass difference between the (lightest) charged sleptons of the first and third generations. ### B General $`SO(10)`$ More generally, we may take the two MSSM Higgs doublets to live in different fundamental representations of $`SO(10)`$: $`\widehat{H}_u\widehat{H}_2`$ and $`\widehat{H}_d\widehat{H}_1`$. Then the superpotential can be written as $$\widehat{f}=f_t\widehat{\psi }\widehat{\psi }\widehat{H}_2+f_b\widehat{\psi }\widehat{\psi }\widehat{H}_1,$$ (67) so that there exist two Yukawa couplings above the GUT scale, and just $`f_b=f_\tau `$ unification, which can occur for a much wider range of $`\mathrm{tan}\beta `$ values, is required. In addition to the usual scalar masses, as in Ref. , we include an off-diagonal mass term $`m_{H_{12}}^2(H_1^{}H_2+H_2^{}H_1)`$. As in minimal $`SO(10)`$, there should also be at least higher dimensional Higgs representations present responsible for $`SO(10)`$ breaking, but again, we ignore these. We give here the RGEs for the general $`SO(10)`$ model, thereby completing the results of Refs. . For the gauge coupling constant we have: $$\frac{dg}{dt}=\frac{g^3}{16\pi ^2}(S24),$$ (68) where $`S`$ again is the sum of the Dynkin indices of the $`SO(10)`$ fields. For just two 10 dimensional Higgs multiplets and 3 generations of matter, $`S=8`$. For gaugino masses, we again have $$\frac{dm_{1/2}}{dt}=\frac{1}{16\pi ^2}2(S24)g^2m_{1/2}.$$ (69) The Yukawa coupling RGEs are: $`{\displaystyle \frac{df_t}{dt}}`$ $`=`$ $`{\displaystyle \frac{f_t}{16\pi ^2}}\left(14f_t^2+14f_b^2{\displaystyle \frac{63}{2}}g^2\right)`$ (70) $`{\displaystyle \frac{df_b}{dt}}`$ $`=`$ $`{\displaystyle \frac{f_b}{16\pi ^2}}\left(14f_t^2+14f_b^2{\displaystyle \frac{63}{2}}g^2\right).`$ (71) The RGEs for the scalar masses are given by $`{\displaystyle \frac{dm_{16}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{10}{16\pi ^2}}[f_t^2(2m_{16}^2+m_{H_2}^2)+f_b^2(2m_{16}^2+m_{H_1}^2)+2f_tf_bm_{H_{12}}^2`$ (73) $`+(A_t^2f_t^2+A_b^2f_b^2){\displaystyle \frac{9}{2}}g^2m_{1/2}^2]`$ $`{\displaystyle \frac{dm_{H_1}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{8}{16\pi ^2}}\left[f_b^2(2m_{16}^2+m_{H_1}^2)+f_tf_bm_{H_{12}}^2+A_b^2f_b^2{\displaystyle \frac{9}{2}}g^2m_{1/2}^2\right]`$ (74) $`{\displaystyle \frac{dm_{H_2}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{8}{16\pi ^2}}\left[f_t^2(2m_{16}^2+m_{H_2}^2)+f_tf_bm_{H_{12}}^2+A_t^2f_t^2{\displaystyle \frac{9}{2}}g^2m_{1/2}^2\right]`$ (75) $`{\displaystyle \frac{dm_{H_{12}}^2}{dt}}`$ $`=`$ $`{\displaystyle \frac{4}{16\pi ^2}}\left[f_tf_b(4m_{16}^2+m_{H_1}^2+m_{H_2}^2+2A_tA_b)+(f_t^2+f_b^2)m_{H_{12}}^2\right].`$ (76) Finally, the RGEs for the $`A`$ parameters are $`{\displaystyle \frac{dA_t}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}(28f_t^2A_t+20f_b^2A_b+63g^2m_{1/2})`$ (77) $`{\displaystyle \frac{dA_b}{dt}}`$ $`=`$ $`{\displaystyle \frac{1}{16\pi ^2}}(28f_b^2A_b+20f_t^2A_t+63g^2m_{1/2}).`$ (78) We show in Fig. 3 the running of SSB parameters in the general $`SO(10)`$ model using $`GUT`$ scale values of $`g=0.717`$, $`f_t=0.534`$ and $`f_b=0.271`$, as in Fig. 1. Except for $`m_{H_{12}}^2`$ which is fixed to be zero at $`Q=M_P`$, the SSB parameters are also as in this figure. The main effect is again a significant splitting between first or second and third generation scalar masses at the $`GUT`$ scale. Some splitting between $`m_{H_u}`$ and $`m_{H_d}`$ also occurs, with $`m_{H_u}^2<m_{H_d}^2`$ as desired. The corresponding weak scale sparticle masses are shown in Table LABEL:gsoten. The $`GUT`$ scale SSB term splitting results in somewhat heavier scalars than in the $`mSUGRA`$ case. For this example, because most of the weak scale squark mass comes from the RG evolution, the effect is more pronounced for sleptons than squarks. In particular, this increase is just a few percent for squarks, but as much as 22% for sleptons. ## VII Supersymmetric missing partner models with hypercolor In this variety of models, the gauge group is of the type $`G_{GUT}\times G_H`$, where the first group is $`SU(5)`$ or $`SO(10)`$ and the second is related to a ‘hypercolor’ interaction . While the weak $`SU(2)`$ is completely contained in the first factor, colour $`SU(3)`$ is not embedded in either of the factors. Although the gauge group is not simple, an approximate unification of the gauge coupling constants of the group $`SU(3)_C\times SU(2)\times U(1)`$ is achieved if the couplings of $`G_H`$ are large enough. These models provide a solution to the doublet-triplet splitting problem by the missing partner mechanism. Since the MSSM gauginos do not belong to a single multiplet of a simple gauge group, their masses do not obey the usual unification condition , resulting in non-universality of gaugino masses. However, if usual squarks and sleptons and the MSSM Higgs fields are singlets of $`G_H`$, universality of scalar masses is still possible, as for instance, in the $`SU(5)_{GUT}\times SU(3)_H\times U(1)_H`$ model of Ref. , where the hypercharge $`U(1)`$ is a combination of $`U(1)_H`$ and a $`U(1)`$ subgroup in the first factor. In this case, the following relations among gauge couplings hold at the unification scale $$\frac{1}{g_1^2}=\frac{1}{g_{GUT}^2}+\frac{1}{15g_{H1}^2},\frac{1}{g_2^2}=\frac{1}{g_{GUT}^2},\frac{1}{g_3^2}=\frac{1}{g_{GUT}^2}+\frac{1}{g_{H3}^2},$$ (79) where $`\sqrt{3/5}g_1`$, $`g_2`$, and $`g_3`$ are the gauge couplings of the $`U(1)_Y`$, $`SU(2)_L`$, and $`SU(3)_C`$ SM groups, and $`g_{GUT}`$, $`g_{H3}`$, and $`g_{H1}`$ are the $`SU(5)_{GUT}`$, $`SU(3)_H`$, and $`U(1)_H`$ unified groups respectively. Clearly from Eq. (79) we see that the unification of the gauge coupling constants from low energy data is achieved if $`g_{H1}^2g_{GUT}^2`$ and $`g_{H3}^2g_{GUT}^2`$. In addition, considering that the prediction for $`\alpha _s`$ at the weak scale in SUSY GUT models (without threshold corrections) is higher than the world averaged experimental value, it was argued that the correction introduced by hypercolor moves the prediction for $`\alpha _s`$ in the correct direction. It was found that : $$\alpha _s(m_Z)0.130\frac{0.014}{\alpha _{H3}}\frac{0.010}{15\alpha _{H1}}$$ (80) where $`\alpha _i=g_i^2/4\pi `$ and threshold corrections have been neglected. In order for $`\alpha _s`$ not to shift too much, we must have $`\alpha _{H3}0.6`$ and $`\alpha _{H1}0.03`$, though for $`\alpha _{H1}`$ as small as 0.03, $`g_2^2g_1^2=0.18g_1^2g_2^2`$. Above the GUT scale there are three gauginos associated to the groups $`SU(5)_{GUT}`$, $`SU(3)_H`$, and $`U(1)_H`$ whose masses we denote $`m_{1/2}`$, $`M_{H3}`$ and $`M_{H1}`$ respectively. Below the GUT scale we have the MSSM and the three MSSM gauginos are a linear combination of the former ones. The masses of the bino, wino, and gluino are then given by A somewhat different model based on the group $`SO(10)_{GUT}\times SO(6)_H`$ also has non-universal MSSM gaugino masses. However, since the hypercolor group is simple, there is one relation between them, $`M_1`$ $`=`$ $`g_1^2\left({\displaystyle \frac{m_{1/2}}{g_{GUT}^2}}+{\displaystyle \frac{M_{H1}}{15g_{H1}^2}}\right),`$ (81) $`M_2`$ $`=`$ $`m_{1/2},`$ (82) $`M_3`$ $`=`$ $`g_3^2\left({\displaystyle \frac{m_{1/2}}{g_{GUT}^2}}+{\displaystyle \frac{M_{H3}}{g_{H3}^2}}\right).`$ (83) The thing to note is that $`M_{H1,3}/\alpha (H_{1,3})`$ are renormalization group invariants (at one loop) so that $`M_{H1,3}/g_{H1,3}^2`$ need not be small even when $`g_{H1,3}^2`$ is large. The relative magnitude of the three masses $`m_{1/2}`$, $`M_{H3}`$ and $`M_{H1}`$ is unknown because it depends on the SUSY breaking mechanism. One might naively suppose that they are of the same order of magnitude; in this case, gaugino masses could be significantly different at the GUT scale, though the magnitude of the non-universality would be limited because, as noted above, the couplings $`g_{H1}`$ and $`g_{H3}`$ have to be considerably larger than $`g_{GUT}`$. There is no reason, however, why $`M_{H1}`$ and $`M_{H3}`$ cannot be much larger than $`m_{1/2}`$. Indeed in scenarios with dilaton dominated SUSY breaking, we have $$\frac{m_{1/2}}{g_{GUT}^2}=\frac{M_{H1}}{g_{H1}^2}=\frac{M_{H3}}{g_{H3}^2},$$ (84) so that gaugino mass splittings of O(100%) are expected. In Fig. 4a we plot non–universal gaugino masses of the MSSM as a function of a common hypercolor gaugino mass $`M_{H1}=M_{H3}M_H`$. We take $`m_{1/2}=200`$ GeV, $`g_{GUT}=0.716`$, and two different choices for the hypercolor gauge couplings: $`\alpha _{H1}=0.1`$ and $`\alpha _{H3}=0.7`$ in solid lines, and $`\alpha _{H1}=0.5`$ and $`\alpha _{H3}=0.8`$ in dashed lines. As indicated in Eq. (82) the wino mass $`M_2`$ is always equal to $`m_{1/2}=200`$ GeV. The other two gaugino masses are larger (smaller) than $`M_2`$ if the hypercolor gaugino mass is larger (smaller) than $`m_{1/2}`$. The gluino mass deviates more from $`M_2`$ compared to the bino mass because of the factor 15 in Eq. (82) and our choice of values for other parameters. The larger the hypercolor gauge couplings, the smaller the deviations from universality. In addition, if the common hypercolor mass is equal to $`m_{1/2}`$ there is no deviation from universality no matter the value of the hypercolor gauge couplings. In Fig. 4b, we show the same gaugino masses but assuming instead that $`M_{H1}/g_{H1}^2=M_{H_3}/g_{H3}^2`$. The three gauge couplings are chosen exactly as in frame a) so that there is a large hierarchy between the masses of the gauginos of the three groups. The cross denotes the dilaton-dominated scenario for which point Eq. (84) is satisfied. Indeed we see that very large non-universality of gaugino masses may be possible. In Fig. 5, we show several weak scale sparticle masses versus the same parameter $`M_H`$ as in Fig. 4 for parameter values corresponding to the solid curves in this figure. The two frames illustrate the results for the same choices of the gaugino masses as in Fig. 4. In frame a) we see that the non-colored sparticle masses hardly vary at all versus $`M_H`$, while the gluino and squark masses can vary by up to 12%. This is presumably because the coloured sparticle masses run considerably more than those of uncoloured sparticles coupled with the fact that $`M_3`$ varies more with $`M_H`$ than $`M_1`$ does, and $`M_2`$ does not change at all. The variation is, of course, much more dramatic in frame b). For very large values of $`M_{H1}`$, the coloured sparticles as well as the heavier chargino and neutralinos become very heavy, and may be in conflict with fine-tuning considerations. We also mention that although $`M_1`$ starts out larger than $`M_2`$ at the GUT scale (but not by a huge amount), $`\frac{M_1}{M_2}`$ is driven to a value close to $`\frac{1}{2}`$ at the weak scale for acceptable values of $`M_{H1}`$: it would be interesting to examine whether precise measurements of masses and mixing angles could lead to observable deviations from expectations in mSUGRA or gauge-mediated SUSY breaking frameworks. In the same vein, we also mention that $`m(\stackrel{~}{e}_R)`$ also increases slowly from 132 GeV in mSUGRA to 134 GeV for the dilaton dominated scenario to 151 GeV for the extreme case with $`M_{H1}=3`$ TeV, while $`m(\stackrel{~}{e}_L)`$ is roughly constant. This is because the RG evolution of $`m(\stackrel{~}{e}_R)^2`$ is due to hypercharge gauge interactions, and $`M_1`$ starts out bigger than $`M_2`$ (which is independent of $`M_{H1}`$). ## VIII Models with effective supersymmetry The SM exhibits accidental global symmetries which inhibit flavor–changing neutral currents (FCNC), lepton flavor violation (LFV), electric dipole moments (EDM) of electron and neutron, and proton decay, as opposed to the MSSM where degeneracy or alignment in the mass matrices has to be invoked. On the other hand, supersymmetry stabilizes the scalar masses under radiative corrections, contrary to the SM where it is hard to understand the hierarchy between the Higgs mass and the Planck scale. The models presented in this section aim to combine the good features of both the SM and the MSSM. There are two mass scales: gauginos, higgsinos, and third generation squarks are sufficiently light ($`1`$ TeV) to naturally stabilize the Higgs mass and the electroweak scale, while the first two generations of squarks and sleptons (whose Yukawa couplings to Higgs are very small) are sufficiently heavy ($`\stackrel{~}{M}5`$ to 20 TeV) to suppress FCNC, LFV, etc.. This class of models, called Effective Supersymmetry, does not invoke degeneracy or alignment in the mass matrices. In one of the realizations of Effective Supersymmetry, the first two generations of squarks and sleptons, together with the down–type Higgs, are composite, with constituents that carry a “superglue” charge, and have a mass $`\stackrel{~}{M}`$. Gauge superfields, third generation superfields and the up–type Higgs superfield are taken to be fundamental and neutral under superglue, with perturbative couplings to the constituents, so that their mass is suppressed relative to the mass of the composites. In this way, the spectrum is characterized as follows. * Gaugino masses are light and can be non-universal with masses given by $`M_i=n_i(\alpha _i/4\pi )\stackrel{~}{M}`$, where $`n_i`$ are numerical factors that can be as large as $`𝒪(10)`$. * Left and right squark and slepton masses for the first two generations are of the order of $`\stackrel{~}{M}`$. * Left and right squark and slepton masses for the third generation are of the order of $`(\lambda _3/4\pi )\stackrel{~}{M}`$; for $`\lambda _31`$, this is an order of magnitude smaller than $`\stackrel{~}{M}`$. * The down–Higgs mass satisfy $`m_{H_d}\stackrel{~}{M}`$. The up–Higgs mass on the other hand, is given by $`m_{H_u}(\lambda _H/4\pi )\stackrel{~}{M}`$, where $`\lambda _H`$ is its perturbative coupling to the constituents. Therefore, there is only one Higgs in the low energy theory and $`\mathrm{tan}\beta 4\pi /\lambda _H`$ is large. * The “$`\mu `$–term” and the “$`B\mu `$–term” respectively satisfy $`\mu (\lambda _H/4\pi )\stackrel{~}{M}`$ and $`B\mu (\lambda _H/4\pi )\stackrel{~}{M}^2`$. To obtain $`m_{H_u}`$ 100 GeV, we require $`\lambda _H/4\pi 10^2`$, while $`\lambda _3/4\pi 10^1`$ ensures $`m_{\stackrel{~}{t}}1`$ TeV. If the hierarchy of scalar masses is already present at the unification scale, then it has been shown that unless the stop mass squared at the unification scale is taken to be well above (1 TeV)<sup>2</sup>, two-loop contributions to scalar renormalization group equations drive the top squark mass squared negative well before the weak scale, resulting in a breakdown of color symmetry. Thus, this simple class of models seems to be ruled out by fine-tuning considerations. To account for this class of constraints, we have implemented the full set of two-loop MSSM RGEs in ISAJET versions $`7.50`$. Very recently, Hisano et al. have identified scenarios in which first and second generation scalars can be much heavier than gauginos and scalars of the third generation, and for which the scalar masses are renormalization group invariant (so that the constraints of Ref. are not relevant) as long as gaugino masses are neglected in the RGEs. These constraints are also inapplicable in models in which the assumption of the scalar hierarchy is made for mass parameters at a scale $`1050`$ TeV, since then there are no large logs that drive $`m^2`$ to negative values. In this case, however, model-dependent finite contributions to $`\delta m^2`$ are no longer negligible, and need to be examined to discuss the viability of any particular model . Yet another possibility has been considered in Ref.. These authors begin with all scalar masses initially at the multi-TeV scale at or above $`M_{GUT}`$, and show that for certain choices of $`M_{GUT}M_{Planck}`$ scale boundary conditions on the scalar masses and $`A`$ parameters– keeping gaugino masses at the weak scale– the third generation sfermion and Higgs masses are driven to weak scale values, while scalars of the first two generations remain heavy. Such a scenario is particularly attractive in the context of minimal $`SO(10)`$. In this case, with Yukawa coupling unification plus a singlet $`\widehat{N}^c`$, particularly simple boundary conditions , $$4m_{16}^2=2m_{10}^2=A^2$$ (85) lead to sub-TeV scale third generation scalar masses, while first and second generation scalar masses can be as high as 20 TeV. If instead the boundary value of $`A`$ is taken to be at the weak scale, the hierarchy generated is somewhat smaller. Examples of sparticle mass spectra were not generated in Ref. , where it was noted that this scenario shares the problem of obtaining correct radiative breaking of electroweak symmetry common to most high $`\mathrm{tan}\beta `$ scenarios: in examples shown in Ref. and Ref. , the two Higgs SSB masses stay positive at all scales in their evolution to the weak scale, with $`m_{H_u}>m_{H_d}`$, contrary to what is needed for REWSB. In a recent analysis it has been shown that if the boundary conditions in Eq. (85) are augmented by $`SO(10)`$ $`D`$-terms, it is possible to obtain the desired inverted mass hierarchy amongst the squarks together with radiative electroweak symmetry breaking. This then yields a calculable model based on the gauge group $`SO(10)`$ with (approximate) unification of Yukawa couplings. The analysis in Ref. took the right-handed neutrino mass to be fixed near $`10^{13}`$ GeV, and obtained “crunch” factor values $`S`$ up to $`57`$ for full $`SO(10)`$ $`D`$-terms, and factors of $`S`$ up to 9 if splittings were applied only to the soft SUSY breaking Higgs masses. The crunch factor $`S`$ is defined as $`S={\displaystyle \frac{3(m_{u_L}^2+m_{d_L}^2+m_{u_R}^2+m_{d_R}^2)+m_{\stackrel{~}{e}_L}^2+m_{\stackrel{~}{e}_R}^2+m_{\stackrel{~}{\nu }_e}^2}{3(m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{b}_1}^2+m_{\stackrel{~}{t}_2}^2+m_{\stackrel{~}{b}_2}^2)+m_{\stackrel{~}{\tau }_1}^2+m_{\stackrel{~}{\tau }_2}^2+m_{\stackrel{~}{\nu }_\tau }^2}}.`$ These values are considerably below those quoted in Ref. , where a more idealized case was considered. Effective supersymmetry is not as mature a framework as mSUGRA or the gauge-mediated SUSY breaking. Except for the inverted hierarchy model of the previous paragraph, all the models discussed in this Section suffer from incompleteness which preclude computations at as thorough a level. The scenario in Ref. involves new unknown strong dynamics at the 10 TeV scale. Models where the splitting between third generation scalars and those of the other generations has a dynamical origin suffer from the fact that this dynamics does not break electroweak symmetry: the mass spectrum thus does not appear to be calculable unless deviations such as non-universality are imposed. These considerations notwithstanding, collider events for generic effective SUSY models can be generated with ISAJET by using the weak-scale $`MSSMi`$ keywords, with independent weak scale SSB masses as inputs. One may enter multi-TeV scale first and second generation scalar masses, while using sub-TeV scale gaugino masses, third generation scalar masses and $`\mu `$ parameters. In the scenario of Ref., $`A`$-terms are $`𝒪(100)`$ GeV or smaller, while $`m_A`$ is very large. Sparticle mass spectra from the radiatively generated inverted mass hierarchy solution due to Bagger et al. are not possible without modifications that allow REWSB to occur. Two possibilities are the non-universalities due to $`SO(10)`$ $`D`$-terms, or ad-hoc Higgs sector splittings. These may be implemented in ISAJET using the NUSUG inputs along with the right-handed neutrino solution. In ISAJET, if a zero physical neutrino mass is entered, then the Yukawa couplings $`f_t`$ and $`f_\nu `$ automatically unify. It remains to be seen whether the resulting inverted mass hierarchy is truly sufficient to solve problems due to FCNCs, LFVs and the EDM of the electron and neutron. ## IX Anomaly-mediated SUSY breaking In most models, soft SUSY breaking parameters of the low energy effective theory are thought to receive contributions from gravitational or gauge interactions which are considered to be messengers of SUSY breaking in a hidden sector. It has recently been recognized that there is an additional contribution, that originates in the super-Weyl anomaly, which is always present when SUSY is broken. In models without SM gauge singlet superfields that can acquire a Planck scale $`vev`$, the usual supergravity contribution to gaugino masses is suppressed by an additional factor $`\frac{M_{SUSY}}{M_P}`$ relative to $`m_{\frac{3}{2}}=M_{SUSY}^2/M_P`$, and the anomaly-mediated contribution can dominate. These contributions are determined in terms of the SUSY breaking scale by the corresponding $`\beta `$ functions. $$M_i=\frac{\beta _g}{g}m_{\frac{3}{2}},$$ (86) where $`\beta _i`$ is the one–loop beta function, defined by $`\beta _{g_i}dg_i/d\mathrm{ln}\mu =b_ig_i^3+\mathrm{}`$. The gaugino masses are not universal, but given by the ratios of the respective $`\beta `$-functions. In general, however, Kähler potential couplings between the observable sector and the hidden sector (Goldstino) field, which are generically not forbidden by a symmetry, result in large gravity contributions ($`m_{\frac{3}{2}}`$) to scalar masses which would completely dominate the corresponding anomaly-mediated contributions. These gravity contributions can be strongly suppressed if the SUSY breaking and visible sectors reside on different branes, and are “sufficiently separated” in a higher dimensional space: in this case, the suppression is the result of geometry and not a symmetry, though then one has to wonder about the dynamics that results in such a geometry. The anomaly-mediated contribution is given by, $$m_{\stackrel{~}{q}}^2=\frac{1}{4}\left\{\frac{d\gamma }{dg}\beta _g+\frac{d\gamma }{df}\beta _f\right\}m_{\frac{3}{2}}^2$$ (87) where $`\beta _g`$ and $`\beta _f`$ are the $`\beta `$ functions for gauge and Yukawa interactions, respectively, and $`\gamma =\mathrm{ln}Z/\mathrm{ln}\mu `$, with $`Z`$ the wave function renormalization constant. Notice that this is comparable to the corresponding contribution to the gaugino masses. Furthermore, since Yukawa interactions are negligible for the first two generations, the anomaly-mediated contributions to scalar masses of the first two generations are essentially equal. Unfortunately, however , the anomaly contribution turns out to be negative for sleptons, necessitating additional sources for the squared masses of scalars. There are several proposals in the literature, but phenomenologically it suffices to add a universal contribution $`m_0^2`$ (which, of course, preserves the degeneracy between the first two generations of scalars) to Eq. (87), and regard $`m_0`$ as an additonal parameter. Finally, in the sign convention of ISAJET This is opposite to that used in Ref. ., the anomaly-mediated contribution to the trilinear SUSY breaking scalar coupling is given by, $$A_f=+\frac{\beta _f}{f}m_{\frac{3}{2}}.$$ (88) It is assumed that the ad hoc introduction of $`m_0^2`$ in Eq. (87) does not affect the other relations. ### A The Minimal Anomaly-Mediated SUSY Breaking Model (AMSB) In this framework, it is assumed that the anomaly-mediated SUSY breaking contributions to the soft-SUSY breaking contributions dominate, and further, that the introduction of the parameter $`m_0^2`$ is sufficent to circumvent the problem of negative squared masses for sleptons. The parameter space of the model consists of $$m_0,m_{3/2},\mathrm{tan}\beta \mathrm{and}sign(\mu ).$$ (89) In this case, gaugino masses are given by $`M_1`$ $`=`$ $`{\displaystyle \frac{33}{5}}{\displaystyle \frac{g_1^2}{16\pi ^2}}m_{3/2},`$ (90) $`M_2`$ $`=`$ $`{\displaystyle \frac{g_2^2}{16\pi ^2}}m_{3/2},\mathrm{and}`$ (91) $`M_3`$ $`=`$ $`3{\displaystyle \frac{g_3^2}{16\pi ^2}}m_{3/2}.`$ (92) Third generation scalar masses are given by $`m_U^2`$ $`=`$ $`\left({\displaystyle \frac{88}{25}}g_1^4+8g_3^4+2f_t\widehat{\beta }_{f_t}\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2,`$ (93) $`m_D^2`$ $`=`$ $`\left({\displaystyle \frac{22}{25}}g_1^4+8g_3^4+2f_b\widehat{\beta }_{f_b}\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2,`$ (94) $`m_Q^2`$ $`=`$ $`\left({\displaystyle \frac{11}{50}}g_1^4{\displaystyle \frac{3}{2}}g_2^4+8g_3^4+f_t\widehat{\beta }_{f_t}+f_b\widehat{\beta }_{f_b}\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2,`$ (95) $`m_L^2`$ $`=`$ $`\left({\displaystyle \frac{99}{50}}g_1^4{\displaystyle \frac{3}{2}}g_2^4+f_\tau \widehat{\beta }_{f_\tau }\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2,`$ (96) $`m_E^2`$ $`=`$ $`\left({\displaystyle \frac{198}{25}}g_1^4+2f_\tau \widehat{\beta }_{f_\tau }\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2,`$ (97) $`m_{H_u}^2`$ $`=`$ $`\left({\displaystyle \frac{99}{50}}g_1^4{\displaystyle \frac{3}{2}}g_2^4+3f_t\widehat{\beta }_{f_t}\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2,`$ (98) $`m_{H_d}^2`$ $`=`$ $`\left({\displaystyle \frac{99}{50}}g_1^4{\displaystyle \frac{3}{2}}g_2^4+3f_b\widehat{\beta }_{f_b}+f_\tau \widehat{\beta }_{f_\tau }\right){\displaystyle \frac{m_{3/2}^2}{(16\pi ^2)^2}}+m_0^2.`$ (99) The $`A`$-parameters are given by $`A_t`$ $`=`$ $`{\displaystyle \frac{\widehat{\beta }_{f_t}}{f_t}}{\displaystyle \frac{m_{3/2}}{16\pi ^2}},`$ (100) $`A_b`$ $`=`$ $`{\displaystyle \frac{\widehat{\beta }_{f_b}}{f_b}}{\displaystyle \frac{m_{3/2}}{16\pi ^2}},\mathrm{and}`$ (101) $`A_\tau `$ $`=`$ $`{\displaystyle \frac{\widehat{\beta }_{f_\tau }}{f_\tau }}{\displaystyle \frac{m_{3/2}}{16\pi ^2}}.`$ (102) In the above, we have $`\widehat{\beta }_{f_t}`$ $`=`$ $`16\pi ^2\beta _t=f_t\left({\displaystyle \frac{13}{15}}g_1^23g_2^2{\displaystyle \frac{16}{3}}g_3^2+6f_t^2+f_b^2\right),`$ (103) $`\widehat{\beta }_{f_b}`$ $`=`$ $`16\pi ^2\beta _b=f_b\left({\displaystyle \frac{7}{15}}g_1^23g_2^2{\displaystyle \frac{16}{3}}g_3^2+f_t^2+6f_b^2+f_\tau ^2\right),`$ (104) $`\widehat{\beta }_{f_\tau }`$ $`=`$ $`16\pi ^2\beta _\tau =f_\tau \left({\displaystyle \frac{9}{5}}g_1^23g_2^2+3f_b^2+4f_\tau ^2\right).`$ (105) The first two generations of squark and slepton masses are given by the corresponding formulae above with the Yukawa couplings set to zero. This model has been implemented in ISAJET versions $`7.45`$, using the $`AMSB`$ keyword, which allows input of the above parameter space set. In $`ISAJET`$, it is easiest to implement the above masses at scale $`Q=M_{GUT}`$, and proceed with evolution to the weak scale. Then the $`B`$ and $`\mu ^2`$ parameters are calculated in accord with the constraint from radiative electroweak symmetry breaking. The most notable feature of this framework is the hierarchy of gaugino masses. The gluino is (as in mSUGRA) much heavier than the electroweak gauginos, but the novel feature is that $`\frac{M_1}{M_2}3.2`$, so that the wino is by far the lightest supersymmetric particle (LSP). The wino LSP scenario has several implications for phenomenology, the most important of which is the near degeneracy of the chargino and the (wino-like) neutralino LSP. One loop corrections, which make the dominant contribution to the chargino-neutralino mass gap, have been included in ISAJET v7.46 (in the gaugino limit). The phenomenology can be sensitive to this mass difference. In Table LABEL:tamsb, we show spectra generated from the minimal AMSB model for two values of $`m_0`$, with other parameters being the same. Note that the parameter $`m_{3/2}`$ should be selected typically above 25,000 GeV to avoid constraints from LEP experiments. ¿From the spectra shown, we immediately see several well-known aspects of the $`AMSB`$ spectrum. Most notably, we see that the $`\stackrel{~}{W}_1`$ and $`\stackrel{~}{Z}_1`$ are nearly degenerate in mass, so that in addition to the usual leptonic decay modes $`\stackrel{~}{W}_1\stackrel{~}{Z}_1\mathrm{}\nu `$, the only other allowed (and in these cases dominant) decay of the chargino is $`\stackrel{~}{W}_1^\pm \stackrel{~}{Z}_1\pi ^\pm `$. The chargino has a very small width, corresponding to a lifetime $`1.5\times 10^9`$ s, so that it would be expected to travel a significant fraction of a meter before decaying . Secondly, the $`\stackrel{~}{\mathrm{}}_L`$ and $`\stackrel{~}{\mathrm{}}_R`$ are nearly mass degenerate. This degeneracy (which seems fortuitous) is much tighter than expected in the mSUGRA framework and certainly in the gauge-mediated SUSY breaking framework. Their mass scale is largely determined by the parameter $`m_0`$, and it is possible that for small enough $`m_0`$ slepton signals may be detectable at the next generation of $`e^+e^{}`$ colliders or even at the LHC. Another interesting feature (which may serve to distinguish the cases shown from mSUGRA) is that the $`\stackrel{~}{\tau }_L\stackrel{~}{\tau }_R`$ mixing is near maximal. The prospects for measuring this have been discussed in Ref.. In the minimal AMSB framework, $`m_{\stackrel{~}{W}_1}m_{\stackrel{~}{Z}_1}`$ is typically bigger than 160 MeV, so that $`\stackrel{~}{W}_1\stackrel{~}{Z}_1\pi `$ is always allowed and the chargino typically decays within the detector . The chargino would then manifest itself only as missing energy, unless the decay length is a few tens of cm, so that the chargino track can be established in the detector. The track would then seem to disappear since the presence of the soft pion would be very difficult to detect. Some parameter regions with $`m_{\stackrel{~}{W}_1}m_{\stackrel{~}{Z}_1}<m_{\pi ^\pm }`$ may be possible; in this case, the chargino would mainly decay via $`\stackrel{~}{W}_1\stackrel{~}{Z}_1e\nu `$ and its decay length (depending on the mass difference) would be typically larger than several metres. It would then show up via a search for long-lived charged exotics. There have been a number of alternative suggestions to cure the negative slepton mass squared problem. Generally, these require the introduction of additional fields at energy scales higher than the weak scale. The mass spectrum in these scenarios differs from that of the minimal AMSB model sketched above, and characteristic features such as $`m_{\stackrel{~}{W}_1}m_{\stackrel{~}{Z}_1}`$ and $`m_{\stackrel{~}{\mathrm{}}_L}m_{\stackrel{~}{\mathrm{}}_R}`$ need not occur. These models are not hard wired into ISAJET, but can be generated using the $`NUSUG_i`$ inputs at a scale dictated by $`SSBCSC`$; in this case, the user must perform the calculation of the SSB masses of MSSM particles. ## X Minimal Gaugino Mediation Very recently, Schmaltz and Skiba have proposed a model based on extra dimensions with branes, which is claimed to provide novel solutions to the SUSY flavour and $`CP`$ problems. Within their framework, chiral supermultiplets of the observable sector reside on one brane whereas the SUSY breaking sector is confined to a different brane . Gravity and gauge superfields propagate in the bulk, and hence, directly couple to fields on both the branes. As a result of their direct coupling to the SUSY breaking sector, gauginos acquire a mass. The scalar components of the chiral supermultiplets, however, can acquire a SUSY breaking mass only via their interactions with gauginos (or gravity) which feel the effects of SUSY breaking: as a result, these masses are suppressed relative to gaugino masses, and may be neglected in the first approximation. The same is true for the $`A`$\- and $`B`$-parameters. In the specific realization, to preserve the success of the unification of gauge couplings, it is assumed that there is grand unification (both $`SU(5)`$ and $`SO(10)`$ are discussed), and further, that the compactification scale $`M_c`$ below which there are no Kaluza-Klein excitations, is larger than $`M_{GUT}`$. Furthermore, since light bulk fields have flavor-blind interactions by construction, it is argued that the scale $`M_cM_{Planck}/10`$ in order to sufficiently suppress flavour violating scalar couplings (due to heavy bulk fields) that would be generically present. Based on the discussion in the previous paragraph, they take the boundary conditions for the soft SUSY breaking parameters of the MSSM to be, $`m_0=A_0=B_0=0`$ at the scale $`M_c`$, and argue that the spectrum is completely specified by the parameter set, $$\mu ,m_{1/2},M_c$$ (106) where it is the grand unification assumption that leads to a universal gaugino mass above $`Q=M_{GUT}`$. They refer to this as the Minimal Gaugino Mediation (MGM) model. The parameters $`m_{1/2}`$ and $`\mu `$ should be comparable, and are chosen to be $`M_{Weak}`$. The REWSB constraints fix $`\mu ^2`$, while the requirement $`B_0=0`$ fixes $`\mathrm{tan}\beta `$. In Ref. it is shown that if $`M_cM_{Planck}/10`$ $`\mathrm{tan}\beta `$ lies between $`12`$ and $`18`$ (12-25) for the $`SU(5)`$ ($`SO(10)`$) model with $`\mathrm{𝟓}+\overline{\mathrm{𝟓}}`$ ($`\mathrm{𝟏𝟔}+\overline{\mathrm{𝟏𝟔}}`$) Higgs supermultiplets in addition to the usual adjoint Higgs multiplet. The LSP may be the stau, the lightest neutralino or the gravitino. However, the latter has a weak scale mass, and as in the mSUGRA framework, is irrelevant for collider phenomenology. Our purpose here is to outline how to generate sample spectra in this framework using ISAJET , and examine some issues that have not been discussed in Ref. . For definiteness, we will choose the $`GUT`$ group to be $`SU(5)`$. This model is then a special case of our discussion in Sec. II, except that the SSB parameters now “unify” at the scale $`M_c`$ rather than $`M_P`$ (where they take on the special values). Our first observation is that the allowed range of $`\mathrm{tan}\beta `$ seems incompatible with $`\mathrm{tan}\beta 30`$ required for the unification of the $`\tau `$ and $`b`$ Yukawa couplings<sup>\</sup><sup>\</sup>\Another possibility is the inclusion of a bilinear $`R`$-parity violating term in the tau sector. In this case, $`b`$-$`\tau `$ Yukawa unification can be achieved at smaller values of $`\mathrm{tan}\beta `$.. For this reason, and also because the prediction for $`\mathrm{tan}\beta `$ could depend on how the $`\mu `$ problem might be solved, we will ignore the $`B_0=0`$ condition and treat $`\mathrm{tan}\beta `$ as a phenomenological parameter.<sup>††</sup><sup>††</sup>††Moreover, if Higgs fields are also allowed to propogate in the bulk , we would expect $`B_0m_{1/2}m_{H_u}m_{H_d}`$. For our analysis, we modify the model parameters <sup>‡‡</sup><sup>‡‡</sup>‡‡There are other coupling constants involving GUT scale physics, but we will see that these do not significantly change the spectrum. to, $$m_{1/2},M_c,\mathrm{tan}\beta ,sign(\mu ).$$ (107) As before, the user will have to obtain the values of the SSB parameters at $`Q=M_{GUT}`$ using the RG equations of Sec. II, and input these into ISAJET for generating mass spectra and/or collider events as desired. As shown in Table LABEL:mgmtable, we fix $`m_{1/2}`$ at the GUT scale and $`\mathrm{tan}\beta `$ at the weak scale. In Fig. 6, we show the evolution of the various SSB parameters of the MSSM, starting with the MGM boundary conditions. Here, the unified gaugino mass is taken to be 300 GeV at $`Q=M_{GUT}`$. The compactification scale is taken to be $`M_c=10^{18}`$ GeV, and other parameters are fixed to be the same as in Fig. 1. We see that RG evolution results in GUT scale scalar masses and $`A`$-parameters that are substantial fractions of $`m_{1/2}`$; i.e. although we have no-scale boundary conditions at the scale $`M_c`$, there are substantial deviations from these at $`M_{GUT}`$. While the inter-generation splitting is small, the splittings between the 5 and the 10 dimensional matter multiplets, as well as between these and the Higgs multiplets is substantial. In Fig. 7, we show the variation of several SSB masses at the scale $`Q=M_{GUT}`$ with the unified gaugino mass $`m_{1/2}`$ for the same values of other parameters as in the previous figure. These masses then serve as inputs for ISAJET. We note that if $`m_{1/2}`$ is too small, the no-scale like boundary conditions lead to incorrect electroweak symmetry breaking or $`m_{\stackrel{~}{\tau }_1}<m_{\stackrel{~}{Z}_1}`$. For instance, if $`\mathrm{tan}\beta =35`$ (this allows unification of the $`b`$ and $`\tau `$ Yukawa couplings) with other parameters as in Fig. 7, only values of $`m_{1/2}`$ larger than 275 GeV are phenomenologically acceptable. In Table LABEL:mgmtable we show a sample spectrum for this model. We choose $`m_{1/2}=300`$ GeV, $`\mathrm{tan}\beta =35`$ and other parameters as in Fig. 7. The spectrum is not unlike that in the mSUGRA framework with small $`m_0`$ so that sleptons are relatively light and squarks are lighter than the gluino. The chargino and $`\stackrel{~}{Z}_2`$ almost exclusively decay via $`\stackrel{~}{W}_1\stackrel{~}{\tau }_1\nu _\tau `$ and $`\stackrel{~}{Z}_2\stackrel{~}{\tau }_1\tau `$, respectively, so that cascade decays of gluinos and squarks will lead to multi-jet plus multi-tau events, with (soft) leptons as daughters of the tau. Except for $`h`$, this scenario is probably beyond the reach of the Tevatron, but it should be straightforward to study $`\stackrel{~}{\mathrm{}}_R`$ and $`\stackrel{~}{\tau }_1`$, and probably also detect $`\stackrel{~}{W}_1`$ and $`\stackrel{~}{\nu }`$, at the NLC. At the LHC a variety of signals should be present. We have also examined how the mass spectrum changes with variation of the superpotential couplings $`\lambda `$ and $`\lambda ^{}`$. These couplings cannot be too large in order that they remain perturbative up to $`M_c`$. For variation in this range, we found that $`m_{10}(GUT)`$ and $`m_5(GUT)`$ were insensitive to the choice of these couplings, while the GUT scale values of $`m__1`$ and $`m__2`$ as well as $`A_t`$ and $`A_b`$ vary by about 20% over the entire range of $`\lambda `$ and $`\lambda ^{}`$ that we examined. The weak scale spectrum and the $`\mu `$ value are, however, insensitive to the choice of these parameters; this is presumably because $`m_{1/2}`$ is significantly larger than the scalar masses at the GUT scale, so that RG evolution between the GUT and weak scales, rather than from $`m_0`$, makes the bulk of the contribution to scalar masses. ## XI Models with non-universal soft terms due to 4–D superstring dynamics Soft supersymmetry breaking terms obtained from $`N=1`$ four–dimensional superstrings, in general, exhibit non-universality at the string scale , a notable exception being when the dilaton is the dominant source of SUSY breaking. The soft supersymmetry breaking terms are determined by the Kähler potential $`K`$ and the gauge kinetic functions $`f_a`$ of the effective supergravity theory obtained from the string. The Kähler potential depends on the hidden sector fields, the dilaton $`S`$ and the moduli $`T`$ (there could be several), and the observable sector fields $`C_i`$, and it has the form, $$K=\mathrm{log}(S+S^{})+K_0(T,T^{})+\stackrel{~}{K}_{ij}(T,T^{})C_iC^j.$$ (108) To avoid potential problems with FCNCs, we will assume that $`\stackrel{~}{K}_{ij}=\stackrel{~}{K_i}\delta _{ij}`$. In addition, the gauge kinetic function in any 4–dimensional superstring is given at tree level by $$f_a=k_aS$$ (109) where $`k_a`$ is the Kac–Moody level of the gauge factor $`G_a`$, with the entire group given by $`G=\mathrm{\Pi }_aG_a`$. The Kac–Moody levels are usually taken $`k_3=k_2=\frac{3}{5}k_1=1`$. Beyond the tree level, $`f_a`$ would in general also contain a dependence on the moduli fields. Supersymmetry is broken when the auxiliary $`F`$-terms of the hidden sector fields acquire vacuum expectation values ($`vev`$). A convenient way to parametrize the vevs (in the case of one modulus) is as follows $`F^S`$ $`=`$ $`\sqrt{3}Cm_{3/2}K_{S\overline{S}}^{1/2}\mathrm{sin}\theta e^{i\gamma _S}`$ (110) $`F^T`$ $`=`$ $`\sqrt{3}Cm_{3/2}K_{T\overline{T}}^{1/2}\mathrm{cos}\theta e^{i\gamma _T}`$ (111) where $`C`$ is a constant defined by $`C^2=1+V_0/3m_{3/2}^2`$, $`V_0`$ is the cosmological constant (the vev of the scalar potential), and $`m_{3/2}`$ is the gravitino mass. Here, $`\mathrm{sin}\theta `$ is the overlap between the goldstino and the fermionic component of the dilaton field. Therefore, $`\mathrm{sin}\theta =1`$ in the limit where the SUSY breaking is completely due to the dilaton: i.e. $`F_S`$ is the only relevant $`vev`$. The matrix $`K_{n\overline{m}}_n_{\overline{m}}K`$ is called the Kähler metric and $`\gamma _S`$ and $`\gamma _T`$ are possible complex phases. The soft masses for scalar particles are determined by the Kähler potential in Eq. (108) and are given by $$m_i^2=2m_{3/2}^2(C^21)+m_{3/2}^2C^2(1+N_i\mathrm{cos}^2\theta ),$$ (112) with $$N_i=\frac{3(\mathrm{log}\stackrel{~}{K}_i)_{T\overline{T}}}{(K_0)_{T\overline{T}}}.$$ We readily see that we can obtain non-universal scalar masses if $`\mathrm{cos}\theta `$ is different from zero. We mention that here we have for simplicity assumed that there is just one modulus field: multiple moduli are treated in Ref. . The gaugino masses are given by $$M_a=\frac{1}{2}(\mathrm{Re}f_a)^1F^m_mf_a=\sqrt{3}Cm_{3/2}(\mathrm{Re}f_a)^1k_a\mathrm{Re}Se^{i\gamma _S}\mathrm{sin}\theta ,$$ (113) where the gauge coupling constants are $`\mathrm{Re}f_a=1/g_a^2`$. In the last equality, we have used the fact that (at tree level) the gauge kinetic function in Eq. (109) depends only on the dilaton field $`S`$, so that the tree level gaugino masses are independent of the moduli sector. Model-dependent corrections to this may, however, be significant, particularly when dilaton contributions to SUSY breaking are small. Expressions for $`A`$-parameters may also be found in Ref. . These depend on additional parameters, and generically also on the unknown phases $`\gamma _S`$ and $`\gamma _T`$ (as well as on additional direction cosines in the multi-moduli case). For the single modulus case, the form of $`A`$ is given by,<sup>\</sup><sup>\</sup>\We have flipped the sign of $`A`$ to conform to our convention where the soft trilinear term is written as $`A_{ijk}f_{ijk}\stackrel{~}{C}_i\stackrel{~}{C}_j\stackrel{~}{C}_k`$ in the Lagrangian and not the scalar potential, with $`f_{ijk}`$ being the corresponding superpotential coupling. $$A_{ijk}=\sqrt{3}m_{3/2}C(e^{i\gamma _S}\mathrm{sin}\theta +e^{i\gamma _T}\omega _{ijk}(T,T^{})\mathrm{cos}\theta ),$$ (114) where $`\omega _{ijk}`$ depend on the Kähler and superpotentials. Fortunately, in many cases of interest, these model-dependent parameters either vanish or assume a simple form. We should mention that these expressions for the soft-SUSY breaking masses and $`A`$-parameters are valid for these parameters renormalized at the string scale. As always, these have then to be evolved down to the weak scale for use in phenomenological analysis. We now consider some special cases to illustrate the forms of (string scale) non-universality that may occur in this general framework. ### A Large–T limit of Calabi–Yau compactifications Because of the complexity of the world–sheet instanton and sigma model contributions, the general form of the Kähler potential of generic Calabi–Yau $`(2,2)`$ compactifications is not known. The gauge group is $`E_6\times E_8`$, with matter in the 27 dimensional representation of $`E_6`$. It is usual to analyze the large $`T`$ (in practice $`23<\|T\|<2030`$, large enough so that world sheet instanton contributions can be neglected, but not so large that string threshold corrections invalidate perturbation theory) limit of these theories. In this limit the Kähler potential takes a simple form : $$K=\mathrm{log}(S+S^{})3\mathrm{log}(T+T^{})+\underset{i}{}\frac{\|C_i\|^2}{T+T^{}},$$ (115) and the gauge kinetic function is given by Eq. (109) at tree level. In this case the gaugino mass is $$m_{1/2}=\sqrt{3}Cm_{3/2}\mathrm{sin}\theta e^{i\gamma _S},$$ (116) while Eq. (112) for the scalar masses reduces to, $$m_0^2=m_{3/2}^2C^2\mathrm{sin}^2\theta +2m_{3/2}^2(C^21)$$ (117) which simplifies even further if the cosmological constant vanishes ($`C=1`$). Notice that we find universality of soft scalar masses, even though we are not in the dilaton dominated SUSY breaking scenario. In the $`C=1`$ case, we see that $`\|m_{1/2}\|=\sqrt{3}m_0`$, so that the gaugino mass always exceeds the scalar mass at the string scale. This relation obviously puts a significant constraint on SUSY phenomenology. Since this is a special case of the mSUGRA scenario whose phenomenological implications have been discussed at length in the literature, we will not mention this any further. There are, however, arguments in the literature that suggest that the observed cosmological constant (which is bounded to be smaller than $`(3meV)^4`$) may not be directly connected to $`V_0`$; then, $`C`$ could differ from unity, and the gaugino mass may (depending on the value of $`C`$ and the goldstino angle $`\theta `$) be even smaller than $`m_0`$, but for an appreciable effect, $`C1`$ would have to deviate by many orders of magnitude<sup>\†</sup><sup>\†</sup>\†It should be appreciated that even $`C=1.1`$ is an enormous value relative to the bound $`C110^{87}`$ that we would get if we took $`V_0`$ to be related to the observed value of $`\mathrm{\Lambda }`$. from the bound that would have resulted assuming $`V_0`$ was the observed cosmological constant. Finally, in this limit, the parameters $`\omega _{ijk}`$ in Eq. (114) vanish so that $$A_{ijk}=\sqrt{3}m_{3/2}Ce^{i\gamma _S}\mathrm{sin}\theta .$$ In the single modulus large $`T`$ case that we have been discussing, effects of the sigma–model loop contribution and the non–perturbative instanton contribution to the Kähler potential are known . We still obtain universality of soft SUSY breaking parameters, with gaugino masses given by Eq. (116) and scalar masses and the $`A`$ parameter (in the case $`C=1`$) modified to, $$m_0^2=m_{3/2}^2\left[1\mathrm{cos}^2\theta \left(1\mathrm{\Delta }(T,T^{})\right)\right],$$ (118) and $$A=\sqrt{3}m_{3/2}\left[e^{i\gamma _S}\mathrm{sin}\theta +\omega (T,T^{})e^{i\gamma _T}\mathrm{cos}\theta \right].$$ (119) Here $`\mathrm{\Delta }`$ and $`\omega `$ corresponds to the sigma–model and instanton contributions (the latter are negligible): the numerical values of these are model dependent, but $`\mathrm{\Delta }0.4`$ and $`\omega =0.17`$ have been quoted for a typical model. Notice that although these corrections do not lead to non-universality, we lose the earlier prediction $`m_{1/2}=\sqrt{3}m_0`$: now, the soft scalar mass may even exceed the corresponding gaugino mass if $`\mathrm{cos}^2\theta `$ is sufficiently large. ### B General Calabi–Yau compactifications There is no reason to believe that there is just a single modulus field $`T`$. In the multi–moduli case the parametrization of the $`vevs`$ of the moduli in Eq. (111) is modified to , $$F^{T_i}=\sqrt{3}Cm_{3/2}K_{T_i\overline{T}_i}^{1/2}\mathrm{cos}\theta \mathrm{\Theta }_ie^{i\gamma _{T_i}},$$ (120) where we have assumed the Kähler metric to be diagonal to avoid any FCNC problems. Here $`\mathrm{\Theta }_i`$ are direction cosines that parametrize the direction of the $`vev`$ in moduli space. Indeed the more general case of an off-diagonal metric has also been examined in Ref. where a more general parametrization of the $`vevs`$ of the moduli may be found. In this general case, the scalar masses are non-diagonal and the mass squared matrix assumes the form, $$m_{ij}^2=m_{3/2}^2\left[\delta _{ij}\mathrm{cos}^2\theta \left(\delta _{ij}\mathrm{\Delta }_{ij}(T_k,T_k^{})\right)\right]$$ (121) where $`\mathrm{\Delta }_{ij}`$ depends on the moduli and on the direction of the $`vev`$ in the moduli space. Notice that the model-dependent $`\mathrm{\Delta }_{ij}`$ would be strongly constrained by experimental data on flavour mixing. We are, however, not aware of a realistic model in which such constraints may be analyzed. We also note that the presence of a (even diagonal) matrix $`\mathrm{\Delta }`$ in Eq. (121) would be a source of non-universality of scalar masses. ### C Orbifold models with large threshold corrections An example of such a model is the so-called $`O`$-$`I`$ model discussed by Brignole et al. . In orbifold compactifications the coefficient $`\stackrel{~}{K}_{ij}`$ which determines the soft masses has the form $`(T+T^{})^{n_i}`$, where $`n_i`$ is the modular weight of the matter field $`C_i`$. The Kähler potential is in this case: $$K=\mathrm{log}(S+S^{})3\mathrm{log}(T+T^{})+\underset{i}{}\|C_i\|^2(T+T^{})^{n_i}$$ (122) Gauge unification in good agreement with low energy data is achieved by assigning the following modular weights for the massless fields: $`n_Q=n_D=1`$, $`n_U=2`$, $`n_L=n_E=3`$, and $`n_{H_d}+n_{H_u}=5`$ or $`4`$, together with a large value for the modulus field, $`\mathrm{T}16`$, which then results in large threshold corrections. Under these conditions the gaugino masses are non-universal at the string scale: $`M_1`$ $`=`$ $`1.18\sqrt{3}m_{3/2}\left[\mathrm{sin}\theta +2.9\times 10^2(B_1^{}/k_1)\mathrm{cos}\theta \right]`$ (123) $`M_2`$ $`=`$ $`1.06\sqrt{3}m_{3/2}\left[\mathrm{sin}\theta +2.9\times 10^2(B_2^{}/k_2)\mathrm{cos}\theta \right]`$ (124) $`M_3`$ $`=`$ $`1.00\sqrt{3}m_{3/2}\left[\mathrm{sin}\theta +2.9\times 10^2(B_3^{}/k_3)\mathrm{cos}\theta \right]`$ (125) where $`B_a^{}b_a^{}k_a\delta _{GS}`$ are given by $`B_1^{}=18k_1\delta _{GS}`$, $`B_2^{}=8k_2\delta _{GS}`$, and $`B_3^{}=6k_3\delta _{GS}`$ if $`n_{H_d}+n_{H_u}=5`$, and $`k_a`$ as specified previously. Here, the parameter $`\delta _{GS}`$ is a model dependent negative integer and $`m_{3/2}`$ and $`\theta `$ are the gravitino mass and the goldstino angle as before. To obtain Eqs. (124), it is assumed that string threshold corrections lead to an apparent unification of the couplings at the “$`GUT`$ scale” rather than at the string scale. Of course, since there is no $`GUT`$ these couplings continue to evolve and diverge when evolved from the “$`GUT`$ scale” to the one order of magnitude larger string scale. The coefficients in front of the gaugino mass formulae reflect just this difference in the gauge couplings at the string scale. In other words, if $`\mathrm{sin}\theta =1`$, gaugino masses (while slightly different at the string scale) would be universal at $`Q=M_{GUT}`$: non-universality of GUT scale gaugino masses occurs only due to the loop correction proportional to $`\mathrm{cos}\theta `$ in Eqs. (124). Finally, we note that if $`n_{H_d}+n_{H_u}=4`$, the gaugino masses are obtained from Eqs. (124) by modifying the coefficients $`B_i^{}`$ to $`B_1^{}=17k_1\delta _{GS}`$ and $`B_2^{}=7k_2\delta _{GS}`$ while $`B_3^{}`$ does not change. The string scale scalar masses and $`A`$ parameters depend on the modular weights, and (assuming zero cosmological constant) are given by, $`m_Q^2=m_D^2`$ $`=`$ $`m_{3/2}^2\left[1(1\delta _{GS}\times 10^3)^1\mathrm{cos}^2\theta \right],`$ (126) $`m_U^2`$ $`=`$ $`m_{3/2}^2\left[12(1\delta _{GS}\times 10^3)^1\mathrm{cos}^2\theta \right],`$ (127) $`m_L^2=m_E^2`$ $`=`$ $`m_{3/2}^2\left[13(1\delta _{GS}\times 10^3)^1\mathrm{cos}^2\theta \right],`$ (128) and $$A_{ijk}=\sqrt{3}m_{3/2}\mathrm{sin}\theta \pm m_{3/2}\mathrm{cos}\theta (1\delta _{GS}\times 10^3)^{1/2}(3+n_i+n_j+n_k),$$ (129) where the terms with $`\delta _{GS}`$ come from radiative corrections, and the sign ambiguity reflects the possible relative phase between $`\gamma _S`$ and $`\gamma _T`$ (we take the $`A`$-parameters to be real). Note that if $`\mathrm{sin}\theta =1`$, the scalar masses and $`A`$-parameters are universal at the string scale: RG evolution would then introduce a small non-universality at $`M_{GUT}`$. In Fig. 8 we plot different soft masses at the string scale as a function of $`\mathrm{sin}\theta `$ in the $`O`$-$`I`$ model. There is a sign ambiguity since $`\mathrm{cos}\theta `$ could be negative. We have chosen $`\mathrm{cos}\theta >0`$ and fixed $`m_{3/2}=200`$ GeV, $`\delta _{GS}=0`$, and, for the evaluation of gaugino masses, $`n_{H_u}+n_{H_d}=5`$. We set the phases $`\gamma _S`$ and $`\gamma _T`$ to be zero. Scalar masses are universal in the dilaton dominated scenario and radiative corrections do not spoil this universality. On the contrary, gaugino masses are not universal at $`\mathrm{sin}\theta =1`$, but as explained above, there is (approximate) universality at $`M_{GUT}`$. Values of $`\mathrm{cos}\theta \stackrel{>}{}\mathrm{\hspace{0.25em}1}/\sqrt{3}`$ ($`\mathrm{sin}\theta \stackrel{<}{}\mathrm{\hspace{0.25em}0.8}`$) yield negative slepton soft squared masses and may be unacceptable;<sup>\‡</sup><sup>\‡</sup>\‡It may be possible to have these squared masses negative at a high scale as long as they are positive near the weak scale. hence the dilaton field is necessarily the most important source of SUSY breaking. Except close to the lowest acceptable values of $`\theta `$, deviations from universality in the scalar sector are thus limited. To facilitate simulation of such a scenario, we have introduced into ISAJET versions $`7.50`$ the “SUSY Boundary Condition Scale” ($`SSBCSC`$ keyword) option into ISAJET that allows the user to input a chosen scale $`Q_{max}`$ up to which the MSSM is assumed to be valid. The values of SUSY breaking masses and $`A`$-parameters of the MSSM as given by any theory valid at the scale beyond $`Q_{max}`$ would then be used as inputs to ISAJET, which would then evolve them down to the weak scale and generate SUSY events as usual. For the case at hand, $`Q_{max}`$ would be the string scale, and the gaugino masses, scalar masses and $`A`$-parameters as given by Eqs. (124) - (129), the boundary conditions for the RGE. We stress, however, that $`Q_{max}`$ need not be larger than $`M_{GUT}`$. For instance, in $`SO(10)`$ models, $`Q_{max}`$ would be the mass of the right-handed neutrino, or in $`E_6`$ models, the mass scale where the additional particles in the 27 dimensional representation and any extra $`Z^{}`$ bosons all decouple, leaving the MSSM spectrum. We give an example of the SUSY spectrum in the O-I scenario in Table LABEL:o-scenarios. In this example, we have fixed $`\mathrm{tan}\beta =4`$, $`\mathrm{sin}\theta =0.85`$ (with $`\mathrm{cos}\theta >0`$) and have taken $`n_{H_u}=3`$, with other parameters as in Fig. 8. Since the value of $`B`$ depends on how $`\mu `$ is generated, we have treated $`\mathrm{tan}\beta `$ as a free parameter, and eliminated $`B`$ in its favour, using the constraints given by radiative electroweak symmetry breaking. We fix the string scale to be $`4\times 10^{17}`$ GeV. Despite the fact that string scale slepton masses are considerably smaller than those of squarks (see Fig. 8), the spectrum is qualitatively very similar to that in the mSUGRA framework with $`m_{\stackrel{~}{q}}m_{\stackrel{~}{g}}`$. ### D Orbifold models with small threshold corrections In the $`O`$-$`I`$ model, $`\mathrm{sin}\theta `$ was restricted to be large, so that the parameters of phenomenological interest were qualitatively similar to the mSUGRA scenario. To allow a wider range of $`\mathrm{sin}\theta `$ we consider a model where all the modular weights are $`1`$. As noted in Ref. string threshold corrections cannot account for gauge coupling unification, which has then to be attributed to some different physics. Unlike the $`O`$-$`I`$ model where a large value of $`\mathrm{R}eT`$ was needed to accommodate coupling constant unification, we will, following Brignole et al. use $`\mathrm{R}eT1.2`$ and refer to this as the $`O`$-$`II`$ model. As before, the gaugino masses are non-degenerate at the string scale (again, for $`\mathrm{sin}\theta =1`$, these would be universal at “$`M_{GUT}`$”) and given by: $`M_1`$ $`=`$ $`1.18\sqrt{3}m_{3/2}\left[\mathrm{sin}\theta +4.6\times 10^4(B_1^{\prime \prime }/k_1)\mathrm{cos}\theta \right]`$ (130) $`M_2`$ $`=`$ $`1.06\sqrt{3}m_{3/2}\left[\mathrm{sin}\theta +4.6\times 10^4(B_2^{\prime \prime }/k_2)\mathrm{cos}\theta \right]`$ (131) $`M_3`$ $`=`$ $`1.00\sqrt{3}m_{3/2}\left[\mathrm{sin}\theta +4.6\times 10^4(B_3^{\prime \prime }/k_3)\mathrm{cos}\theta \right]`$ (132) with $`B_1^{\prime \prime }=11k_1\delta _{GS}`$, $`B_2^{\prime \prime }=1k_2\delta _{GS}`$, and $`B_3^{\prime \prime }=3k_3\delta _{GS}`$. On the other hand, the scalar masses ($`V_0=0`$) and $`A`$ parameters are all degenerate and equal to $$m_Q^2=m_D^2=m_U^2=m_L^2=m_E^2=m_{3/2}^2\left[1(1\delta _{GS}\times 10^3)^1\mathrm{cos}^2\theta \right],$$ (133) and $$A_{ijk}=\sqrt{3}m_{3/2}\mathrm{sin}\theta ,$$ (134) at the string scale. If $`\mathrm{sin}\theta 1`$ the spectra should be the same as in the $`O`$-$`I`$ model discussed previously. For smaller values of $`\mathrm{sin}\theta `$, the degeneracy in the string scale scalar masses still remains. The most important difference between the two scenarios is that very small values of $`\mathrm{sin}\theta `$ are now permitted; i.e. the dilaton contribution need not necessarily dominate SUSY breaking. If $`\mathrm{sin}\theta `$ is very small so that the $`\mathrm{cos}\theta `$ terms are the dominant contributions to the gaugino mass, we see that (depending on the value of $`\delta _{GS}`$) the GUT scale gluino mass may be much smaller than the corresponding electroweak gaugino masses. Indeed it is possible to arrange scenarios where the gluino is the LSP. The additional parameters also allow the possibility $`M_1M_2`$ so that the lighter chargino and the two lighter neutralinos (and sometimes also the gluino) are all very degenerate. Such scenarios pose interesting experimental challenges . In Fig. 9 we illustrate the gaugino and scalar soft masses at the string scale as a function of $`\mathrm{sin}\theta `$ in the $`O`$-$`II`$ model. Again, we take $`m_{3/2}=200`$ GeV, $`\mathrm{cos}\theta >0`$, and ignore all phases. We choose $`\delta _{GS}=5`$. The masses decrease as $`\mathrm{sin}\theta `$ decrease but they do not vanish at $`\mathrm{sin}\theta =0`$ due to one–loop effects. Of course, for very small values of $`\mathrm{sin}\theta `$ phenomenological considerations require $`m_{3/2}`$ to be significantly larger. In the extreme case of moduli-dominated SUSY breaking, gaugino masses can be smaller than scalar masses, but generally speaking scalar masses are smaller than gaugino masses at the unification scale. In the last three columns of Table LABEL:o-scenarios we illustrate three examples of $`O`$-$`II`$ model spectra. Again, we fix $`\mathrm{tan}\beta =4`$, $`\mu >0`$ and take $`\delta _{GS}=5`$, to be in the region which can potentially yield roughly equal masses for all the MSSM gauginos. First, we choose an $`O`$-$`II`$ scenario with parameters close to those of the $`O`$-$`I`$ model in the previous column: $`m_{3/2}=200`$ GeV and $`\mathrm{sin}\theta =0.85`$. This is the “typical” case for such a model. In this case, the $`\mathrm{sin}\theta `$ terms in Eq. (131) completely dominate, and the resulting spectrum is again very similar to that in the mSUGRA framework (with $`m_{\stackrel{~}{q}}m_{\stackrel{~}{g}}`$). In the next column, we show a spectrum for the case $`\mathrm{sin}\theta =0`$, the extreme caseIt does not matter whether we take $`\theta =0`$ or $`\theta =\pi `$ since the sign of the gaugino mass has no import for physics. of moduli-dominated SUSY breaking. Here, because of the small coefficient $`4.6\times 10^4`$ in the expressions for gaugino masses, we have to choose $`m_{3/2}`$ to be large. We fix $`m_{3/2}=60`$ TeV. For this case, we have taken $`m_t=180`$ GeV, since we found that electroweak symmetry was not brokenThe scalars start at a very large mass at the string scale, and the top Yukawa is not large enough to drive a Higgs mass squared eigenvalue negative at the scale $`Q=\sqrt{m_{\stackrel{~}{t}_L}m_{\stackrel{~}{t}_R}}`$ where the effective potential is evaluated in ISAJET . We should mention that this is sensitive to the top mass radiative corrections that have been included in ISAJET versions $`7.48`$. These radiative corrections decrease the top Yukawa coupling by a few percent, and in this case, this is just sufficient to preclude electroweak symmetry breaking. for $`m_t=175`$ GeV. Since the (common) string-scale scalar mass is much bigger than the corresponding gaugino masses, the scalars are all roughly degenerate, and their spectrum is close to that of the corresponding mSUGRA spectrum with $`m_{\stackrel{~}{q}}m_\stackrel{~}{\mathrm{}}m_{\stackrel{~}{g}}`$ (i.e. $`m_0m_{1/2}`$). The gluinos, charginos and neutralino spectrum is quite different from that in the mSUGRA model: even though the lighter chargino and neutralinos are gaugino-like, $`m_{\stackrel{~}{W}_1}=m_{\stackrel{~}{Z}_2}=m_{\stackrel{~}{Z}_1}=0.7m_{\stackrel{~}{g}}`$. This is because by choosing $`\delta _{GS}`$ we can adjust $`M_1:M_2:M_3`$ at the string scale. By a careful adjustment of parameters the gluino mass can even be brought closer to the chargino and neutralino masses. Experiments at the Fermilab Tevatron may be sensitive to this scenario. To emphasize that the novel scenarios shown in Ref. obtain only for a very limited range of parameters, in the last column we show the spectrum for $`\mathrm{sin}\theta =0.005`$ (with $`\mathrm{cos}\theta >0`$), with all other parameters (including $`m_t`$) as for the $`\mathrm{sin}\theta =0`$ case. We see that even for this tiny value of $`\mathrm{sin}\theta `$, the $`\mathrm{sin}\theta `$ terms in Eq. (131) are comparable to (or even dominate) the $`\mathrm{cos}\theta `$ terms, and the spectrum is qualitatively different. While the sfermions are once again extremely heavy, the gluino, chargino and neutralino masses are now approximately as in the mSUGRA framework. Sparticle detection in this scenario would only be possible at the LHC. Our purpose in showing this (possibly unacceptably heavy) spectrum is only to emphasize the qualitative difference from the $`\mathrm{sin}\theta =0`$ case. Of course, if $`m_{3/2}`$ is chosen to be 15 TeV, many more sparticles would be in the accessible range, but the spectrum would then be much like the canonical mSUGRA case with large $`m_0`$. ## XII Models with non-universal soft terms due to M–theory dynamics It was proposed that M–theory, i.e., an 11–dimensional supergravity on a manifold where two $`E_8`$ gauge multiplets are restricted to the two 10–dimensional boundaries, is equivalent to the strong coupling limit of $`E_8\times E_8`$ heterotic string theory . It may be argued that M–theory is a better candidate than the weakly coupled string to explain low energy physics and unification. After compactifying the 11–dimensional M–theory, a 4–dimensional effective theory emerges which can reconcile the reduced Planck scale $`M_P2.4\times 10^{18}`$ GeV, the grand unification scale $`M_{GUT}3\times 10^{16}`$ GeV, and $`\alpha _{GUT}`$, in a way that the weakly coupled heterotic string theory cannot. An interesting feature of the 4–dimensional effective SUGRA is that, in first approximation, the gauge kinetic function, the superpotential, and the Kähler potential do not change when moving from the weakly coupling heterotic string case to the M–theory case by changing the value of the modulus field. ### A One modulus case Supersymmetry is broken when the auxiliary components of the dilaton field $`S`$ and the modulus field $`T`$ aquire non–zero $`vevs`$, as discussed in the last Section. The low energy effective supergravity theory obtained from a specific Calabi-Yau compactification in M–theory is a Yang-Mills gauge theory with $`E_6`$ as the gauge group. The gauge kinetic function is given by, $$f_{E_6}=S+\alpha T$$ (135) where $`\alpha `$ is an integer while the corresponding Kähler potential is $$K=\mathrm{log}(S+S^{})3\mathrm{log}(T+T^{})+\left[\frac{3}{T+T^{}}+\frac{\alpha }{S+S^{}}\right]\underset{i}{}\|C_i\|^2,$$ (136) where $`C_i`$ again denote the observable fields. Adopting the same parametrization as in Eq. (111) above, we find that with the Kähler potential of Eq. (136) the soft SUSY breaking parameters are universal and given by, $$m_{1/2}=\frac{\sqrt{3}Cm_{3/2}}{1+x}\left[\mathrm{sin}\theta e^{i\gamma _S}+\frac{x}{\sqrt{3}}\mathrm{cos}\theta e^{i\gamma _T}\right],$$ (137) $`m_0^2`$ $`=`$ $`m_{3/2}^2(3C^22){\displaystyle \frac{3C^2m_{3/2}^2}{(3+x)^2}}[x(6+x)\mathrm{sin}^2\theta +(3+2x)\mathrm{cos}^2\theta `$ (139) $`2\sqrt{3}x\mathrm{sin}\theta \mathrm{cos}\theta \mathrm{cos}(\gamma _S\gamma _T)],`$ and $$A=\frac{\sqrt{3}Cm_{3/2}}{3+x}\left[(3x)\mathrm{sin}\theta e^{i\gamma _S}+\sqrt{3}x\mathrm{cos}\theta e^{i\gamma _T}\right],$$ (140) where, $$x\frac{\alpha (T+T^{})}{S+S^{}},$$ (141) The range of $`x`$ is $`0x1`$. Note that in the weak coupling limit $`x0`$, we recover from Eqs. (137), (139) and (140) the gaugino and scalar masses as well as the $`A`$-parameter in the large $`T`$–limit of Calabi–Yau compactifications in Eqs. (116) and (117) respectively. In Fig. 10 we show the dependence on the goldstino angle of the universal gaugino and scalar masses, $`m_{1/2}`$ and $`m_0`$ respectively. We consider zero cosmological constant and three values of $`x`$. The strong coupling limit corresponds to $`x=1`$ and for comparison $`x=0.5`$ and $`x=0`$ are also plotted. We remind the reader that the soft parameters obtained above are for an $`E_6`$ gauge theory. In order to obtain a realistic low energy theory, we have to know how the symmetry group is reduced to the MSSM gauge group, which in turn will depend on the details of the theory at the high scale. It is possible that there may be additional $`TeV`$ scale supermultiplets in the particle spectrum, or even extra gauge bosons. Moreover, depending on how $`E_6`$ breaks to $`SU(3)\times SU(2)\times U(1)`$, additional $`D`$-term contributions (see Sec. V) which break the universality of scalar masses may also be present. ### B Multi-moduli case As before, the situation in the multi-moduli case can be more complicated. A toy example with three moduli fields and three observable fields has been considered in Ref. . The Kähler potential and gauge kinetic function of the effective theory is written as, $$K=\mathrm{log}(S+S^{})\underset{j=1}{\overset{3}{}}\mathrm{log}(T_j+T_j^{})+\left[2+\frac{2}{3}\underset{j=1}{\overset{3}{}}\frac{\alpha _j(T_j+T_j^{})}{S+S^{}}\right]\underset{i=1}{\overset{3}{}}\frac{\|C_i\|^2}{T_i+T_i^{}}.$$ (142) and $$f_a=S+\underset{i=1}{\overset{3}{}}\alpha _iT_i$$ (143) This then yields a universal mass for the gaugino and a universal $`A`$-parameter, but non-universal masses (and no mixing) for the scalars. While the gaugino and scalar masses as well as the $`A`$-parameter depend on the parameters and fields in the Kähler potential, the splitting $`\delta m^2`$ (between the scalars) appears to depend only on the orientation of the $`vevs`$ of the auxiliary components of the moduli and on the goldstino angle $`\theta `$. Since our focus is on sources of non-universality in realistic scenarios that can potentially be of phenomenological interest, we merely note that multiple moduli could be a source of non-universality of scalar masses, but do not exhibit results for this toy model here. ## XIII Concluding Remarks While weak scale supersymmetry is a well-motivated idea, the physical principles that fix the multitude of SUSY breaking parameters are not known. Without any sparticle signals to provide clues, we do not have any guidance as to what these might be. The scale of this new physics may be as low as a few hundred TeV as in models with low energy SUSY breaking mediated by gauge interactions, or as high as $`M_{GUT}M_P`$ as in frameworks where SUSY breaking is mainly mediated by gravity. Observable sparticle masses and mixing patterns, and via these weak scale SUSY phenomenology, are determined by the physics behind SUSY breaking and how this is communicated to the observable sector. Turning this around, measurement of sparticle properties may provide clues about physics at energy scales that would be inaccessible to experiments in the foreseeable future. Most early phenomenological analyses have been done within the framework of the mSUGRA model or the mSUGRA-motivated MSSM (where ad hoc relations between SSB parameters were assumed). In the last few years, phenomenological aspects of gauge-mediated SUSY breaking have also been examined in some detail. Both these models rest upon untested assumptions about physics at high energies. The good thing is that some of these assumptions will be directly testable if sparticles are discovered and their properties are measured . Nevertheless, it seems worthwhile to look at other viable alternatives for physics at energy scales much beyond the weak scale, with a view to see if there are direct ramifications for sparticle signals in future experiments. A serious study of this would entail SUSY simulation at colliders in a wide variety of models with features different from the mSUGRA paradigm, which is characterized by universality of SSB parameters at a scale $`QM_{GUT}`$. Our study represents a first step in this direction. Here, we have surveyed a number of proposals for high scale physics that lead to non-universality of the SSB parameters in the MSSM, which we regard as the effective theory at a sufficently low mass scale. These range from relatively minor modifications of the mSUGRA $`SU(5)`$ GUT model, where, e.g. unification of scalar masses and $`A`$-parameters is assumed to occur at $`M_P`$ (so that RG evolution induces some non-universality at $`Q=M_{GUT}`$), to major modifications involving conceptually new ideas for high scale physics (new hypercolour interactions, string physics) or the mediation of SUSY breaking to the observable sector (anomaly mediated SUSY breaking, gaugino mediated SUSY breaking). Other proposals that fall somewhat between these two extremes include models with larger unifying groups that naturally have additional non-universal contributions to scalar masses, or models where special boundary conditions on SSB parameters lead to unusual RG evolution and non-degeneracy of sparticle masses. For each of these scenarios, we have outlined the underlying physical ideas, delineated the parameter space in terms of which SUSY phenomenology might be analyzed, and discussed SUSY event generation using the simulation program ISAJET . A variety of improvements to the ISAJET program have been made to allow event generation in the models discussed in this paper. These improvements are characterized by ISAJET keyword inputs, including $`NUSUGi`$ for non-universal masses, $`SUGRHN`$ for models with a right-handed neutrino contribution, such as $`SO(10)`$, $`SSBCSC`$ for user choice as to when the MSSM becomes valid, and $`AMSB`$ for anomaly-mediated SUSY breaking models. Where possible, we present sample spectra, and allude to the important phenomenological differences from the reference mSUGRA framework. A detailed phenomenological analysis of each one of these scenarios is beyond the scope of the present work. Our hope though is that this study will facilitate and spur such analyses. Except for unusual cases where extreme degeneracies between sparticle masses result , we do not expect the reach of various future facilities (expressed in terms of physical sparticle masses) to qualitatively differ between the various scenarios. However, a careful examination of these will help us assess what we can hope to learn about high scale physics if sparticles are discovered and their properties measured. Careful examination of physical implications of a variety of viable alternatives for the underlying theory will also help increase our understanding of the sort of analyses that might be needed to discriminate between these. In view of the potential pay-off, we believe that such studies will be very worthwhile. ###### Acknowledgements. We thank B. de Carlos, M. Drees, L. Ibañez, P. Mercadante, C. Muñoz, and T. Li for helpful communications, and M. Drees again for valuable comments on the manuscript. This research was supported in part by the U. S. Department of Energy under contract number DE-FG02-97ER41022 and DE-FG-03-94ER40833. M.A.D. was also supported in part by CONICYT grant 1000539.
warning/0002/quant-ph0002016.html	ar5iv	text	# Three-Qubit Gate Realization Using Single Quantum Particle ## 1 Introduction Toffoli proved that a classical reversible computer can be constructed using a universal three-bit gate CCNOT (“controlled controlled NOT” or “double controlled NOT”). Deutsch has shown that in quantum information theory the generalized CCNOT gate - CCUT gate (“controlled controlled unitary transformation” or “double controlled unitary transformation”) \- also is a universal one. However to realize CCUT is not an easy task since there is no three body interaction in Nature. A way round was found Barenco et al. : it was proved that CCUT can be realized using five two qubit CUT (“controlled unitary transformation”) gates. Such an approach was used in the NMR implementation of Toffoli gate \[4, Eq. (38) and (39)\]. The virtual spin formalism developed in is used in the present paper to physically implement CCUT in a simplest possible way - using one pulse and one quantum particle. ## 2 The quantum information notation for gates of system of three real spins-$`1/2`$ In quantum information the gates for the three real spins-1/2 system usually are written in Hilbert space $`\mathrm{\Gamma }`$ which is a direct product $`\mathrm{\Gamma }=\mathrm{\Gamma }_Q\mathrm{\Gamma }_R\mathrm{\Gamma }_S`$ of the Hilbert spaces of three real spins-1/2 $`Q`$, $`R`$, $`S`$. As a basis of $`\mathrm{\Gamma }`$ the following eight functions can be chosen which will be used later $$\begin{array}{cccc}\|0=\|000,& \|1=\|001,& \|2=\|010,& \|3=\|011,\\ \|4=\|100,& \|5=\|101,& \|6=\|110,& \|7=\|111,\end{array}$$ (1) where $`\|M=\|m_Q,m_R,m_S`$, for example, $`\|5=\|m_Q=+1/2,m_R=1/2,m_S=+1/2`$ and so on. Let us consider all possible gates with NOT operation. In three spin system there are three NOT gates - the gate, which invert one spin leaving two others unchanged, $`NOT_Q`$ $`=`$ $`\|1m_Rm_S0m_Rm_S\|+\|0m_Rm_S1m_Rm_S\|,`$ $`NOT_R`$ $`=`$ $`\|m_Q1m_Sm_Q0m_S\|+\|m_Q0m_Sm_Q1m_S\|,`$ (2) $`NOT_S`$ $`=`$ $`\|m_Qm_R1m_Qm_R0\|+\|m_Qm_R0m_Qm_R1\|.`$ There are six CNOT gates. For example, when “master” spin R state controls “slave” spin Q state leaving spin S state unchanged $`CNOT_{RQ}`$ $`=`$ $`\|00m_S00m_S\|+\|11m_S01m_S\|+`$ (3) $`\|10m_S10m_S\|+\|01m_S11m_S\|.`$ Reverse gate $`CNOT_{QR}`$ $`=`$ $`\|00m_S00m_S\|+\|01m_S01m_S\|+`$ (4) $`\|11m_S10m_S\|+\|10m_S11m_S\|.`$ Other four CNOT gates can be written for pairs $`RS`$ and $`QS`$ by analogy. There are three CCNOT gates. For example, when two “master” spins $`Q`$ and $`R`$ control one “slave” $`S`$ one has $$\begin{array}{ccc}\hfill CCNOT_{Q,RS}& =& \|000000\|+\|001001\|+\hfill \\ & & \|010010\|+\|011011\|+\hfill \\ & & \|100100\|+\|101101\|+\hfill \\ & & \|111110\|+\|110111\|.\hfill \end{array}$$ (5) The gates $`CCNOT_{R,SQ}`$ and $`CCNOT_{S,QR}`$ can be written by analogy. ## 3 Virtual spin formalism In the realized NMR quantum gates one qubit NOT operation is implemented as a spin rotation under influence of a resonant pulse. Whereas conditional quantum dynamics, which is necessary for two qubit CNOT gate, is implemented using system with two body spin-spin interaction . It will be shown later in this paper that conditional quantum dynamics can be realized on one spin-7/2 without using spin-spin interactions. Another specific feature of such an approach is that all three types of gates NOT, CNOT and CCNOT can be implemented using just one RF pulse. These advantages are possible due to special information coding onto spin-7/2 states - this coding was introduced in under the name virtual spin formalism. The main idea of the virtual spin formalism can be expressed as follows. As a basis for the spin-7/2 Hilbert space $`\mathrm{\Gamma }_{7/2}`$ the eigen functions $`\|\chi _m`$ of $`I_z`$ operator $`(m=\pm \frac{1}{2},\pm \frac{3}{2},\pm \frac{5}{2},\pm \frac{7}{2})`$ or eigen functions $`\|\psi _m`$ of total spin energy operator (defined below) can be chosen. Let us instead $$m=\frac{7}{2},\frac{5}{2},\frac{3}{2},\frac{1}{2},+\frac{1}{2},+\frac{3}{2},+\frac{5}{2},+\frac{7}{2}$$ use notation $`M=0,1,2,3,4,5,6,7`$ and therefore instead $`\|\psi _m`$ to use $`\|\psi _M`$. Then formally $`\|\psi _M`$ can be mapped to $`\|M=\|m_Q,m_R,m_S`$ which is a state of virtual spins-1/2 $`Q`$, $`R`$, $`S`$. Then in order to implement all above mentioned gates it is necessary to find such external influence on real spin-7/2, which would have in the basis $`\|\psi _M`$ the evolution propagator matrices of the same operational forms as the above written gates in the basis $`\|M`$. ## 4 Spin-7/2 system and physical realization of gates Let us consider the system of nuclear spin $`I=7/2`$ placed in constant magnetic field $`H_0`$ and axially symmetrical electric crystal field: $$\begin{array}{cc}=_z+_Q,\hfill & \\ _z=\mathrm{}\omega _0I_z,\hfill & _Q=\frac{1}{2}\mathrm{}\omega _Q\underset{a=0,\pm 1,\pm 2}{}Q_\alpha q_\alpha ,\hfill \\ Q_0=I_z^2\frac{1}{3}I(I+1),\hfill & q_0=3\mathrm{cos}^2\theta 1,\hfill \\ Q_{\pm 1}=I_zI_{\pm 1}+I_{\pm 1}I_z,\hfill & q_{\pm 1}=\mathrm{sin}\theta \mathrm{cos}\theta e^{\pm i\varphi },\hfill \\ Q_{\pm 2}=I_{\pm 1}^2,\hfill & q_{\pm 2}=\frac{1}{2}\mathrm{sin}2\theta e^{\pm i2\varphi },\hfill \\ \omega _Q=\frac{3eqQ}{\left[2I\left(2I1\right)\mathrm{}\right]},\hfill & I_{\pm 1}=I_x\pm iI_y,\hfill \end{array}$$ (6) where $`e`$ \- the electron charge, $`Q`$ \- nuclear quadrupole moment, $`I_\beta (\beta =x,y,z)`$ \- spin components, $`2eq`$ \- the electric field gradient value, $`\theta `$ and $`\varphi `$ \- the polar angles, which define its axes orientation in laboratory system frame. Let us consider a case when $`\omega _Q<<\omega _0`$ so quadrupole interaction influence can be calculated using perturbation theory. It is supposed that quadrupole interaction much greater than the spin resonance line width, so the spectrum consists of seven well separated resonance lines. The first approximation gives the following energy levels and eigen functions: $$\begin{array}{c}E_m\mathrm{}\epsilon _m=\mathrm{}\omega _0m+\mathrm{}\omega _Qq_0(m^2\frac{21}{4}),\hfill \\ \|\psi _m=\chi _m+\underset{mk}{}\frac{\chi _k\left\|_Q\right\|\chi _m}{\mathrm{}\omega _0\left(km\right)}\chi _k,\hfill \end{array}$$ (7) here the normalization factor of $`\|\psi _m`$ is omitted. For simplicity the projective operators $`P_{mn}`$ notation will be used. They are matrices $`8\times 8`$ with all elements $`p_{kl}`$ equal zero except one $`p_{mn}=1`$. They have very simple relations: $$P_{kl}P_{mn}=\delta _{lm}P_{kn},P_{mn}=P_{nm}^+,P_{mn}\|\psi _k=\delta _{nk}\|\psi _m.$$ (8) An RF pulse, which is resonant for energy levels $`E_m`$ and $`E_n(E_m>E_n)`$, results in evolution operator $$\begin{array}{c}V_X(\varphi _{mn},f)=\widehat{1}+(P_{nn}+P_{mm})(\mathrm{cos}\frac{1}{2}\varphi _{mn}1)+\hfill \\ i(P_{mn}e^{if}+P_{nm}e^{if})\mathrm{sin}\frac{1}{2}\varphi _{mn},\hfill \\ \varphi _{mn}=2(tt_0)\gamma H_{rf}\|n\left\|I_x\right\|m\|,\widehat{1}=_mP_{mm},\hfill \end{array}$$ (9) where oscillating magnetic field is parallel to $`X`$ axis, and $`H_{rf}`$, $`f`$ and $`\mathrm{\Omega }(=\mathrm{\Omega }_{mn}(E_mE_n)/\mathrm{})`$ \- its amplitude, phase and frequency, and $`\widehat{1}`$ \- unit operator in space $`\mathrm{\Gamma }_{7/2}`$. A case when the field is parallel to $`Y`$ axis results in replacing $`f`$ by $`f+\frac{1}{2}\pi `$ in (9). Let us consider the realization of logic gates on the spin-7/2 states in order of increasing their complexity. The CCNOT gate requires one single frequency pulse. For example, the gate $`CCNOT_{Q,RS}`$ is realized using pulse with frequency $`\mathrm{\Omega }_{67}`$ and with rotation angle $`\pi `$. According to (9) the evolution operator at such conditions has the following form $$V_X(\pi _{67},0)=\widehat{1}(P_{77}+P_{66})+i(P_{67}+P_{76}).$$ (10) Taking into account the above mentioned isomorphism one can see that the following equality takes place: $$P_{67}+P_{76}=\|67\|+\|76\|=\|110111\|+\|111110\|.$$ (11) It means that the matrix of evolution operator $`V_X(\pi _{67},0)`$ is equal to matrix of the gate $`CCNOT_{Q,RS}`$ (Eq. (5)): $$V_X(\pi _{67},0)=CCNOT_{Q,RS}$$ (12) up to the phase factor $`i`$ for non diagonal elements. For other two gates one has $$\begin{array}{c}V_X(\pi _{75},0)=CCNOT_{Q,SR},\hfill \\ V_X(\pi _{73},0)=CCNOT_{R,SQ}.\hfill \end{array}$$ (13) It is necessary to point out that in comparison with transition $`\mathrm{\Omega }_{67}`$, the transitions $`\mathrm{\Omega }_{57}`$ and $`\mathrm{\Omega }_{37}`$ between the states $`\chi _M`$ are forbidden and become non zero in the first order of parameter $`\omega _Q/\omega _0`$. To get the $`\pi `$ rotation in these cases longer pulses or stronger RF field are necessary. However, numerical calculations show, that when $`\omega _Q`$ and $`\omega _0`$ are of the same order of magnitude, the expressions for rotation angles also have the matrix elements of the same orders. The CNOT gate requires one double frequency pulse, the evolution operator of which is a product of two operators: $$\begin{array}{c}V_X(\pi _{23},0)V_X(\pi _{67},0)=CNOT_{RS},\hfill \\ V_X(\pi _{13},0)V_X(\pi _{57},0)=CNOT_{SR},\hfill \\ V_X(\pi _{45},0)V_X(\pi _{67},0)=CNOT_{QS},\hfill \\ V_X(\pi _{15},0)V_X(\pi _{37},0)=CNOT_{SQ},\hfill \\ V_X(\pi _{46},0)V_X(\pi _{57},0)=CNOT_{QR},\hfill \\ V_X(\pi _{26},0)V_X(\pi _{37},0)=CNOT_{RQ}.\hfill \end{array}$$ (14) The NOT gate requires one four-frequency pulse, the evolution operator of which is a product of four operators: $$\begin{array}{c}V_X(\pi _{04},0)V_X(\pi _{15},0)V_X(\pi _{26},0)V_X(\pi _{37},0)=NOT_Q,\hfill \\ V_X(\pi _{02},0)V_X(\pi _{13},0)V_X(\pi _{46},0)V_X(\pi _{57},0)=NOT_R,\hfill \\ V_X(\pi _{01},0)V_X(\pi _{23},0)V_X(\pi _{45},0)V_X(\pi _{67},0)=NOT_S.\hfill \end{array}$$ (15) The physical realizations of gates, expressed in Eq. (12)-(15), in comparison with adopted in quantum information notation have additional phase factor $`i`$ for non diagonal elements. It should be taken into account later, when these gates will be used for constructing complex algorithms. Above for simplicity specific values of parameters $`\varphi `$ and $`f`$ have been used in evolution operators. The expressions (12)-(15) for CCNOT, CNOT, NOT can be easily generalized to get expressions for CCUT, CUT, UT. For example, if one uses Eq. (9) with arbitrary parameters $`\varphi `$ and $`f`$, then (12) gives the expression for CCUT. The corresponding calculations are straightforward but rather complicated, and would hide the main idea of the paper.
warning/0002/quant-ph0002086.html	ar5iv	text	# Pressure dependence of the Mg 3⁢𝑠⁢4⁢𝑠³⁢𝑆₁→3⁢𝑠⁢3⁢𝑝³⁢𝑃_{0,1,2} transition in superfluid 4He ## 1 Introduction Superfluid helium is a quantum substance with unique features, like the phenomenon of superfluidity or the unusual dispersion curve henshaw . Despite a successful history and expanded research in this field important properties of this quantum liquid remain still unexplained. Different experimental methods have been employed so far to study superfluid helium. In general, they can be divided into two groups of conceptually distinguishable approaches. Firstly, the superfluid itself is under investigation, which means parameters like its density, its friction or its phase diagram are measured. Secondly, the interaction of probe particles with the quantum fluid can be studied, e.g. the dispersion curve has been measured with neutron scattering. This group comprises experiments, where the experimental signal is derived from internal degrees of freedom of the microscopic probes. Foreign atoms and ions can be implanted and the changes in their spectra reveal information about the interactions of the probes with the helium environment toennies ; weis ; berndrev . In the experiment reported here magnesium atoms are introduced into the bulk superfluid and electronic transitions within them are observed. Foreign atoms or ions generally perturb the helium environment. Depending on the interaction between the probe particles and the superfluid helium distinctly different defect structures are formed. If the density around the foreign particle is lowered compared to the unperturbed helium bulk a void with the foreign atom in its center forms; such structures are known as bubbles. They are determined by the interplay of repulsive and attractive interactions, the Pauli repulsion between the electrons of the probe and the helium atoms surrounding it, as well as the volume, respectively the surface energy of the bubble. In contrast, there may be a strongly increased density, even larger than the solidification density. These objects have been named snowballs. They originate mainly from the very strong attractive polarization forces between the foreign particle and the surrounding liquid and are typically observed for positively charged particles, particularly for most of the positive ions due to strong monopole induced dipole interactions atkins ; toennies ; weis ; berndrev . From spectroscopy measurements mgpap it is known that magnesium atoms form bubble like structures under saturated vapour pressure. As a surprising feature magnesium atoms show in liquid helium an unusual three times longer lifetime for the $`3s3p^3P_13s^2{}_{}{}^{1}S_{0}^{}`$ intercombination transition compared to this transition in vacuum mgpap . For other systems such a behaviour has not been as pronounced as in this case. Therefore atomic magnesium has been chosen to study the influence of an increased helium pressure on a bubble-like structure in order to investigate whether this object is stable at higher helium pressures and may even undergo observable structure changes. Due to the interaction of the magnesium atoms with the surrounding superfluid its electronic states are perturbed and the emission as well as the absorption lines of corresponding electronic transitions are shifted with respect to their vacuum values. Further they are broadened and have asymmetric shapes berndrev . The wavelength of the electronic transitions and the mean bubble size can be predicted in the framework of a straightforward theoretical approach, the standard bubble model. This is based on macroscopic quantities such as surface and volume energies bubb1 . It has been successfully applied to singlet states so far exclusively berndrev . Here it is employed to describe triplet states as well. ## 2 Experimental set-up A copper pressure cell (inner volume $`=600`$cm<sup>3</sup>) is mounted inside a helium bath cryostat (see figure 2). Its temperature is maintained between $`1.2`$ and $`1.4`$K. The cell is connected with a helium gas reservoir via a capillary system (inner diameter $`=1.5`$mm) to allow filling by condensation of helium gas. The liquid pressure can be adjusted by applying a corresponding helium pressure from the gas reservoir. Optical access to the cell is possible through three quartz windows (diameter $`=39`$mm) which are sealed by indium gaskets up to $`40`$bar helium pressure at $`1.2`$K. In the experiment the sample material under investigation has a typical size of $`5`$x$`5`$x$`5`$mm. Ions are produced by laser ablation from the surface of the sample with a focused Nd:YAG laser (focal diameter $`4.4\mu `$m). The laser energy is $`8`$mJ per pulse with a pulse width of $`68`$ns at wavelength $`1064`$nm inapap . The ions are drawn by an electric field towards the bottom of the experimental chamber, where most of them recombine with electrons from a field emission tip. The tip voltage was varied between $`0.9`$ and $`2.8`$kV and the probe voltage between $`0.6`$ and $`1.2`$kV. These voltages were adjusted for each pressure to maximize the signal to noise ratio. These parameters correspond to electric fields between $`0.4`$ and $`1.0`$ $`\frac{kV}{cm}`$ for a drift length of $`42`$mm. The light emitted from the electron cascade after recombination is imaged onto the entrance of a grating monochromator (Czerny-Turner type) with a wavelength resolution of $`0.025`$nm. A photomultiplier tube (EMI S 20 extended) serves as detector, the signal of which is digitized and recorded time resolved in $`400`$ bins of a width of $`1.0`$ms. The recombination method as well as the implantation and production of ions directly in the liquid based on the use of laser ablation are both well established techniques berndrev . In this experiment they were combined for the first time. Experimental data were taken at pressures in the full accessible pressure range of the experimental method up to $`24`$bars, where close to the solidification point the ion mobility drops dramatically with increasing pressure. A typical spectrum of the $`3s4s^3S_13s3p^3P_{0,1,2}`$ transition at a helium pressure of $`1.5`$bar is displayed in figure 3. The mean wavelength of the three emission lines can be obtained by a fit of three convoluted Gaussian line shapes. ## 3 Calculation of the emission wavelength with the standard bubble model The bubble model allows a prediction of the bubble size as well as of the energy shift of electronic transitions compared with the free atomic case. The total energy of the defect $`E_{tot}`$ is the sum of two terms, the electronic contribution of the free atom $`E_{free}`$ and the so called defect energy $`E_{defect}`$ bubb1 : $$E_{tot}=E_{defect}+E_{free}=E_{bubble}+E_{int}+E_{free}$$ (1) The defect part includes the bubble energy $`E_{bubble}`$ which is needed to form the void and the pairwise interaction $`E_{int}`$ between the defect atom and surrounding helium atoms. The bubble energy consists of macroscopic terms like volume $`E_{vol}`$, surface $`E_{surf}`$ and volume kinetic $`E_{vk}`$ energies where the later is due to the helium density gradient at the bubble surface bubb1 ; bubb2 : $`E_{bubble}`$ $`=`$ $`E_{vol}+E_{surf}+E_{vk}`$ $`=`$ $`{\displaystyle \frac{4\pi }{3}}pR_B^3\left(R_0\right)+4\pi \sigma R_B^2\left(R_0\right)+`$ $`+`$ $`{\displaystyle \frac{\mathrm{}^2}{8m_{He}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{(\rho (r,R_0,\alpha ))^2}{\rho (r,R_0,\alpha )}}d^3r`$ (3) with the helium pressure $`p`$, the equilibrium bubble radius $`R_B`$, the radius $`R_0`$ where the liquid density approaches zero, the width of the transition region from the bubble to the helium environment $`1/\alpha `$, the surface density $`\sigma `$ and the density $`\rho (r,R_0,\alpha )`$. The density follows an assumed parametrization bubb2 : $`\rho =\{\begin{array}{cc}0\hfill & r<R_0\hfill \\ \rho _0\left[1\left[1+\alpha \left(rR_0\right)\right]e^{\alpha \left(rR_0\right)}\right]\hfill & rR_0\hfill \end{array}`$ (6) with the constant helium density $`\rho _0=0.146\frac{g}{cm^3}`$. This ansatz assumes that helium is incompressible as $`\rho (r,R_0,\alpha )`$ can’t be larger than $`\rho _0`$. The bubble model has been successfully applied to describe experiments at elevated helium pressures, e.g. for electron bubbles the pressure dependence of electronic transitions can be very well calculated ebubb . Further, there is a less than $`20`$% change in $`\rho _0`$ Donn over the whole pressure range covered in this experiment and the associated relative difference in the calculated pressure shift, which arises from the last term in eq.(3), is below $`210^3`$. Therefore we find this assumption motivated in our case. The defect energy is obtained by adding the interaction energy $`E_{int}`$ of the states involved and the bubble energy. Multi particle interactions are neglected in this approach and only pairwise magnesium- helium interactions are taken into account beau : $`E_{int}\left(S\right)`$ $`=`$ $`4\pi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}V_S\left(r\right)\rho (r,R_0,\alpha )r^2𝑑r`$ (7) $`E_{int}\left(P\right)`$ $`=`$ $`4\pi {\displaystyle \underset{0}{\overset{\pi }{}}}sin\theta d\theta {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}[\left(cos\theta \right)^2V_P^\sigma \left(r\right)+`$ (8) $`+`$ $`\left(sin\theta \right)^2V_P^\pi \left(r\right)]\rho (r,R_0,\alpha )r^2,dr`$ berndrev with the interatomic pair potentials $`V_S`$, $`V_P^\sigma `$, $`V_P^\pi `$, where S stands for s-states and P for p-states, which are in the case of magnesium atoms in triplet p-states only known without fine structure splitting pair . The fine structure splitting arising from spin orbit interactions is assumed not to depend on the externally applied helium pressure, therefore a prediction for all three emission lines can be made. As the energy of the free atom is only an additive contribution to the total energy, it can be neglected for the calculation of the radial dependence of the defect energy, but has to be added for the calculation of the wavelength of the electronic transitions. An example of the calculated defect energies of the two interesting states $`3s4s^3S_1`$ and $`3s3p^3P_1`$ for two different helium pressures ($`2.9`$, $`25`$bar) is shown in figure 4. The radius at the minimum of the defect energy is the mean equilibrium radius of the defect structure in the specific state. It decreases with increasing pressure (see figure 5) for the $`3s4s^3S_1`$ state from $`8.34`$Å at $`2.9`$mbar to $`7.68`$Å at $`25`$bar and for the $`3s3p^3P_1`$ state from $`4.85`$Å to $`4.56`$Å . This decrease in the equilibrium radius with increasing pressure is qualitatively similar to the behaviour of an electron bubble at an enhanced helium pressure ebubb . Additionally the model predicts the width of the transition region from the bubble to the helium environment to be $`0.45`$Å . The sum of the free energy and the difference of the two defect state energies yields a prediction of the pressure dependent emission wavelength of the transition $`3s4s^3S_13s3p^3P_1`$: $$\lambda (p)=(516.48\pm 0.01)nm(0.08\pm 0.01)\frac{nm}{bar}p[bar]$$ (9) The wavelength of the other two emission lines is obtained by adding the respective fine structure splitting ( $`+1.1`$nm for $`{}_{}{}^{3}P_{0}^{}`$, $`0.53`$nm for $`{}_{}{}^{3}P_{2}^{}`$) to the zero pressure wavelength of $`516.48`$nm. ## 4 Experimental results Typical measured emission spectra for helium pressures $`3`$mbar, $`1.5`$, $`8`$ and $`22`$bar are shown in figure 6. The values below $`1`$bar were measured with another experimental cell inapap as the pressure cell allows measurements only at helium pressures above $`1`$bar. The spectra shift with increasing pressure to smaller wavelength, in accordance with the bubble model. The central emission wavelength of the three transitions is given in figure 7, 8 and 9 as a function of the applied helium pressure. The error bars result from the line shape fits. The uncertainty of the wavelength calibration of the monochromator is 0.1 nm common to all points. The dotted line is the calculated wavelength predicted by the standard bubble model (see chapter 3). The pressure dependence of the three emission lines is: * $`3s4s^3S_13s3p^3P_0`$: $`\lambda =(517.11\pm 0.04)nm(0.09\pm 0.01)\frac{nm}{bar}p[bar]`$ * $`3s4s^3S_13s3p^3P_1`$: $`\lambda =(517.51\pm 0.06)nm(0.06\pm 0.01)\frac{nm}{bar}p[bar]`$ * $`3s4s^3S_13s3p^3P_2`$: $`\lambda =(518.52\pm 0.04)nm(0.06\pm 0.01)\frac{nm}{bar}p[bar]`$ The deviation in wavelength between the theoretical and the experimental curves is due to the precision of the pair potentials and is rather small compared to other calculations mbook , e.g. for barium atoms in superfluid helium the deviation is about $`14`$nm mdiss . The quality of the agreement of the calculated and measured pressure shifts for all three lines can be tested with a statistical hypothesis test, the students test. The deviation of the three values is compatible with statistical fluctuations. Therefore a mean pressure line shift of ($`0.07\pm 0.01`$nm$`/`$bar) can be derived. This very good consistency between the experimental and the theoretical values allows the conclusion that the magnesium atom seem to maintain a bubble like structure under increased helium pressures. The pressure shift is monotonous. ## 5 Discussion As a consequence of the higher pressure the bubble like defect shrinks, i.e. the equilibrium radius decreases. The repulsive part of the pair potential energies due to Pauli forces rises in the upper P state already at larger radii than for the lower S state which implies a smaller wavelength for emitted radiation. Up to now only few pressure dependent measurements of electronic transitions of foreign particles implanted into superfluid helium exist (see table 1). A quantitative comparison between the published line shifts and the results presented in this paper is not possible for the line shift themselves, because different types of transitions have been investigated. Since the foreign atom-helium interactions potentials are not comparable with each other, the different shifts for the various elements are not surprising. Interesting is a comparison concerning the relative pressure shift in wavelength which is much larger for the electron bubble than for any other structure. This reflects the fact that the electron bubble is much more compressible than the other bubbles. The similarity of the relative line shifts, i.e. the change of wavelength with pressure relative to the transition wavelength at saturated vapor pressure, for Mg, Rb, Ba, Tm and He<sub>2</sub> may be taken as indication that in all these cases bubbles are formed with similar size and compressibility. The within statistics linear behaviour of the pressure shifts suggests smooth and continuous change in the size and structure of the defect caused by all these systems. ###### Acknowledgements. This work was supported in part by the Deutsche Forschungsgemeinschaft (DFG). We would like to express our thanks to B. Tabbert and M. Foerste for their input at an early stage and their constant interest and suggestions. A. K. would like to acknowledge an Alexander von Humboldt postdoctoral fellowship.
warning/0002/math0002037.html	ar5iv	text	# Wave invariants for non-degenerate closed geodesics ## 0. Introduction This paper is a continuation of \[Z.1\]. There, we gave an effective method for putting the Laplacian $`\mathrm{\Delta }`$ of a Riemannian manifold $`(M,g)`$ into a quantum Birkhoff normal form around a non-degenerate elliptic closed geodesic $`\gamma `$, and applied it to the calculation and characterization of the wave invariants $`a_{\gamma k}`$ at $`\gamma .`$ The wave invariants, we recall, are the coefficients in the singularity expansion \| $`TrU(t)=e_o(t)+_\gamma e_\gamma (t)`$ \| \| --- \| \| $`e_\gamma (t)a_{\gamma 1}(tL_\gamma +i0)^1+_{k=0}^{\mathrm{}}a_{\gamma k}(tL_\gamma +i0)^klog(tL_\gamma +i0)`$ \| of the trace of the wave group $`U(t)=e^{it\sqrt{\mathrm{\Delta }}}`$ at lengths $`t=L_\gamma `$ of closed geodesics. The first purpose of this article is to extend the methods and results of \[Z.1\] to general non-degenerate closed geodesics, i.e. to $`\gamma `$ whose Poincare map $`P_\gamma `$ is any symplectic sum of non-degenerate elliptic, hyperbolic, or loxodromic parts. Our second purpose is to generalize to the full non-degenerate case the inverse result of Guillemin that the quantum normal form coefficients at non-degenerate elliptic closed geodesics are spectral invariants \[G.1,2\]. It will follow that, for metrics with simple length spectra, the classical Birkhoff normal form of the metric around any non-degenerate closed geodesic is a spectral invariant of the Laplacian. Let us state the results precisely. The first is that the wave invariants $`a_{\gamma k}`$ in the general non-degenerate case are essentially analytic continuations of the expressions obtained in the elliptic case. They may be written in the form $`(0.1)`$ $$a_{\gamma k}=_{k,1}(D)Ch(x)\|_{x=P_\gamma }$$ where $$Ch(x)=\frac{i^\sigma }{\sqrt{\|det(Ix)\|}}$$ is the character of the metaplectic representation (with $`\sigma `$ a certain Maslov index) and where $`_{k,1}(D)`$ is an invariant partial differential operator on the metaplectic group $`Mp(2n,)`$ which is canonically fashioned from the germ of the metric $`g`$ at $`\gamma `$. The exact expression for $`_{k,1}(D)`$ will be given in §5 and leads to the following characterization of the wave invariants (cf. \[Z.1, Theorem A\]): Theorem I Let $`\gamma `$ be a non-degenerate closed geodesic. Then $`a_{\gamma k}=_\gamma I_{\gamma ;k}(s;g)𝑑s`$ where: (i) $`I_{\gamma ;k}(s,g)`$ is a homogeneous Fermi-Jacobi-Floquet polynomial of weight -k-1 in the data $`\{y_{ij},\dot{y}_{ij},D_{s,y}^\beta g\}`$ with $`\|\beta \|2k+4`$ ; (ii) The degree of $`I_{\gamma ;k}`$ in the Jacobi field components is at most 6k+6; (iii) At most 2k+1 indefinite integrations over $`\gamma `$ occur in $`I_{\gamma kr}`$; (iv) The degree of $`I_{\gamma ;k}`$ in the Floquet invariants $`\beta _j`$ is at most k+2. The relevant terminology and notation will be recollected in §1. From (0.1) and from the formulae for $`_{k,1}(D)`$ and $`Ch`$, we give a rather simple proof (§6) of the following inverse result (strictly speaking, proven only for non-degenerate elliptic closed geodesics in \[G.2, Theorem 1.4\]): Theorem II Let $`\gamma `$ be a non-degenerate closed geodesic. Then the entire quantum Birkhoff normal form around $`\gamma `$ is a spectral invariant; in particular the classical Birhoff normal form is a spectral invariant. We thus have: Corollary II.1 Suppose $`(M,g)`$ is a compact Riemannian manifold with simple length spectrum $`Lsp(M,g)`$ and with all closed geodesics non-degenerate. Then from $`Spec(M,g)`$ one can recover the quantum (and hence classical) Birkhoff normal forms around all closed geodesics. The hypotheses of the corollary are of course satisfied by generic Riemannian metrics (cf. \[Kl, Lemma 4.4.3\]). The corollary therefore answers affirmatively the third question in \[Z.2, p.692\], which asks whether isospectral manifolds in this class of metrics are locally Fourier isospectral near corresponding closed geodesics. Let us now briefly discuss the main ideas in the proofs, and in particular the novel aspects caused by the hyperbolic and loxodromic parts of $`P_\gamma `$. As in \[G.1,2\]\[Z.1\], the wave invariants at a closed geodesic $`\gamma `$ will be expressed as non-commutative residues of the wave group and its time derivatives at $`t=L_\gamma .`$ For simplicity we will often abbreviate $`L_\gamma `$ by $`L`$. Then we have: $$a_{\gamma k}=resD_t^ke^{it\sqrt{\mathrm{\Delta }}}\|_{t=L}:=Res_{s=0}TrD_t^ke^{it\sqrt{\mathrm{\Delta }}}\sqrt{\mathrm{\Delta }}^s\|_{t=L}.$$ Since $`res`$ is invariant under conjugation by (microlocal) unitary Fourier integral operators, the $`a_{\gamma k}`$’s may be calculated by putting the wave group into a microlocal (quantum Birkhoff) normal form around $`\gamma `$ and by by determining the residues of the resulting wave group of the normal form. The primary step in the analysis of the wave invariants $`a_{\gamma k}`$ is therefore to put $`\mathrm{\Delta }`$ into this microlocal normal form around $`\gamma .`$ In the case of non-degenerate elliptic closed geodesics, we recall, the normal form was a polyhomogeneous function in the (microlocally elliptic) element $`D_s`$ with coefficients in the transverse (elliptic) harmonic oscillators $`\widehat{I}_j^e=\frac{1}{2}(D_{y_j}^2+y_j^2)`$ \[Z.1,Theorem B\]. In the general non-degenerate case, the normal form will involve a greater variety of quadratic normal forms or ‘action operators:’ in addition to the elliptic action operator $`\widehat{I}_j^e`$ there can also occur the real hyperbolic action operators $`\widehat{I}_j^h`$ and complex hyperbolic (or loxodromic) action operators $`\widehat{I}_j^{ch,Re},\widehat{I}_j^{ch,Im}.`$ These hyperbolic actions cause several complications to the arguments in the elliptic case: First, they have continuous spectra, and so the construction of the intertwining operator to the normal form has to be modified in several ways (§3,4). Second, the wave group of the normal form has continuous spectrum and this alters the calculation of its residues (§5). Third, the presence of real parts in the Floquet exponents of $`P_\gamma `$ complicates the process of determining the normal form coefficients from the wave invariants (§6). To get acquainted with these action operators and the normal form algorithm, let us consider the very first step of “linearizing” $`\sqrt{\mathrm{\Delta }}`$ around $`\gamma `$ and of putting the ‘linearization’ $$=D_s\frac{1}{2}(\underset{j=1}{\overset{n}{}}D_{y_j}^2+\underset{ij=1}{\overset{n}{}}K_{ij}y_iy_j)$$ into quantum quadratic normal form. Here, $`n=dimM1`$, the coordinates $`(s,y_j)`$ are the (re-scaled) Fermi normal coordinates around $`\gamma `$, $`D_s=\frac{}{is}`$ and $`K_{ij}`$ is the curvature operator $`g(R(_s,_{y_i})_s,_{y_j})`$. The linearization $``$ of $`\sqrt{\mathrm{\Delta }}`$ is a quadratic Hamiltonian and is the Weyl quantization of a quadratic classical Hamiltonian (see \[Ho III\] and \[Ho\] for background on Weyl quantizations and normal forms for quadratic Hamiltonians). Hence its symbol may be conjugated into normal form by an element of $`𝒲Sp(2n,).`$ The operator $``$ itself may be put into normal form by conjugating with the metaplectic operator $`\mu (𝒲)`$ quantizing $`𝒲`$. As will be seen in §1 , this linear symplectic map is the Wronskian matrix $`𝒲`$ whose columns consist of the Jacobi eigenfields of $`P_\gamma `$. The normal form of the linearized $`\sqrt{\mathrm{\Delta }}`$ is therefore given by $`=\mu (𝒲)^{}\mu (𝒲)^1.`$ In the elliptic case \[Z.1\] $`P_\gamma `$ was a direct sum of rotations, and the quantum normal form of $``$ had the form $$^e=D_s+\frac{1}{L}H_\alpha ,H_\alpha =\underset{j=1}{\overset{n}{}}\alpha _j\widehat{I}_j^e$$ where the spectrum $`\sigma (P_\gamma )=\{e^{\pm i\alpha _j}\}`$. In the general non-degenerate case the normal form will similarly depend on the spectral decomposition of $`P_\gamma `$. Recall that, since $`P_\gamma `$ is symplectic, its eigenvalues $`\rho _j`$ come in three types: (i) pairs $`\rho ,\overline{\rho }`$ of conjugate eigenvalues of modulus 1; (ii) pairs $`\rho ,\rho ^1`$ of inverse real eigenvalues; and (iii) 4-tuplets $`\rho ,\overline{\rho },\rho ^1\overline{\rho }^1`$ of complex eigenvalues. We will often write them in the forms: (i) $`e^{\pm i\alpha _j}`$, (ii)$`e^{\pm \lambda _j}`$, (iii) $`e^{\pm \mu _j\pm i\nu _j}`$ respectively (with $`\alpha _j,\lambda _j,\mu _j,\nu _j`$), although a pair of inverse real eigenvalues $`\{e^{\pm \lambda }\}`$ could be negative. Here, and throughout, we make the assumption that $`P_\gamma `$ is non-degenerate in the sense that $$\mathrm{\Pi }_{i=1}^{2n}\rho _i^{m_i}1,(\rho _i\sigma (P_\gamma ),(m_1,\mathrm{},m_{2n})𝐍^{2n}).$$ Each type of eigenvalue then determines a different type of quadratic action, both on the classical and quantum levels (cf. \[Ho, Theorem 3.1\],\[Ar\]): \| Eigenvalue type \| Classical Normal form \| Quantum normal form \| \| --- \| --- \| --- \| \| (i) Elliptic type \| $`I^e=\frac{1}{2}\alpha (\eta ^2+y^2)`$ \| $`\widehat{I}^e:=\frac{1}{2}\alpha (D_y^2+y^2)`$ \| \| $`\{e^{i\pm \alpha }\}`$ \| \| \| \| (ii) Real hyperbolic type \| $`I^h=2\lambda y\eta `$ \| $`\widehat{I}^h:=\lambda (yD_y+D_yy)`$ \| \| $`\{e^{\pm \lambda }\}`$ \| \| \| \| (iii) Complex hyperbolic \| $`I^{ch,Re}=2\mu (y_1\eta _1+y_2\eta _2)`$ \| $`\widehat{I}^{ch,Re}=\mu (y_1D_{y_1}+D_{y_1}y_1+y_2D_{y_2}+D_{y_2}y_2),`$ \| \| (or loxodromic type) \| $`I^{ch,Im}=\nu (y_1\eta _2y_2\eta _1)`$ \| $`\widehat{I}^{ch,Im}=\nu (y_1D_{y_2}y_2D_{y_1})`$ \| \| $`\{e^{\pm \mu +i\pm \nu }\}`$ \| \| \| In the case where the Poincare map $`P_\gamma `$ has p pairs of complex conjugate eigenvalues of moduls 1, q pairs of inverse real eigenvalues and c quadruplets of complex hyperbolic eigenvalues, the linearized $`\sqrt{\mathrm{\Delta }}`$ will have the form: $$=D_s+\frac{1}{L}[\underset{j=1}{\overset{p}{}}\alpha _j\widehat{I}_j^e+\underset{j=1}{\overset{q}{}}\lambda _j\widehat{I}_j^h+\underset{j=1}{\overset{c}{}}\mu _j\widehat{I}_j^{ch,Re}+\nu _j\widehat{I}_j^{ch,Im}].$$ The full quantum Birkhoff normal form is then given by the analogue of Theorem B of \[Z.1\]: Theorem B Assuming $`\gamma `$ non-degenerate, there exists a microlocally elliptic Fourier integral operator $`W`$ from the conic neighborhood of $`^+\gamma `$ in $`T^{}(N_\gamma )`$ to the corresponding cone in $`T_+^{}S^1`$ in $`T^{}(S^1\times ^n)`$ such that $$W\sqrt{\mathrm{\Delta }}W^1D_s+\frac{1}{L}[\underset{j=1}{\overset{p}{}}\alpha _j\widehat{I}_j^e+\underset{j=1}{\overset{q}{}}\lambda _j\widehat{I}_j^h+\underset{j=1}{\overset{c}{}}\mu _j\widehat{I}_j^{ch,Re}+\nu _j\widehat{I}_j^{ch,Im}]+$$ $$+\frac{p_1(\widehat{I}_1^e,\mathrm{},\widehat{I}_p^e,\widehat{I}_1^h,\mathrm{},\widehat{I}_q^h,\widehat{I}_1^{ch,Re},\widehat{I}_1^{ch,Im},\mathrm{},\widehat{I}_c^{ch,Re},\widehat{I}_c^{ch,Im})}{D_s}+$$ $$+\frac{p_{k+1}(\widehat{I}_1^e,\mathrm{},\widehat{I}_c^{ch,Im})}{D_s^k}+$$ where the numerators $`p_j(\widehat{I}_1^e,\mathrm{},\widehat{I}_p^e,\widehat{I}_1^h,\mathrm{},\widehat{I}_c^{ch,Im})`$ are polynomials of degree j+1 in the variables $`(\widehat{I}_1^e,\mathrm{},\widehat{I}_c^{ch,Im})`$ and where the kth remainder term lies in the space $`_{j=o}^{k+2}O_{2(k+2j)}\mathrm{\Psi }^{1j}`$ Here, $`O_n\mathrm{\Psi }^r`$ is the space of pseudodifferential operators of order r whose complete symbols vanish to order n at $`(y,\eta )=(0,0).`$ Thus, the remainder terms are ‘small’ in that they combine in some mixture a low pseudodifferential order or a high vanishing order along $`\gamma `$. Some remarks now on the contents and organization of this paper. Since it is a continuation of \[Z.1\], we have tried to avoid duplicating arguments and calculations which are essentially unchanged from the elliptic case. Many of the arguments which remain are still quite analogous to the elliptic case and inevitably produce a sense of deja-vu. Our excuse for drawing them out to their present length is that it is not apriori clear that the arguments of the elliptic case generalize so neatly to the hyperbolic and loxodromic cases. It may even be viewed as a virtue of the method of \[Z.1\] that it adapts so effortlessly to the general case. It should be noted here that Guillemin was aware that the arguments of the elliptic case should extend to the general non-degenerate case and stated his main result, Theorem 1.4 of \[G.2\], for the general case. However, we also note that the methods used here are extensions of the methods of \[Z.1\], which in many significant respects differ from the methods of \[G.2\]. The organization of this paper is as follows: In §1 we will review the symplectic and microlocal ingredients required to construct a ‘linear model’ for the Laplacian near a closed geodesic $`\gamma `$. In §2, we introduce the semi-classically scaled Laplacian and the linearized Laplacian and conjugate them to the model space. In §3, we conjugate the resulting semi-classical model Laplacian to a semi-classical normal form to infinite order. In §4 we show how this semi-classical normal form induces a bona-fide quantum Birkhoff normal form for the Laplacian near $`\gamma `$. In §5, we use the normal form to give the explicit formula (0.1) for the wave invariants. In §6, we show that the quantum normal form coefficients can be determined from the special values of (0.1) corresponding to $`\gamma `$ and its iterates. ## 1. Preliminaries This section begins with a resume of the symplectic linear algebra underlying the Jacobi equation, the linear Poincare map, and the symplectic classification of non-degenerate quadratic forms (§1.1). It then summarizes the quantum aspects of the linear theory, in particular the behaviour of the quantum action operators (§1.2) §1.1: Symplectic preliminaries §1.1a: Closed geodesics, linear Poincare maps and Jacobi fields Throughout this paper, $`\gamma `$ will denote a primitive closed geodesic of $`(M,g)`$; its iterates will be denoted by $`\gamma ^m`$. The space $`𝒥_\gamma ^{}`$ of (real) orthogonal Jacobi fields along $`\gamma `$ is then the real symplectic vector space, of dimension 2n, of solutions of the Jacobi equation $`Y^{\prime \prime }+R(T,Y)T=0`$ (with $`T`$ the unit tangent vector along $`\gamma .`$) The symplectic structure is given by the Wronskian $$\omega (X,Y)=g(X,\frac{D}{ds}Y)g(\frac{D}{ds}X,Y).$$ The linear Poincare map $`P_\gamma `$ is the (real) linear symplectic map on $`(𝒥_\gamma ^{},\omega )`$ defined by $`P_\gamma Y(t)=Y(t+L_\gamma ).`$ To diagonalize it, we also complexify it as a complex symplectic map $`P_\gamma ^{}`$ on the space $`𝒥_\gamma ^{}`$ of complex orthogonal Jacobi fields. Here, the symplectic form is extended to the complexified space as a complex bilinear form $`\omega ^{}`$. Since $`P_\gamma ^{}Sp(𝒥_\gamma ^{},\omega )`$, its spectrum $`\sigma (P_\gamma ^{})`$ is stable under inverse and complex conjugation: thus, if $`\rho \sigma (P_\gamma ^{})`$, then also $`\rho ^1,\overline{\rho },\overline{\rho }^1\sigma (P_\gamma ^{})`$. As mentioned above, we will assume that $`P_\gamma ^{}`$ is non-degenerate in the following strong sense: $`(1.1a.1)`$ $$\rho _1^{m_1}\mathrm{}\rho _n^{m_n}=1m_i=0(i,m_i).$$ In particular, the eigenvalues are simple and $`\pm 1\sigma (P_\gamma ^{}).`$ The eigenspace of $`P_\gamma ^{}`$ of eigenvalue $`\rho `$ will be denoted by $`𝒥_\gamma ^,(\rho )𝒥_\gamma ^{}.`$ We then have the symplectic orthogonal decomposition $`(1.1a.2)`$ $$𝒥_\gamma ^,=𝒥_{nc}^,𝒥_{co}^,$$ into the ‘non-compact’ symplectic subspace $`(1.1a.3.nc)`$ $$𝒥_{nc}^,=_{\rho :\|\rho \|1}𝒥_\gamma ^,(\rho )$$ where $`P_\gamma ^{}`$ does not belong to a compact subgroup of $`Sp`$ and the ‘compact’ symplectic subspace $`(1.1a.3.co)`$ $$𝒥_{co}^,=_{\rho :\|\rho \|=1}𝒥_\gamma ^,(\rho )$$ where $`P_\gamma ^{}`$ does belong to a compact subgroup of $`Sp.`$ This and the following decompositions are described in more detail in Klingenberg \[Kl\], but also somewhat differently since in \[Kl\] the symplectic form is extended to the complexification as a sesquilinear form rather than as a complex bilinear form. The non-compact subspace has the further symplectic orthogonal decompositon $`(1.1a.4)`$ $$𝒥_{nc}^,=\underset{\rho :\|\rho \|<1}{}[𝒥_\gamma ^,(\rho )𝒥_\gamma ^,(\rho ^1)]$$ into symplectic complex 2-planes. We may rewrite this decomposition in the form $`(1.1a.5)`$ $$𝒥_{nc}^,=𝒥_s^,𝒥_u^,$$ where $$𝒥_s^,=\underset{\rho :\|\rho \|<1}{}𝒥_\gamma ^,(\rho ),𝒥_u^,=\underset{\rho :\|\rho \|>1}{}𝒥_\gamma ^,$$ are the symplectically dual stable, resp. unstable Lagrangean subspaces. The compact subspace has the further symplectic decompositon $`(1.1a.6)`$ $$𝒥_{co}^,=\underset{\rho :\rho =e^{i\alpha },\alpha (0,\pi )}{}𝒥_\gamma ^,(\rho )𝒥_\gamma ^,(\overline{\rho }).$$ Any choice of one $`\rho `$ from a pair $`\{\rho ,\overline{\rho }\}`$ determines a splitting of $`𝒥_{co}^,`$ into a pair of dual Lagrangean subspaces. On the level of real symplectic spaces, we have the closely related $`P_\gamma `$-invariant symplectic decomposition $`(1.1a.7)`$ $$𝒥_\gamma ^{}=𝒥_s^r𝒥_u^r𝒥_{ce}^{,2p}$$ into the stable, unstable and center stable real subspaces of dimensions $`r,r,2p`$ respectively. By definition, $`(1.1a.8)`$ $$𝒥_s^r=\underset{\rho ,\|\rho \|<1}{}𝒥_\gamma ^,(\rho )\underset{\rho ,\|\rho \|<1}{}𝒥_\gamma ^,(\rho )$$ where $$P_\gamma \|_{𝒥_\gamma ^,(\rho )}=\rho (\rho )$$ respectively $$P_\gamma \|_{𝒥_\gamma ^,(\rho )}=e^\mu \left(\begin{array}{cc}cos\nu \hfill & sin\nu \hfill \\ sin\nu \hfill & cos\nu \hfill \end{array}\right)\rho =e^{\mu +i\nu },\mu ,\nu >0.$$ In the latter case, $`𝒥_\gamma ^,(\rho )`$ is the real symplectic 2-plane whose complexification equals $`𝒥_\gamma ^,(\rho )𝒥_\gamma ^,(\rho ^1)`$. Similarly for the case of the unstable subspace. In the center stable case, $`(1.1a.9)`$ $$𝒥_{\gamma ,ce}^{,2p}=\underset{j=1}{\overset{p}{}}𝒥_\gamma ^,(\alpha _j)$$ where $$P_\gamma \|_{𝒥_\gamma ^,(\alpha )}=\left(\begin{array}{cc}cos\alpha \hfill & sin\alpha \hfill \\ sin\alpha \hfill & cos\alpha \hfill \end{array}\right)$$ and with $`𝒥_\gamma ^,(\alpha )`$ the symplectic two plane whose complexification equals $`𝒥_\gamma ^,(\rho )𝒥_\gamma ^,(\overline{\rho })`$ with $`\rho =e^{i\alpha }.`$ We will put $`r=q+2c`$ where $`q=\mathrm{\#}\{\rho ,\|\rho \|<1\}`$ and where $`2c=\mathrm{\#}\{\rho ,\|\rho \|<1\}`$ and say that $`\gamma `$ has type $`(p,q,c)`$ if it has $`p`$ pairs of conjugate eigenvalues of modulus one $`\{e^{i\alpha },e^{i\alpha }\},`$ $`q`$ pairs of real inverse eigenvalues $`\{e^\lambda ,e^\lambda \}`$, and $`c`$ quadruples of non-real complex eigenvalues $`\{e^{\pm \mu +\pm i\nu }\}.`$ Here, we have assumed the real eigenvalues are positive for brevity of notation. Finally we may write: $`(1.1a.10)`$ $$𝒥_\gamma ^{}=𝒥_\gamma ^e𝒥_\gamma ^h𝒥_\gamma ^{ch}$$ where $`𝒥_\gamma ^e=𝒥_{\gamma ,ce}^{,2p}`$ is the elliptic (or real center stable) subspace, where $`(1.1a.11)`$ $$𝒥_\gamma ^h=\underset{\rho ,\|\rho \|<1}{}𝒥_\gamma ^,(\rho )𝒥_\gamma ^,(\rho ^1)$$ is the real hyperbolic subspace and where $`(1.1a12)`$ $$𝒥_\gamma ^{ch}=\underset{\rho ,\|\rho \|<1}{}𝒥_\gamma ^,(\rho )𝒥_\gamma ^,(\rho ^1)$$ is the complex hyperbolic (or loxodromic) subspace. §1.1b: Jacobi eigenvectors and Wronskian matrix As mentioned in the introduction, the intertwining operator to the quantum normal form will involve a certain Wronskian matrix of the Jacobi equation. Roughly speaking, it is the real symplectic matrix whose entries are given by the real and imaginary parts of the Jacobi eigenvectors and their time derivatives relative to a normal frame. More precisely, we fix a symplectic orthonormal basis of Jacobi eigenvectors as follows: \| complex subspace \| eigenvectors \| normalization \| \| --- \| --- \| --- \| \| elliptic plane \| $`P_\gamma Y_j^e=e^{i\alpha _j}Y_j^e,P_\gamma \overline{Y^e}_j=e^{i\alpha _j}\overline{Y^e}_j`$ \| $`\omega (Y_j^e,\overline{Y}_j^e)=1.`$ \| \| real hyp. plane \| $`P_\gamma Y_j^+=e^{\lambda _j}Y_j^+,P_\gamma Y_j^{}=e^{\lambda _j}Y_j^{}`$ \| $`\omega (Y_j^+,Y_j^{})=1.`$ \| \| cx.hyp.4-plane \| $`P_\gamma Y_j^{++}=e^{\mu +i\nu }Y_j^{++},P_\gamma Y_j^{}=e^{\mu i\nu }Y_j^{},`$ \| $`\omega (Y^{++},Y^{})=1`$ \| \| \| $`P_\gamma Y_j^+=e^{\mu i\nu }Y_j^+,P_\gamma Y_j^+=e^{\mu +i\nu }Y_j^+,`$ \| $`\omega (Y^+,Y^+)=1`$ \| The normalization makes sense since complex 2-planes spanned by eigenvectors corresponding to inverse eigenvalues are symplectic. We now fix a parallel normal frame $`e(s):=(e_1(s),\mathrm{},e_n(s))`$ along $`\gamma \|_{[O,L)}`$ and denote by $`Y,e_j`$ the Riemannian inner product of a vector $`Y`$ along $`\gamma `$ with the the jth normal vector. Corresponding to the splitting of $`𝒥_\gamma ^{}`$ into its elliptic, hyperbolic and and complex hyperbolic (real) subspaces, we then get a real symplectic 2n x 2n Wronskian matrix $`(1.1b.1)`$ $$𝒲=[𝒲^e\|𝒲^h\|𝒲^{ch}]$$ formed by the 2n x 2p elliptic Wronskian matrix $`(1.1b.2e)`$ $$𝒲^e(s):=\left(\begin{array}{cc}ReY_i^e,e_j\hfill & ImY_i^e,e_j\hfill \\ Re\dot{Y_i}^e,e_j\hfill & Im\dot{Y_i^e},e_j\hfill \end{array}\right)_{i=1,\mathrm{},p;j=1,\mathrm{},n}$$ the 2n x 2q real hyperbolic Wronskian matrix $`(1.1b.2h)`$ $$𝒲^h(s):=\left(\begin{array}{cc}Y_i^+,e_j\hfill & Y_i^{},e_j\hfill \\ \dot{Y_i}^+,e_j\hfill & \dot{Y_i}^{},e_j\hfill \end{array}\right)_{i=1,\mathrm{},q;j=1,\mathrm{},n}$$ and finally the complex hyperbolic 2n x 4c Wronskian matrix $`(1.1b.2ch)`$ $$𝒲^{ch}(s):=\left(\begin{array}{cccc}ReY_i^{++},e_j\hfill & ImY_i^{++},e_j\hfill & ReY_i^{},e_j\hfill & ImY_i^{},e_j\hfill \\ Re\dot{Y_i}^{++},e_j\hfill & Im\dot{Y_i}^{++},e_j\hfill & Re\dot{Y_i}^{},e_j\hfill & Im\dot{Y_i}^{},e_j\hfill \end{array}\right)_{i=1,\mathrm{},c;j=1,\mathrm{},n}.$$ To see that $`𝒲`$ is indeed symplectic, and to better understand its properties, we reconsider the Jacobi equation and Poincare map from the Riemannian viewpoint. Thus, we let $``$ denote the Riemannian connection, and recall that it determines a horizontal subbundle of $`T(S^{}M)`$ complementary to the vertical subbundle of the projection $`\pi :S^{}MM.`$ Together with the symplectic structure, we get a splitting $$T(S^{}M)=\overline{H}\overline{V}\overline{T}$$ where $`\overline{T}`$ is the real span of $`\dot{\gamma }`$, and $`\overline{H}\overline{V}`$ is the horizontal plus vertical decomposition of the kernel of the contact form $`\alpha =\xi dx`$ (or equivalently, of the symplectic orthogonal of $`T`$ and the cone axis). The subspaces $`\overline{H},\overline{V}`$ are symplectically paired Lagrangean subspaces of $`T(T^{}M)`$. Given a vector $`XN_{\gamma (t)}`$, we denote by $`X^h`$ the horizontal lift of $`X`$ to $`\overline{H}_{\gamma (t)}`$ and by $`X^v`$ the vertical lift to $`\overline{V}_{\gamma (t)}`$. The correspondence $$Y(t)(Y(t)^h,\dot{Y}(t)^v)$$ then defines an isomorphism between the spaces of Jacobi fields and geodesic flow invariant vector fields along $`(\gamma (t),\dot{\gamma }(t)`$ (cf \[Kl, Lemma 3.1.6\]). That is, $$dG_{(\gamma (0),\dot{\gamma }(0))}^s:(Y(0)^h,\dot{Y}(0)^v)(Y(s)^h,\dot{Y}(s)^v)$$ where $`Y(s)`$ is the Jacobi field with the given initial conditions. Moreover, since $`G^t`$ is a Hamiltonian flow $`dG^s`$ is a linear symplectic mapping from $`(\overline{H}\overline{V})_{\gamma (0),\dot{\gamma }(0)}`$ to $`(\overline{H}\overline{V})_{\gamma (s),\dot{\gamma }(s)}.`$ Relative to the basis $`\{e_j(s)\}`$, the Jacobi equation is equivalent to the linear system $`\frac{D}{ds}(Y,P)=JH(Y,P)`$, where: (i) $`P=\frac{DY}{ds},`$ (ii) $`J`$ is the standard complex structure on $`R^{2n},`$ and where (iii) $$H=\left(\begin{array}{cc}K\hfill & 0\hfill \\ 0\hfill & I\hfill \end{array}\right)$$ with $`K`$ the curvature matrix and with $`I`$ the identity matrix. Moreover the basis $`\{e_j(s)\}`$ induces a moving symplectic frame $`\{e_j^h(s),e_j^v(s)\}`$ of $`(\overline{H}\overline{V})_{\gamma (s),\dot{\gamma }(s)}.`$ The evolution operator for the linear system is then just $`dG^s`$ expressed as a matrix relative to the moving symplectic frame. Now consider the above basis $`\{ReY_j^e,ImY_j^e,Y_j^+,Y_j^{},ReY_j^{++},ImY_j^{++},ReY_j^{},ImY_j^{}\}`$ of $`𝒥_\gamma ^{}`$ formed by the eigenvectors of $`P_\gamma `$. By construction, it is a symplectic basis relative to the Wronskian form $`\omega `$. Hence the pairs $`(Y^h(s),P^v(s))`$ consisting of these eigenvectors and their time derivatives form a moving symplectic basis of $`\overline{H}\overline{V}.`$ Expressed in terms of the frame $`\{e_j^h(s),e_j^v(s)\}`$ we then get a symplectic basis of $`^{2n}`$ relative to the standard symplectic structure. The Wronskian matrix $`𝒲(s)`$ is just the matrix whose columns are formed by these basis elements, and it is therefore symplectic for each $`s`$. Consider now the monodromy aspect of $`𝒲(s)`$, i.e. its transformation law under time translation $`ss+L`$ thru one period. Let $`Y_i`$ denote one of the complex eigenvectors of $`P_\gamma `$. Then the matrix element $`Y_i(s),e_j(s)`$ (or with $`\dot{Y}_i`$ in place of $`Y_i`$) satisfies $`(1.1b.3)`$ $$Y_i(s+L),e_j(s+L)=\rho _i\underset{k=1}{\overset{n}{}}t_{jk}Y_i(s),e_k(s)$$ where $`T:=(t_{jk})`$ is the holonomy matrix, $$e_j(s+L)=\underset{k=1}{\overset{n}{}}t_{jk}e_k(s).$$ It follows that $`(1.1b.4)`$ $$𝒲(s+L)=𝒲(s)T^{}P_\gamma .$$ §1.1c: Symplectic equivalence of quadratic Hamiltonians Let $`(^{2n},\omega )`$ be the standard symplectic vector space, endowed with linear coordinates $`x=(q_1,\mathrm{},q_n,p_1,\mathrm{},p_n)`$ such that $`\omega =dq_idp_i.`$ A quadratic (real) Hamiltonian is by definition a quadratic form $$H(q,p)=\frac{1}{2}Ax,x=\frac{1}{2}\omega (JAx,x)$$ where $``$ is the Euclidean scalar product, where $`A`$ is a 2n x 2n real symmetric matrix and where $`J=\left(\begin{array}{cc}0\hfill & I\hfill \\ I\hfill & 0\hfill \end{array}\right).`$ Then $`JAsp(^{2n},\omega )`$, and hence its spectrum decomposes into purely imaginary pairs $`(i\alpha ,i\alpha )`$, into real pairs $`(\lambda ,\lambda )`$, and into complex quadruples $`(\pm \mu \pm i\nu ).`$ We will assume, as above, that the eigenvalues are simple and not equal to $`\pm 1.`$ Then $`H(q,p)`$ decomposes into sums of terms of the following types of quadratic Hamiltonians, or classical ‘actions’: the elliptic type $`(1.1c.1e)`$ $$I^e(q_1,p_1):=\frac{1}{2}\alpha (q_1^2+p_1^2),$$ the real hyperbolic type $`(1.1c.1h)`$ $$I^h(q_1,p_1):=\frac{1}{2}\lambda q_1p_1,$$ and the complex hyperbolic (or loxodromic) type $`(1.1c.1ch)`$ $$I^{ch}(q_1,p_1,q_2,p_2)=\frac{1}{2}\mu (q_1p_1+q_2p_2)+\frac{1}{2}\nu (q_1p_2q_2p_1).$$ Note that $$q_1p_1+q_2p_2=Re(q_1+iq_2)(p_1ip_2),q_2p_1q_1p_2=Im(q_1+iq_2)(p_1ip_2)$$ and $$\{q_1p_1+q_2p_2,q_1p_2q_2p_1\}=0$$ To unify these expressions, we observe that they all have the form $`\frac{1}{2}Resa^{}a`$ where $`s`$ and where $`a^{},a`$ denote symplectically dual complex linear coordinates. Indeed, in the elliptic case, $`a=q_1+ip_1,a^{}=q_1ip_1,s=\alpha ;`$ in the real hyperbolic case, $`a=q_1,a^{}=p_1,s=\lambda ;`$ and in the loxodromic case, $`s=\mu +i\nu ,a=(q_1+iq_2),a^{}=(p_1ip_2).`$ §1.2: Microlocal preliminaries As mentioned above, the wave invariants only involve the metric and Laplacian $`\mathrm{\Delta }`$ in a tubular neighborhood of $`\gamma `$. In fact, as discussed in \[G.1\]\[Z.1\] they only involve the microlocalization of $`\mathrm{\Delta }`$ to the conic neighborhood $`(\mathrm{1.2.1})`$ $$\|y\|ϵ,\|\eta \|<ϵ\sigma $$ of $`T^{}S_L^10`$ in $`T^{}(S^1\times ^n)`$. Here,$`(s,\sigma ,y,\eta )`$ denote the symplectic Fermi coordinates and $`\psi `$ denotes a smooth homogeneous cut-off function on $`T^{}(S_L^1\times ^n)0`$ which equals 1 in some conic neighborhood $`V`$ of $`T^{}S_L^10`$ and vanishes identifically off of some slightly larger conic neighborhood. As in the case of elliptic closed geodesics, to put $`\mathrm{\Delta }`$ into a microlocal (quantum Birkhoff) normal form around $`\gamma S^1`$ is first of all to conjugate it to a distinguished maximal abelian subalgebra $`𝒜_\gamma `$ of the algebra $`\mathrm{\Psi }^{}(S_L^1\times ^n)`$ of pseudo-differential operators on the model space $`S_L^1\times ^n`$. This distinguished subalgebra will depend on the type of the geodesic $`\gamma `$. Roughly speaking, it will consist of the tangential operator $`D_s:=\frac{1}{i}\frac{}{s}`$ together with an appropriate set of quantized quadratic normal forms or ‘action operators’ in the transversal directions. §1.2.1: The model algebras To specify this ‘appropriate set’ of action operators, we begin by recalling that the Schrodinger representation of the (complexified) Heisenberg algebra $`𝐡_n`$ on the transverse space $`L^2(^n)`$, is generated by the self-adjoint operators $`Y_j=`$ “multiplication by $`y_j`$” and by $`D_j=\frac{}{iy_j}`$. Equivalently it is generated by the creation/annihilation operators $`Y_j+iD_j,Y_jiD_j.`$ For our purposes, however, it will be more natural to choose a different set of generators depending on the $`(q,p,c)`$ type of the closed geodesic. Corresponding to the $`2p`$ dimensional elliptic symplectic subspace we will use as generators the above (elliptic) annihilation/creation operators $`(\mathrm{1.2.1.1}e)`$ $$Z_j:=Y_j+iD_{y_j}Z_j^{}=Y_jiD_{y_j}(j=1,\mathrm{},p)$$ which satisfy the commutation relations $$[Z_j,Z_k]=[Z_j^{},Z_k^{}]=0[Z_j,Z_k^{}]=2\delta _{ij}I.$$ We would like to use the real (resp. complex) hyperbolic analogues in the hyperbolic subspaces. To determine the analogues we note that $`Z_j`$, resp. $`Z_j^{}`$ are the Weyl quantizations of the symplectically dual (modulo a factor of 2) complex linear coordinates $`z_j:=y_j+i\eta _j,`$ resp. $`z_j^{}:=y_ji\eta _j`$. We use the ‘dagger’ notation rather than the adjoint notation $`Z_j^{}`$ to emphasize that the dual operators are symplectically dual; they are also adjoints of each other, but this property will not extend to the hyperbolic cases. Indeed, corresponding to the $`2q`$ dimensional real hyperbolic subspace, the natural generators are the hyperbolic annihilation/creation operators $`(\mathrm{1.2.1.1}h)`$ $$Y_j,D_{y_j}(j=p+1,\mathrm{},q)$$ which of course are also symplectically dual. And corresponding to the $`4c`$ dimensional complex hyperbolic subspace, we the natural generators are the complex hyperbolic annihilation/creation operators $`(\mathrm{1.2.1}ch)`$ $$W_j:=Y_j+iY_{c+j},W_j^{}:=D_{y_j}iD_{y_{c+j}},\overline{W_j}:=Y_jiY_{c+j},\overline{W_j}^{}:=D_{y_j}+iD_{y_{c+j}}(j=p+q+1,\mathrm{},c).$$ We note that they satisfy the commutation relations: $$[W_j,W_j^{}]=2,[\overline{W_j},\overline{W_j}^{}]=2$$ with all other brackets zero. We will not bother to renormalize the operators to be precisely dual. The enveloping algebra of the Heisenberg algebra is then generated by all the above annihilation/creation operators, $`(\mathrm{1.2.1.2})`$ $$:=<Z_1,\mathrm{}Z_p,Z_1^{},\mathrm{},Z_p^{},Y_1,\mathrm{},Y_q,D_{y_1},\mathrm{},D_{y_q},W_1,\overline{W}_1,\mathrm{},W_c,\overline{W}_c,W_1^{},\overline{W}_1^{},\mathrm{},W_c^{},\overline{W}_c^{}>$$ and is of course the algebra of partial differential operators on $`^n`$ with polynomial coefficients. We will denote by $`^n`$ the subspace of polynomials of degree n in the generators. The microlocalization of this algebra is the isotropic Weyl algebra $`𝒲^{}`$ of pseudo-differential operators on $`^n`$, in which the generators are assigned the order $`\frac{1}{2}`$, so that $$^n𝒲^{n/2}$$ $$[^m,^n]^{m+n2}.$$ The symplectic algebra $`\mathrm{𝐬𝐩}(n,)`$ is then represented in $`^2`$ by homogeneous quadratic polynomials in the generators, which have degree 1. In particular it contains the following elliptic, resp. hyperbolic, resp. complex hyperbolic ( loxodromic) ‘action’ operators: $`(\mathrm{1.2.1.3})`$ $$\widehat{I}_j^e:=Z_j^{}Z_j,\widehat{I}_j^h=\frac{1}{2}(Y_jD_{y_j}+D_{y_j}Y_j),\widehat{I}_j^{ch,Re}=\frac{1}{2}Re(W_j^{}W_j+(W_j^{}W_j)^{}),\widehat{I}_j^{ch,Im}=ImW_j^{}W_j.$$ The complex hyperbolic action operators can also be written in the form $`(\mathrm{1.2.1.4})`$ $$\widehat{I}_j^{ch,Re}=\frac{1}{2}(Y_jD_{y_j}+D_{y_j}Y_j+Y_{j+c}D_{y_{j+c}}+D_{y_{j+c}}Y_{j+c}),\widehat{I}_j^{ch,Im}=(Y_jD_{y_{j+c}}Y_{j+c}D_{y_j})$$ where the coordinates are indexed so that the $`dy_jdy_{j+c}d\eta _jd\eta _{j+c}`$-planes are the $`P_\gamma `$-invariant complex hyperbolic 4-planes. It is then natural to introduce polar coordinates $`r_j,\varphi _j`$ on the $`(y_j,y_{j+c})`$-plane so that the loxodromic actions operators simplify to $`(\mathrm{1.2.1.5})`$ $$I_j^{ch,Re}=\frac{1}{2}(r_jD_{r_j}+D_{r_j}r_j),I^{ch,Im}=D_\theta .$$ We now introduce the distinguished $`(p,q,c)`$ maximal (transverse) abelian subalgebra of $`𝒲`$, given by $`(\mathrm{1.2.1.6})`$ $$_{p,q,c}:=<I_1^e,\mathrm{},I_p^e,I_1^h,\mathrm{},I_q^h,I_1^{ch,Re},\mathrm{},I_c^{ch,Re},I_1^{ch,Im},\mathrm{},I_c^{ch,Im}>.$$ Together with the tangential operator we get the (p,q,c)- maximal abelian subalgebra given by $`(\mathrm{1.2.1.7})`$ $$𝒜_{p,q,c}:=<D_s,I_1^e,\mathrm{},I_p^e,I_1^h,\mathrm{},I_q^h,I_1^{ch,Re},\mathrm{},I_c^{ch,Re},I_1^{ch,Im},\mathrm{},I_c^{ch,Im}>.$$ §1.2.2: Model eigefunctions An orthonormal basis of $`L^2(S_L^1\times ^n)`$ of joint eigenfunctions of $`𝒜_{p,q,c}`$ is given as follows: corresponding to the (p,q,c)- type of $`P_\gamma `$ we can write $$L^2(S_L^1\times ^n)=L^2(S_L^1)L^2(^p)L^2(^q)L^2(^{2c})$$ and construct the eigenfunctions as (tensor) products of the eigenfunctions on the factors. In the elliptic factors, the eigenfunctions are the normalized Hermite functions $`\gamma _q`$ (cf.\[F\], or \[Z.1\] for a context similar to the one here). In the real hyperbolic factors,the action operators are the generators of the unitary dilations on $`L^2()`$ given by $$U(\theta )f(x)=\theta ^{\frac{1}{2}}f(\theta x),(\theta ^+).$$ Their generalized eigenfunctions are the temperate distributions $$x_+^{\frac{1}{2}+ia},x_{}^{\frac{1}{2}+ia},(a)$$ and any $`fL^2(,dx)`$ has the eigenfunction expansion $$f(x)=_{}\widehat{f}_+(a)x_+^{\frac{1}{2}+ia}𝑑a+_{}\widehat{f}_{}(a)x_{}^{\frac{1}{2}+ia}𝑑a$$ with $`\widehat{f}_{\pm a}:=f,x_\pm ^{\frac{1}{2}+ia}.`$ In the complex hyperbolic (i.e. loxodromic) factors, the actions operators are given by the unitary dilations in polar coordinates $$U(\rho )f(r,\theta )=\rho f(\rho r,\theta ),(\rho ^+)$$ together with rotations. The joint eigenfunctions are the temperate distributions on $`^2`$ given by $$r^{it1}e^{in\theta },(t)$$ and in a notation similar to that of the real hyperbolic case a function $`fL^2(^2,rdrd\theta )`$ may be expressed in the form $$f(r,\theta )=\underset{n}{}_{}\widehat{f}(t,\theta )r^{it1}e^{in\theta }𝑑t.$$ For future reference we summarize the situation in the following table. \| Factor \| Action \| Eigenfunction \| \| --- \| --- \| --- \| \| $`L^2(S_L^1)`$ \| $`D_s`$ \| $`e^{is\frac{2\pi k}{L}}`$ \| \| $`L^2(^p)`$ \| $`I_j^e`$ \| Hermite functions $`\gamma _m,m^n;`$ \| \| $``$ \| $``$ \| $`\gamma _o(y)=\gamma _{iI}(y):=e^{\frac{1}{2}\|y\|^2}`$ \| \| $``$ \| $``$ \| $`\gamma _m:=C_ma_1^{m_1}\mathrm{}a_n^{m_n}\gamma _o`$ \| \| $`L^2(^q)`$ \| $`I_j^h`$ \| $`\mathrm{\Pi }_{j=1}^ry_{j\pm }^{ia_j\frac{1}{2}},a^r`$ \| \| $`L^2(^{2c})`$ \| $`I_j^{chRe},I_j^{ch,Im}`$ \| $`\mathrm{\Pi }_{j=1}^cr_j^{it_j1}e^{in_j\theta _j},t^c`$ \| As in the introduction, we put: $`(\mathrm{1.2.2.1})`$ $$H_{\alpha ,\lambda ,(\mu ,\nu )}:=[\underset{j=1}{\overset{p}{}}\alpha _j\widehat{I}_j^e+\underset{j=1}{\overset{q}{}}\lambda _j\widehat{I}_j^h+\underset{j=1}{\overset{c}{}}\mu _j\widehat{I}_j^{ch,Re}+\nu _j\widehat{I}_j^{ch,Im}]$$ $$:=\frac{1}{L}(LD_s+H_{\alpha ,\lambda ,(\mu ,\nu )})$$ and note that $`(\mathrm{1.2.2.2}a)`$ $$e^{is\frac{2\pi k}{L}}\gamma _m(x)[\mathrm{\Pi }_{j=1}^ry_{j\pm }^{ia_j\frac{1}{2}}][\mathrm{\Pi }_{j=1}^cr_j^{it_j1}e^{in_j\theta _j}]=r_{kmnat}e^{is\frac{2\pi k}{L}}\gamma _m(x)[\mathrm{\Pi }_{j=1}^ry_{j\pm }^{ia_j\frac{1}{2}}][\mathrm{\Pi }_{j=1}^cr_j^{it_j1}e^{in_j\theta _j}]$$ with $`(\mathrm{1.2.2.2}b)`$ $$r_{kmnat}=\frac{1}{L}(2\pi k+[\underset{j=1}{\overset{p}{}}\alpha _j(m_j+\frac{1}{2})+\underset{j=1}{\overset{r}{}}\lambda _ja_j+\underset{j=1}{\overset{c}{}}\mu _jt_j+\nu _jn_j]).$$ ## 2. The semi-classically scaled Laplacian The significance of the maximal abelian algebra $`𝒜_{pqc}`$ will appear as soon as we semi-classically ‘rescale’ the Laplacian and conjugate the principal part, the ‘linearized $`\sqrt{\mathrm{\Delta }}`$’, to its normal form. Let us briefly recollect this rescaling, which proceeds exactly as in the purely elliptic case \[Z.1,§2\]. We first prepare the Laplacian by putting it in Fermi normal coordinates $`(s,y)`$. It is then self-adjoint relative to the volume density $`J(s,u)\|ds\|\|dy\|`$ in these coordinates. To simplify, we then conjugate it to the unitarily equivalent (1/2-density-) Laplacian $$\mathrm{\Delta }_{1/2}:=J^{1/2}\mathrm{\Delta }J^{1/2},$$ which is self-adjoint with respect to the Lesbesgue density $`\|dsdy\|`$. We thus have: $`(\mathrm{2.1.1})`$ $$\mathrm{\Delta }_{1/2}=J^{1/2}_sg^{oo}J_sJ^{1/2}+\underset{ij=1}{\overset{n}{}}J^{1/2}_{y_i}g^{ij}J_{y_j}J^{1/2}$$ $$g^{oo}_s^2+\mathrm{\Gamma }^o_s+\underset{ij=1}{\overset{n}{}}g^{ij}_{u_i}_{y_j}+\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }^i_{y_i}+\sigma _o.$$ Semi-classical rescaling then involves two conjugations: First, by $`M_h=`$ multiplication by $`e^{\frac{is}{Lh}}`$, $$M_h^{}\mathrm{\Delta }M_h=(hL)^2g^{oo}+2i(hL)^1g^{oo}_s+i(hL)^1\mathrm{\Gamma }^o+\mathrm{\Delta }$$ and then by the semi-classical dilation $`T_hf(s,y)=f(s,h^{\frac{1}{2}}y)`$. The complete conjugation $`T_h^{}M_h^{}\mathrm{\Delta }M_hT_h`$ results in the semi-classically scaled Laplacian $`(2.1,2)`$ $$\mathrm{\Delta }_h=(hL)^2g_{[h]}^{oo}+2i(hL)^1g_{[h]}^{oo}_s+i(hL)^1\mathrm{\Gamma }_{[h]}^o+h^1(\underset{ij=1}{\overset{n}{}}g_{[h]}^{ij}_{y_i}_{y_j})+h^{\frac{1}{2}}(\underset{i=1}{\overset{n}{}}\mathrm{\Gamma }_{[h]}^i_{y_i})+(\sigma )_{[h]},$$ the subscript $`[h]`$ indicating to dilate the coefficients of the operator in the form, $`f_h(s,y):=f(s,h^{\frac{1}{2}}y).`$ Expanding the coefficients in Taylor series at $`h=0`$, we obtain the semi-classical expansion $`(\mathrm{2.1.3})`$ $$\mathrm{\Delta }_h\underset{m=0}{\overset{\mathrm{}}{}}h^{(2+m/2)}_{2m/2}$$ where $`_2=L^2,`$ $`_{3/2}=0`$ and where $`(\mathrm{2.1.4}).`$ $$_1=2L^1[i\frac{}{s}+\frac{1}{2}\{\underset{j=1}{\overset{n}{}}_{y_j}^2\underset{ij=1}{\overset{n}{}}K_{ij}(s)y_iy_j\}]$$ We will denote the bracketed operator, the ‘linearized $`\sqrt{\mathrm{\Delta }}`$’ by $``$. It is of order 1 in the sense of pseudodifferential operators (using the Weyl filtration in the transverse variables) and as will be seen below is the principal term in the semi-classical expansion of the square root of the rescaled Laplacian. (2.1.A) Appendix on metric scaling In addition to semi-classical scaling, we have also just introduced an independent scaling, metric scaling, which has to do with the behaviour of objects under dilations $`gϵ^2g`$ of the metric. As discussed in detail in \[Z.1\], the wave invariants have well-defined weights under metric rescaling and in analysing them it is very convenient to rescale all objects to be weightless. For instance, as discussed in \[Z.1, §1.4\], an $`\omega `$\- symplectic basis of Jacobi fields has weight $`\frac{1}{2}`$ and its time derivative has weight $`\frac{1}{2}.`$ To render it weightless a Jacobi eigenfield $`Y`$ should be replaced by $`L^{\frac{1}{2}}Y`$, $`\dot{Y}`$ by $`L^{\frac{1}{2}}\dot{Y}`$ etc. The resulting weightless Wronskian matrix is then denoted by $`𝒲_L`$. It is essentially the weightless matrix denoted $`𝒜_L`$ in \[Z.1\]. To render the coordinates $`(y,\eta )`$ weightless under metric rescaling, we also change variables to $`x=L^1y`$ and rewrite $`\mathrm{\Delta }_h`$ and the $`_{2\frac{n}{2}}`$’s in terms of the $`x`$-variables. For instance, $``$ then takes the form: $$=i\frac{}{s}+\frac{1}{2}[\underset{j=1}{\overset{n}{}}L^1_{x_j}^2\underset{ij=1}{\overset{n}{}}LK_{ij}(s)x_ix_j].$$ The symplectic coordinates on the symplectic normal space $`T^{}^n`$ to $`^+\gamma `$ will henceforth be denoted $`(x,\xi ).`$ For the sake of brevity we will not draw much attention to metric scaling in the various steps to come in the normal form algorithm. In all cases, the role of metric scaling is identical to that in the elliptic case \[Z.1\]. §2.2: Conjugating $`\mathrm{\Delta }_h`$ to the model We now conjugate the semi-classically scaled Laplacian $`\mathrm{\Delta }_h`$ from $`L^2(N_\gamma )`$ to the model $`L^2(S_L^1\times ^n)`$ by means of the moving metaplectic operator $`\mu (𝒲_L)`$, $$\mu (𝒲_L)f(s,y):=\mu (𝒲_L(s))f(s,y).$$ The motivation for this conjugation comes from: (2.2.1) Proposition $$=\mu (𝒲_L^{})D_s\mu (𝒲_L^{})^1$$ where $`𝒲_L`$ is the weightless Wronskian matrix and $`\mu `$ is the metaplectic representation. Proof: First, let us ignore the scaling parameter $`L`$, i.e. let us put $`L=1`$. The right side is then equal to $`(D_s+\mu (𝒲(s))^{}D_s\mu (𝒲(s)).`$ To evaluate the second term, we recall that the columns of $`𝒲`$ are Jacobi fields, and that Jacobi’s equation is equivalent to the linear system $`\frac{D}{ds}(Y,P)=JH(Y,P)`$ with $`P=\frac{DY}{ds},`$ and with $$H=\left(\begin{array}{cc}K\hfill & 0\hfill \\ 0\hfill & I\hfill \end{array}\right).$$ Hence, the second term is $`\frac{1}{i}d\mu (JH)`$ with $`d\mu `$ the derived metaplectic representation. But $`\frac{1}{i}d\mu (JH)=1/2(_{i=1}^n_{y_i}^2_{ij=1}^nK_{ij}(s)y_iy_j)`$ \[F\]. Re-inserting $`L`$ to make all objects weightless, we get the formed claimed above. ∎ Thus, conjugation by $`𝒲_L`$ puts the principal term $``$ of $`\mathrm{\Delta }_h`$ into the simple normal form $`D_s`$. This suggests conjugating the full rescaled Laplacian by $`\mu (𝒲_L)`$ to the ‘twisted model’ semi-classical Laplacian $`(\mathrm{2.2.2})`$ $$𝒟_h=\mu (𝒲_L^{})^1\mathrm{\Delta }_h\mu (𝒲_L^{})$$ which has the asymptotic expansion $`(\mathrm{2.2.3})`$ $$𝒟_h\underset{m=o}{\overset{\mathrm{}}{}}h^{(2+\frac{m}{2})}𝒟_{2\frac{m}{2}}$$ with $`𝒟_2=I,𝒟_{\frac{3}{2}}=0,𝒟_1=D_s`$. Thus, $`𝒟_h`$ is a small perturbation of $`D_s`$, and one may expect that perturbation theory can be used to find a good normal form for the whole of $`𝒟_h`$. Before doing so, we must consider which Hilbert space is the natural domain for $`𝒟_h`$. The point is that the conjugation has non-trivial monodromy (§1.1b) and hence the conjugate will act on functions transforming correctly under the monodromy group. We can describe the Hilbert space in terms of quantum mapping cylinders \[Z.1\]. First, we consider the holonomy aspect, put $`(\mathrm{2.2.4})`$ $$C_T^{\mathrm{}}(\times ^n):=\{fC^{\mathrm{}}(\times ^n):f(s+L,u)=\mu (T)f(s,u)\}$$ and let $`_T`$ denote its closure with respect to the obvious inner product over $`[0,L)\times ^n`$. Note that the metaplectic operator $`\mu (T)`$ is simply $$\mu (T)f(u)=f(t^1u)$$ and hence that $$C_T^{\mathrm{}}(\times ^n)C^{\mathrm{}}(N_\gamma )$$ where the isomorphism is simply the pull-back by the exponential map defined by the frame $`e(s)`$. In other words, expressed in terms of Fermi coordinates relative to a normal frame, $`L^2(N_\gamma )`$ becomes the quantum mapping cylinder of $`\mu (T^{}).`$ Let us note however that $`\mathrm{\Delta }_h`$ and hence all the $`_{2\frac{k}{2}}`$’s are invariant under under $`\mu (T)`$, so that it will play an insignificant role for our purposes. On the other hand, the quantized linear Poincare map $`\mu (P_\gamma )`$ will play an essential role. Hence we introduce its quantized mapping cylinder $`(\mathrm{2.2.5})`$ $$_\gamma :=\{fL_{loc}^2(\times ^n):\tau _Lf=\mu (P_\gamma )f\}$$ and note that $$\mu (𝒲_L):L^2(N_\gamma )_\gamma $$ is a unitary equivalence. Hence, the natural domain for $`𝒟_h`$ is the quantum mapping cyliner of $`\mu (P_\gamma )`$. In the calculation of traces, it is simpler to work in the original model $`L^2(S_L^1\times ^n)`$. Hence we will also consider the conjugate of $`𝒟_h`$ under a conjugation which untwists the mapping cylinder of $`\mu (P_\gamma ).`$ That is, we connect $`P_\gamma `$ to the identity by a segment of the one-parameter subgroup $`P_\gamma (s)`$ thru $`I`$ and $`P_\gamma `$, which exists by our non-degeneracy assumption on $`P_\gamma `$. Indeed, after diagonalizing $`P_\gamma `$ and consulting the list of symplectic equivalence classes of quadratic forms, we see that $$P_\gamma =exp(\mathrm{\Xi }_{H_{\alpha ,\lambda ,(\mu ,\nu )}})$$ where $`\mathrm{\Xi }_f`$ denotes the Hamilton vector field of $`f`$ and where $`exp\mathrm{\Xi }`$ denotes the exponential map from $`sp(n,)Sp(n,R)`$, with $`sp(n,)`$ viewed as the Poisson algebra of quadratic functions on $`^{2n}`$. We then have $$P_\gamma (s)=exp(s\mathrm{\Xi }_{H_{\alpha ,\lambda ,(\mu ,\nu )}})$$ and quantize this subgroup as $`(\mathrm{2.2.6})`$ $$\mu (P_\gamma (s))=e^{isH_{\alpha ,\lambda ,(\mu ,\nu )}}=e^{is[_{j=1}^p\alpha _jI_j^e+_{j=1}^q\lambda _jI_j^h+_{j=1}^c\mu _jI_j^{ch,Re}+\nu _jI_j^{ch,Im}]}.$$ Conjugation by $`\mu (P_\gamma )`$ transforms $`𝒟_h`$ into the model semi-classically scaled Laplacian $`(\mathrm{2.2.7})`$ $$_h:=\mu (P_\gamma )\mu (𝒲_L)\mathrm{\Delta }_h\mu (𝒲_L)^{}\mu (P_\gamma )^{}\underset{m=o}{\overset{\mathrm{}}{}}h^{(2+\frac{m}{2})}_{2\frac{m}{2}}$$ with $`_2=I,_{\frac{3}{2}}=0,`$ and with $`_1:=.`$ All coefficients of terms in $``$ are periodic in $`s`$ and have weight -2 under metric rescaling. ## 3. Semi-classical normal form We now wish to put $`_h`$ into semi-classical normal form, in the sense of \[Z.1, Lemma 2.22\]. This is the key transitional step in putting $`\mathrm{\Delta }`$ into microlocal normal form and is the source of the connections to local geometric invariants. The method is essentially the same as in the elliptic case, both in method and in detail. We therefore present only the first two steps in the proof and refer to \[Z.1, loc.cit\] for the inductive argument. As in the elliptic case, we state the result in terms of the $``$-operators since the trace will later be analysed in this model. However, most of the proof will take place in the twisted model, where the ‘linearized Laplacian’ is $`D_s`$ and the equations simplify most. In the following, the notation $`\|_o`$ means to restrict to functions in the kernel of $``$, that is, to elements of $``$-weight zero. In the twisted model, these are simply functions independent of $`s`$. In the passage from the semi-classical normal form to the microlocal (quantum Birkhoff) normal form, the various operators will only be applied to such weightless elements. This explains the rather complicated statement to follow; the result is only simple and natural when restricted to elements of weight zero. (3.1) Lemma (cf. \[Z.1, Lemma 2.22\] There exists an $`L`$-dependent $`h`$-pseudodifferential operator $`W_h=W_h(s,x,D_x)`$ on $`L^2(S_L^1\times ^n)`$ such that, for each $`sS_L^1`$, $$W_h(s,x,D_x):L^2(^n)L^2(^n)$$ is unitary, and such that $$W_h^{}_hW_hh^2L^2+2h^1L^1+\underset{j=0}{\overset{\mathrm{}}{}}h^{\frac{j}{2}}_{2\frac{j}{2}}^{\mathrm{}}(s,D_s,x,D_x)$$ where (i) $`_{2\frac{j}{2}}^{\mathrm{}}(s,D_s,x,D_x)=_{2\frac{j}{2}}^{\mathrm{},2}^2+_{2\frac{j}{2}}^{\mathrm{},1}+_{2\frac{j}{2}}^{\mathrm{},o},`$ with $`_{2\frac{j}{2}}^{\mathrm{},k}C^{\mathrm{}}(S_L^1,_ϵ^{j2k});`$ (ii) $`_{2j}^{\mathrm{}}(s,D_s,x,D_x)\|_o=_{2j}^{\mathrm{},o}(s,x,D_x)\|_o=f_j(I_1^e,\mathrm{},I_p^e,I_1^h,\mathrm{},I_r^h,I_1^{chRe},\mathrm{}I_c^{chRe},I_1^{chIm},\mathrm{},I_c^{chIm})\|_o`$ for certain polynomials $`f_j`$ of degree j+2 on $`^n,`$ i.e. $`f_j(I_1^e,\mathrm{},I_p^e,I_1^h,\mathrm{},I_r^h,I_1^{chRe},\mathrm{}I_c^{chRe},I_1^{chIm},\mathrm{},I_c^{chIm})𝒫_{}^{j+2}`$ (iii) $`_{2\frac{2k+1}{2}}^{\mathrm{}}(s,D_s,x,D_x)\|_o=_{2\frac{2k+1}{2}}^{\mathrm{},o}(s,x,D_x)\|_o=0;`$ (iv) The terms $`I_j^e,I_j^h,I_j^{ch}`$ are weightless under metric scalings and all of the $``$’s have weight -2. Proof: As in the elliptic case \[Z.1,Lemma 2.22\], the operator $`W_h`$ will be constructed as the asymptotic product $`(3.2)`$ $$W_h:=\mu (P_\gamma )^{}\mathrm{\Pi }_{k=1}^{\mathrm{}}W_{h\frac{k}{2}}\mu (P_\gamma )$$ of weightless unitary $`h`$-pseudodifferential operators on $`^n`$, with $`(3.3)`$ $$W_{h\frac{k}{2}}:=exp(ih^{\frac{k}{2}}Q_{\frac{k}{2}})$$ and with $`h^{\frac{k}{2}}Q_{\frac{k}{2}}h^{\frac{k}{2}}^{\mathrm{}}(S_L^1)^{k+2}`$ of total order 1. The product will converge, for each s, to a unitary operator in $`\mathrm{\Psi }_h^o(^n)`$ (we refer to \[Sj\] for a discussion of such asymptotic products). We first construct a weightless $`Q_{\frac{1}{2}}(s,x,D_x)C^{\mathrm{}}(S_L^1)_ϵ^3`$ such that $`(3.4a)`$ $$e^{ih^{\frac{1}{2}}Q_{\frac{1}{2}}}_he^{ih^{\frac{1}{2}}Q_{\frac{1}{2}}}\|_o=[h^2L^2+2h^1L^1+_o^{\frac{1}{2}}+\mathrm{}]\|_o$$ where the dots $`\mathrm{}`$ indicate higher powers in $`h`$. The operator $`Q_{\frac{1}{2}}`$ then must satisfy the commutation relation $`(3.4b)`$ $$\{[L^1,Q_{\frac{1}{2}}]+_{\frac{1}{2}}\}\|_o=0.$$ To solve for $`Q_{\frac{1}{2}}`$, we conjugate back to the $`𝒟_{2\frac{m}{2}}`$’s of the twisted model by $`\mu (P_\gamma )`$, which transforms $``$ into $`D_s`$. The commutation relation thus becomes $`(3.4c)`$ $$\{[L^1D_s,\mu (P_\gamma ))^{}Q_{\frac{1}{2}}\mu (P_\gamma )]+𝒟_{\frac{1}{2}}\}\|_o=0,$$ that is, $`(3.4d)`$ $$L^1_s\{\mu (P_\gamma )^{}Q_{\frac{1}{2}}\mu (P_\gamma ))\}\|_o=i\{𝒟_{\frac{1}{2}}\}\|_o$$ where $`_sA`$ is the Weyl operator whose complete symbol is the $`s`$-derivative of that of $`A`$. Since (3.4d) is simpler than (3.4b), we henceforth conjugate everything by $`\mu (P_\gamma ))`$, and relabel the operators $`\mu (P_\gamma )^{}Q\mu (P_\gamma )`$ by $`\stackrel{~}{Q}.`$ The resulting $`𝒟`$’s then transform under $`\tau _L`$ like operators on the quantum mapping cylinder of $`\mu (P_\gamma )`$. Our problem is thus to solve (3.4d) with an operator $`\stackrel{~}{Q}_{\frac{1}{2}}`$ satisfying $$\tau _L\stackrel{~}{Q}_{\frac{1}{2}}\tau _L^{}=\mu (P_\gamma )\stackrel{~}{Q}_{\frac{1}{2}}\mu (P_\gamma )^{}.$$ To solve the equation (3.4d) we rewrite it in terms of complete Weyl symbols. We will use the notation $`A(s,x,\xi )`$ for the complete Weyl symbol of the operator $`A(s,x,D_x)`$. Then (3.4d) becomes $`(3.5a)`$ $$L^1_s\stackrel{~}{Q}_{\frac{1}{2}}(s,x,\xi )=i𝒟_{\frac{1}{2}}\|_o(s,x,\xi )$$ with $$\stackrel{~}{Q}_{\frac{1}{2}}(s+L,x,\xi )=\stackrel{~}{Q}_{\frac{1}{2}}(s,P_\gamma (x,\xi )).$$ We solve (3.5a) with the Weyl symbol $$\stackrel{~}{Q}_{\frac{1}{2}}(s,x,\xi )=\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )+L_0^si𝒟_{\frac{1}{2}}\|_o(u,x,\xi )du$$ where $`\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )`$ is determined by the consistency condition $`(3.5b)`$ $$\stackrel{~}{Q}_{\frac{1}{2}}(L,x,\xi )\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )=L_0^Li𝒟_{\frac{1}{2}}\|_o(u,x,\xi )du$$ or in view of the periodicity condition in (3.5a), $`(3.5c)`$ $$\stackrel{~}{Q}_{\frac{1}{2}}(0,P_\gamma (x,\xi ))\stackrel{~}{Q}_{\frac{1}{2}}(0,x,\xi )=L_0^Li𝒟_{\frac{1}{2}}\|_o(u,x,\xi )du.$$ To solve, we use that $`𝒟_{\frac{1}{2}}\|_o(u,x,\xi )`$ is a polynomial of degree 3 in $`(x,\xi )`$. It will be most convenient to express this polynomial in coordinates relative to the eigenvectors of the Poincare map. In the elliptic planes, we use the complex coordinates $`z_j=x_j+i\xi _j`$ and $`\overline{z}_j=x_ji\xi _j`$ ($`j=1,\mathrm{},p`$) in which the action of $`P_\gamma `$ is diagonal. In the real hyperbolic planes we use the real coordinates $`(y_j,\eta _j)=(x_j,\xi _j),(j=p+1,\mathrm{}p+q)`$ in which the real hyperbolic part of $`P_\gamma `$ is diagonal. Finally, in the complex hyperbolic (loxodromic) 4-spaces we use the coordinates $`w_j=x_j+ix_{c+j},\overline{w}_j=x_jix_{c+j},\omega _j=\xi _ji\xi _{c+j},\overline{\omega _j}=\xi _j+i\xi _{c+j},(j=p+q+1,\mathrm{}p+q+c)`$ in which the complex hyperbolic part of $`P_\gamma `$ is diagonal. We will denote the Weyl symbols in these coordinates by their previous expressions. We also suppress the subscripts by using vector notation $`z,\overline{z},y,\eta ,w,\overline{w},\omega ,\overline{\omega }`$. Thus, (3.5c) becomes $$\stackrel{~}{Q}_{\frac{1}{2}}(0,e^{i\alpha }z,e^{i\alpha }\overline{z},e^\lambda y,e^\lambda \eta ,e^{\mu +i\nu }w,e^{\mu i\nu }\overline{w},e^{\mu +i\nu }\omega ,e^{\mu i\nu }\overline{\omega })\stackrel{~}{Q}_{\frac{1}{2}}(0,z,\overline{z},y,\eta ,w,\overline{w},\omega ,\overline{\omega })=$$ $`(3.6)`$ $$=L_0^Li𝒟_{\frac{1}{2}}\|_o(u,z,\overline{z},y,\eta ,w,\overline{w},\omega ,\overline{\omega })du.$$ We now use that $`𝒟_{\frac{1}{2}}(u,z,\overline{z},y,\eta ,w,\overline{w},\omega ,\overline{\omega })`$ is a polynomial of degree 3 to solve (3.5c). If we put $`(3.7a)`$ $$\stackrel{~}{Q}_{\frac{1}{2}}(s,z,\overline{z},y,\eta ,w,\overline{w},\omega ,\overline{\omega }):=\underset{\|a\|+\|\overline{a}\|+\|b\|+\|c\|+\|\overline{c}\|3}{}q_{\frac{1}{2};a\overline{a}bc\overline{c}}(s)z^a\overline{z}^{\overline{a}}y^{b_1}\eta ^{b_2}w^{c_1}\omega ^{c_2}\overline{w}^{\overline{c}_1}\overline{\omega }^{\overline{c}_2}$$ and $`(3.7b)`$ $$𝒟_{\frac{1}{2}}\|_o(s,z,\overline{z},y,\eta ,w,\overline{w},\omega ,\overline{\omega })du:=\underset{\|a\|+\|\overline{a}\|+\|b\|+\|c\|+\|\overline{c}\|3}{}d_{\frac{1}{2};a\overline{a}bc\overline{c}}(s)z^a\overline{z}^{\overline{a}}y^{b_1}\eta ^{b_2}w^{c_1}\omega ^{c_2}\overline{w}^{\overline{c}_1}\overline{\omega }^{\overline{c}_2}$$ then (3.6) becomes $$\underset{\|a\|+\|\overline{a}\|+\|b\|+\|c\|+\|\overline{c}\|3}{}(1e^{i(a\overline{a})\alpha +i(c_1\overline{c}_1)\nu +(b_1b_2)\lambda +(c_2\overline{c}_2)\mu })q_{\frac{1}{2};a\overline{a}bc\overline{c}}(0)z^a\overline{z}^{\overline{a}}y^{b_1}\eta ^{b_2}w^{c_1}\omega ^{c_2}\overline{w}^{\overline{c}_1}\overline{\omega }^{\overline{c}_2}=$$ $`(3.8).`$ $$=iL^2\underset{\|a\|+\|\overline{a}\|+\|b\|+\|c\|+\|\overline{c}\|3}{}\overline{d}_{\frac{1}{2};a\overline{a}bc\overline{c}}z^a\overline{z}^{\overline{a}}y^{b_1}\eta ^{b_2}w^{c_1}\omega ^{c_2}\overline{w}^{\overline{c}_1}\overline{\omega }^{\overline{c}_2}$$ Under the non-degeneracy assumption on $`P_\gamma `$, we can solve with $`(3.9)`$ $$q_{\frac{1}{2};a\overline{a}bc\overline{c}}(0)=iL^2(1e^{i(a\overline{a})\alpha +i(c_1\overline{c}_1)\nu +(b_1b_2)\lambda +(c_2\overline{c}_2)\mu })^1d_{\frac{1}{2};a\overline{a}bc\overline{c}}$$ since $`i(a\overline{a})\alpha +i(c_1c_2)\nu )+(b_1b_2)\lambda +(\overline{c}_1\overline{c}_2)\mu )=2\pi ik`$ only if $`a=\overline{a},b_1=b_2,c=\overline{c}`$ and there are no such $`(a,\overline{a},b_1,b_2,c,\overline{c})`$ in an odd-index equation. Precisely as in the purely elliptic case of \[Z.1\], we see that $`\stackrel{~}{Q}_{\frac{1}{2}}`$ is a pseudodifferential operator on $`^n`$ with the same order, same order of vanishing, and same parity as the restriction of $`𝒟_{\frac{1}{2}}`$ to elements of weight zero. We then extend it as a pseudodifferential operator of the form $$\stackrel{~}{Q}_{\frac{1}{2}}\mathrm{\Psi }^o(^1)_ϵ^3$$ on all of $`_\gamma `$ by decreeing that it commute with $`s`$. The conjugate by $`\mu (P_\gamma )`$ then defines a unitary operator $`W_{h\frac{1}{2}}\mathrm{\Psi }_h^o(S^1\times ^n)`$ satisfying (3.4a). The corresponding twisted unitary operator with exponent $`\stackrel{~}{Q}_{\frac{1}{2}}`$, i.e. the image of $`W_{h\frac{1}{2}}`$ under conjugation by $`\mu (P_\gamma )`$, will be denoted $`\stackrel{~}{W}_{h\frac{1}{2}}.`$ The effect of this first conjugation is precisely as in the elliptic case: Since $`h^{\frac{1}{2}}\stackrel{~}{Q}_{\frac{1}{2}}`$ is of total order 1, $`h^{\frac{1}{2}}ad(\stackrel{~}{Q}_{\frac{1}{2}})`$ (with $`ad(A)B:=[B,A]`$) preserves the total order in $`\mathrm{\Psi }_h^{(,,)}`$, and hence $`\stackrel{~}{W}_{h\frac{1}{2}}`$ is an order-preserving automorphism of the model pseudodifferential algebra. It is moreover independent of $`D_s`$ and has an odd polynomial Weyl symbol, so that $`(3.10)`$ $$h^{\frac{1}{2}}ad(\stackrel{~}{Q}_{\frac{1}{2}}):h^{\frac{k}{2}}\mathrm{\Psi }^l()_ϵ^mh^{\frac{k+1}{2}}[\mathrm{\Psi }^{l1}()_ϵ^{m+3}+\mathrm{\Psi }^l()_ϵ^{m+1}].$$ Finally, the $`d_{\frac{1}{2};m,n}`$’s have weight -2, the variables $`z`$ have weight 0 and hence the $`q_{\frac{1}{2};m,n}`$’s have weight 0. Consider now the element $$𝒟_h^{\frac{1}{2}}:=\stackrel{~}{W}_{h\frac{1}{2}}^{}𝒟_h\stackrel{~}{W}_{h\frac{1}{2}}\mathrm{\Psi }_h^2(^1\times ^n)$$ which can be expanded in the semi-classical series $`(3.11)`$ $$𝒟_h^{\frac{1}{2}}\underset{n=o}{\overset{\mathrm{}}{}}h^{2+\frac{n}{2}}\underset{j+m=n}{}\frac{i^j}{j!}(ad\stackrel{~}{Q}_{\frac{1}{2}})^j𝒟_{2\frac{m}{2}}$$ $$:=h^2L^2+h^1L^1D_s+\underset{n=3}{\overset{\mathrm{}}{}}h^{2+\frac{n}{2}}𝒟_{2\frac{n}{2}}^{\frac{1}{2}}.$$ An obvious induction as in the elliptic case \[loc.cit.\] gives that $$ad(\stackrel{~}{Q}_{\frac{1}{2}})^j𝒟_{2\frac{m}{2}}C^{\mathrm{}}(,_ϵ^{m+j4})D_s^2+C^{\mathrm{}}(,_ϵ^{m+j2})D_s+C^{\mathrm{}}(,_ϵ^{m+j}).$$ It follows that $`𝒟_{2\frac{n}{2}}^{\frac{1}{2}}`$ has the same filtered structure as $`𝒟_{2\frac{n}{2}}.`$ We carry this procedure out one more step before referring to \[Z.1\] for the inductive argument, since the even steps behave differently from the odd ones. We thus seek an element $`\stackrel{~}{Q}_1(s,x,D_x)\mathrm{\Psi }^{}(S^1\times ^n)`$ and an element $`f_o(I_1^e,\mathrm{},I_p^e,I_1^h,\mathrm{},I_r^h,I_1^{ch,Re},I_1^{ch,Im},\mathrm{},I_c^{ch,Re},I^{ch,Im})𝒜`$ so that $`(3.12)`$ $$𝒟_h^1:=\stackrel{~}{W}_{h1}^{}𝒟^{\frac{1}{2}}\stackrel{~}{W}_{h1}=h^2L^2+h^1L^1D_s+h^{\frac{1}{2}}𝒟_{\frac{1}{2}}^{\frac{1}{2}}+𝒟_o^1(s,D_s,x,D_x)+\mathrm{}$$ with $`(3.13a)`$ $$𝒟_o^1(s,D_s,x,D_x)\|_o=f_o(I_1^e,\mathrm{},I_c^{ch,Im})$$ with $`\stackrel{~}{W}_{h1}=e^{ih\stackrel{~}{Q}_1},`$ and where the dots signify terms of higher order in $`h`$. Note that $`𝒟_{\frac{1}{2}}^1=𝒟_{\frac{1}{2}}^{\frac{1}{2}}`$, so that (3.12) implies that $`(3.13b)`$ $$\{h^{\frac{1}{2}}𝒟_{\frac{1}{2}}^1+𝒟_o^1\}\|_o=f_o(I_1^e,\mathrm{},I_c^{ch,Im}).$$ The condition on $`\stackrel{~}{Q}_1`$ is then $`(3.14a)`$ $$\{[D_s,\stackrel{~}{Q}_1]+𝒟_o^{\frac{1}{2}}\}\|_o=f_o(I_1^e,\mathrm{},I_c^{ch,Im})$$ or equivalently $`(3.14b).`$ $$_s\stackrel{~}{Q}_1\|_o=\{𝒟_o^{\frac{1}{2}}+f_o(I_1^e,\mathrm{},I_c^{ch,Im})\}\|_o$$ We solve (3.14b) by again expressing everything in terms of complete Weyl symbols relative to the eigenvector coordinates. Thus we rewrite (3.14b)) in the form $$L^1_s\stackrel{~}{Q}_1(s,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })=$$ $`(3.15a)`$ $$=i\{𝒟_o^{\frac{1}{2}}\|_o(s,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })f_o(\|z_1\|^2,\mathrm{},\|z_p\|^2,y_1\eta _1,\mathrm{},y_r\eta _r,Rew\omega ,Imw\omega )\}$$ or equivalently $$\stackrel{~}{Q}_1(s,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })=\stackrel{~}{Q}_1(0,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })$$ $`(3.15b)`$ $$iL_0^s[𝒟_o^{\frac{1}{2}}\|_o(u,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })f_o(\|z_1\|^2,\mathrm{},\|z_p\|^2,y_1\eta _1,\mathrm{},y_r\eta _r,Rew\omega ,Imw\omega )]𝑑u$$ and solve simeltaneously for $`\stackrel{~}{Q}_1`$ and $`f_o`$. The consistency condition determining a unique solution is that $$\stackrel{~}{Q}_1(L,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })=\stackrel{~}{Q}_1(0,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })$$ $`(3.16a)`$ $$iL_0^L[𝒟_o^{\frac{1}{2}}\|_o(u,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })f_o(\|z_1\|^2,\mathrm{},\|z_p\|^2,y_1\eta _1,\mathrm{},y_q\eta _q,Rew\omega ,Imw\omega )]𝑑u.$$ or in view of the twisted periodicity condition $$\stackrel{~}{Q}_1(0,e^{i\alpha }z,e^{i\alpha }\overline{z},e^\lambda y,e^\lambda \eta ,e^{\mu +i\nu }w,e^{\mu i\nu }\omega ,e^{\mu +i\nu }\overline{w},e^{\mu i\nu }\overline{\omega })\stackrel{~}{Q}_1(0,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })=$$ $`(3.16b)`$ $$iL\{_0^L𝒟_o^{\frac{1}{2}}\|_o(u,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })duLf_o(\|z_j\|^2,y_j\eta _j,Rew_m\omega _m,Imw_n\omega _n)\}.$$ In the spirit of the previous step, we use that $`𝒟_o^{\frac{1}{2}}\|_o(u,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })`$ is a polynomial of degree 4 to solve the equation. We put $$\stackrel{~}{Q}_1(s,z,\overline{z},y,\eta ,w,\omega ,\overline{w},\overline{\omega })=\underset{\|a\|+\|\overline{a}\|+\|b\|+\|c\|+\|\overline{c}\|4}{}q_{1;a\overline{a}bc\overline{c}}(s)z^a\overline{z}^{\overline{a}}y^{b_1}\eta ^{b_2}w^{c_1}\omega ^{c_2}\overline{w}^{\overline{c}_1}\overline{\omega }^{\overline{c}_2},$$ and in an abbreviated notation, $`(3.17)`$ $$f_o(\|z\|^2,y\eta ,Rew\omega ,Imw\omega )=\underset{\|k\|+\|\mathrm{}\|+\|n\|2}{}c_{ok\mathrm{}n}\|z\|^{2k}(y\eta )^{\mathrm{}}(Rew\omega )^{n_1}(Imw\omega )^{n_2},$$ and $$𝒟_o^{\frac{1}{2}}\|_o(s,z,\overline{z},y,\eta )du:=\underset{\|a\|+\|\overline{a}\|+\|b\|+\|c\|+\|\overline{c}\|4}{}d_{o;a\overline{a}bc\overline{c}}^{\frac{1}{2}}(s)z^a\overline{z}^{\overline{a}}y^{b_1}\eta ^{b_2}w^{c_1}\omega ^{c_2}\overline{w}^{\overline{c}_1}\overline{\omega }^{\overline{c}_2},$$ and finally $$\overline{d}_{o;a\overline{a}bc\overline{c}}^{\frac{1}{2}}:=\frac{1}{L}_o^Ld_{o;a\overline{a}bc\overline{c}}^{\frac{1}{2}}(s)𝑑s$$ As above, we can solve for the off-diagonal coefficients where either $`ab`$ or $`mn`$ $`(3.18a)`$ $$q_{1;a\overline{a}bc\overline{c}}(0)=iL^2(1e^{i(a\overline{a})\alpha +i(c_1\overline{c}_1)\nu +(b_1b_2)\lambda )+(c_2\overline{c}_2)\mu })^1\overline{d}_{o;a\overline{a}bc\overline{c}}^{\frac{1}{2}}$$ and must set the diagonal coefficients with $`a=\overline{a},c_1=\overline{c}_1,b_1=b_2,c_2=\overline{c}_2`$ equal to zero. The expression in (3.18a) is well-defined by the non-degeneracy assumption. The coefficients $`c_{ok\mathrm{}}`$ are then determined by $`(3.18b)`$ $$c_{ok\mathrm{}n}=\overline{d}_{1;kk\mathrm{}\mathrm{}nn}^{\frac{1}{2}}.$$ It is evident that $`\stackrel{~}{Q}_1`$ and $`f_o(I_1^e,\mathrm{},I_c^{ch,Im})`$ are even polynomial pseudodifferential operators of degree 4 in the variables $`(x,D_x)`$, that $`\stackrel{~}{Q}_1`$ is weightless under metric rescalings and that the coefficients $`c_{ok\mathrm{}n}`$ are of weight -2. The rest proceeds as in the elliptic case.∎ ## 4. Normal form of the Laplacian: Proof of Theorem I We now use the semi-classical normal forms to put the Laplacian into quantum Birkhoff normal form. Essentially this amounts to taking direct sums (or integrals) of the semi-classical normal form over various internal Planck constants. Proof of Theorem I: As in the elliptic case, we make the transition from the semi-classical normal form to the quantum Birkhoff normal form by using genreralized eigenfunction expansions for the model algebra. From the table in §1.2.1 we see that a function $`fL^2(S_L^1\times _x^p\times _y^q\times _{r,\theta }^{2c})`$ can be expanding in terms of joint $`𝒜_{pqc}`$-eigenfunctions as: $$f(s,x,y,r,\theta )=\underset{\pm }{}\underset{(k,m,n)𝐍^{1+p+c}}{}_^q_{^{+c}}\widehat{f}_\pm (k,m,n,a,t)e^{ir_{kmnat}}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}r^{it1}𝑑a𝑑t.$$ Here as in §3, we have used the notation $`x`$ for linear coordinates on the elliptic factors, $`y`$ for those on the real hyperbolic factors, and polar coordinates $`w_j=r_je^{i\theta _j}`$ in each $`w_j`$-plane of the complex hyperbolic factors. We also employ a multi-index notation. We now assemble the semi-classical intertwining operators into the Fourier-Hermite-Mellin -series-integral intertwining operator $`(4.1)`$ $$W_\gamma :L^2(S_L^1\times ^p\times ^q\times ^{2c},dsdxdydw),L^2(S_L^1\times ^p\times ^q\times ^{2c},dsdxdydw)$$ $$W_\gamma \underset{\pm }{}\underset{(k,m,n)𝐍^{1+p+c}}{}_^q_{^{+c}}\widehat{f}(k,m,n,a,t)e^{ir_{kmnat}s}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}r^{it1}𝑑a𝑑t=$$ $$\underset{\pm }{}\underset{(k,m,n)𝐍^{1+p+c}}{}_^q_{^{+c}}\widehat{f}(k,m,n,a,t)e^{ir_{kmnrt}s}W_{kmnrt}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ir\frac{1}{2}}r^{it1}𝑑a𝑑t$$ with $$W_{kmnrt}:=\mu (\stackrel{~}{𝒲}(s)^{})W_{r_{kqnat}^1}\mu (\stackrel{~}{𝒲}_s)^1.$$ Also, the dilation operators will be assembled into the operator $`(4.2)`$ $$T:L^2(S_L^1\times ^p\times ^q\times ^{2c},dsdxdydw)L^2(S_L^1\times ^p\times ^q\times ^{2c},dsdxdydw)),$$ $$T\underset{\pm }{}\underset{(k,m,n)𝐍^{1+p+c}}{}_^q_{^{+c}}\widehat{f}(k,m,n,a,t)e^{ir_{kmnat}s}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}r^{it1}𝑑a𝑑t=$$ $$=\underset{\pm }{}\underset{(k,m,n)𝐍^{1+p+c}}{}_^q_{^{+c}}\widehat{f}(k,m,n,a,t)r_{kmnrt}^{i\frac{1}{2}(\|a\|+\|t\|)}e^{ir_{kmnat}s}e^{i<n,\theta >}\gamma _m(\sqrt{r_{kmnrt}}x)y_\pm ^{ia\frac{1}{2}}r^{it1}𝑑a𝑑t.$$ Here, we used that the hyperbolic eigenfunctions are eigenfunctions of dilation operators. It follows, formally, from the semi-classical normal form and from the eigenfunction expansion that $`(4.3)`$ $$W_\gamma ^1T^1\mathrm{\Delta }TW_\gamma ^2+f_o(I_1^e,\mathrm{},I_{2c}^{ch,Im})+\frac{f_1(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{}+\mathrm{}.$$ We now show that the intertwining operator is actually a standard Fourier Integral operator (in the Weyl operator, or isotropic, sense) and that (4.3) holds modulo the the kind of error stated in Theorem I. The proof is again similar to the elliptic case, so we concentrate on the novel aspects and refer the reader to \[Z.1, Proposition 3.4\] for the remaining details. As before, we will not be as careful here as in \[Z.1\] to express things in weightless terms relative to metric rescalings. (4.4) Propostion $`TW_\gamma T^1`$ is a (standard) Fourier integral operator, well-defined and invertible on the microlocal neighborhood (0.1) in $`T^{}(S_L^1\times ^n)`$. Sketch of Proof: We first consider the unitarily equivalent operator $`\stackrel{~}{T}\stackrel{~}{W}\stackrel{~}{T}^1`$ in the microlocal neighborhood (1.2.1) in the twisted model, with $`(4.5)`$ $$\stackrel{~}{W}:_\alpha _\alpha $$ $$\stackrel{~}{W}(e^{ir_{kmnat}s}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}\rho ^{it1}):=e^{ir_{kmnat}s}W_{r_{kmnat}^1}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}r^{it1},$$ and with $`\stackrel{~}{T}`$ the dilation operator analogous to (4.2) but relative to the basis $`e^{ir_{kmnat}s}e^{i<n,\theta >}\gamma _q(x)y_\pm ^{ia\frac{1}{2}}r^{it1}`$. We then factor $`\stackrel{~}{T}\stackrel{~}{W}\stackrel{~}{T}^1`$ as the product $`\stackrel{~}{T}\stackrel{~}{W}\stackrel{~}{T}^1=j^{}\stackrel{~}{T}V\stackrel{~}{T}^1`$ where: $`(4.6)`$ $$V:_\alpha L_{loc}^2(\times \times ^n)$$ $$V=\mathrm{\Pi }_{j=o}^{\mathrm{}}exp[iD_s^{\frac{j}{2}}Q_{\frac{j}{2}}(s^{},y,D_y)]$$ that is, $$Ve^{ir_{kmnat}s}e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}r^{it1}):=e^{ir_{kmnat}s}W_{r_{kmnat}^1}(s^{},x,y,w,D_x,D_y,D_w)e^{i<n,\theta >}\gamma _m(x)y_\pm ^{ia\frac{1}{2}}r^{it1},$$ and where $`(4.7)`$ $$j^{}:C^{\mathrm{}}(\times \times ^n)C^{\mathrm{}}(\times ^n)$$ $$j^{}f(s,x)=f(s,s,x)$$ is the pullback under the partial diagonal embedding. The discussion of $`V`$ and the proof that $`j^{}V`$ is a standard Fourier Integral operators goes precisely as in \[Z.1, Proposition 3.4\]. The effect of the dilation is to convert the isotropic calculus into the pure polyhomogenous calculus (see also \[G.1\] for this aspect) and then the power of $`D_s`$ insures that the phases are all homogeneous of degree 1 and vanishing to higher and higher order along $`\gamma `$ (by one step as the index j increases by one unit). Hence the phase of the infinite product has only finitely many terms of a given vanishing order and converges as a formal power series in the transverse variable. A convergent product can be defined (by Borel summation) of the phase (cf. \[Sj\]). The proposition then follows by expressing $$TW_\gamma T^1=T\mu (𝒲)^{}\stackrel{~}{T}^1\stackrel{~}{T}W\stackrel{~}{T}^1\stackrel{~}{T}\mu (𝒲)T^1$$ and noting that $`\stackrel{~}{T}\mu (𝒲)T^1`$ is also a standard Fourier Integral operator. ∎ We now complete the proof of the quantum normal form Theorem I for $`\sqrt{\mathrm{\Delta }}`$, stated in an equivalent form in terms of $`W_\gamma .`$ As in the introduction, the notation $`AB`$ means that the complete (Weyl) symbol of $`AB`$ vanishes to infinite order at $`\gamma `$ and $`O_j\mathrm{\Psi }^m`$ denotes the pseudodifferential operators of order m whose Weyl symbols vanish to order j at $`(y,\eta )=(0,0)`$. Here, pseudodifferential operator can refer to either the standard polyhomogeneous kind, or to the mixed polyhomogeneous-isotropic kind as in $`\mathrm{\Psi }^k(S_L^1)𝒲^l`$, in which case the total order is defined to be $`m=k+l`$. To simplify notation, we will denote the space of mixed operators of order m by $`\mathrm{\Psi }_{mx}^m(S_L^1\times ^n)`$. (4.8) Lemma Let $`TW_\gamma T^1`$ be the Fourier Integral operator of Proposition (4.4), defined over a conic neighborhood of $`R^+\gamma `$ in $`T^{}(S_L^1\times ^n)`$. Then: $$W_\gamma ^1T^1\sqrt{\mathrm{\Delta }}TW_\gamma P_1(,I_1^e,\mathrm{},I_{2c}^{ch,Im})+P_o(,I_1^e,\mathrm{},I_{2c}^{ch,Im})+\mathrm{}\text{mod}_{k=o}^{m+1}O_{2(m+1k)}\mathrm{\Psi }_{mx}^{1k}(S_L^1\times ^n),$$ where $`(4.9)`$ $$P_1(,I_1^e,\mathrm{},I_{2c}^{ch,Im})+\frac{p_1^{[2]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{L}+\frac{p_2^{[3]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^2}+\mathrm{}$$ $$P_m(,I_1^e,\mathrm{},I_{2c}^{ch,Im})\underset{k=m}{\overset{\mathrm{}}{}}\frac{p_k^{[km]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^j}$$ with $`p_k^{[km]}`$, for m=-1,0,1,…, homogenous of degree l-m in the variables $`(I_1^e,\mathrm{},I_{2c}^{ch,Im})`$ and of weight -1. Proof: As a semi-classical expansion in the “parameter” $`h=\frac{1}{L}`$, (4.3) may be rewritten in the form : $`(4.10)`$ $$W_\gamma ^1T^1\sqrt{\mathrm{\Delta }}TW_\gamma +\frac{p_1(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{L}+\frac{p_2(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^2}+\mathrm{}.$$ From the fact that the numerators $`f_j(I_1^e,\mathrm{},I_{2c}^{ch,Im})`$ in (4.3) are polynomials of degree j+2 and of weight -2, the numerators $`p_k(I_1^e,\mathrm{},I_{2c}^{ch,Im})`$ are easily seen to be polynomials of degree $`k+1`$ and of weight -1. Hence they may be expanded in homogeneous terms $`(4.11)`$ $$p_k=p_k^{[k+1]}+p_k^{[k]}+\mathrm{}p_k^{[o]},$$ with $`p_k^{[j]}`$ the term of degree j and still of weight -1. The right side of (4.12) can then be expressed as a sum of homogeneous operators: $`(4.12)`$ $$P_1(,I_1^e,\mathrm{},I_{2c}^{ch,Im})+P_o(,I_1^e,\mathrm{},I_{2c}^{ch,Im})+\mathrm{}$$ with $`(4.13)`$ $$P_1(,I_1^e,\mathrm{},I_{2c}^{ch,Im})+\frac{p_1^{[2]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{L}+\frac{p_2^{[3]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^2}+\mathrm{}$$ $$P_m(,I_1^e,\mathrm{},I_{2c}^{ch,Im})\underset{k=m}{\overset{\mathrm{}}{}}\frac{p_k^{[km]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^k}.$$ We claim that: $`(4.14)`$ $$W_\gamma ^1T^1\sqrt{\mathrm{\Delta }}TW_\gamma [+\frac{p_1(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{L}+\frac{p_2(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^2}+\mathrm{}+\frac{p_m(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^m}]$$ $$_{k=o}^{m+1}O_{2(m+1k)}\mathrm{\Psi }_{mx}^{1k}(S_L^1\times ^n).$$ Indeed, from the analysis of the remainder terms in the semi-classical normal form (see Lemma (3.1 (i)) and \[Z.1, Lemma 2.22\]), we have $`(4.15)`$ $$P_1(,I_1^e,\mathrm{},I_{2c}^{ch,Im})[+\frac{p_1^{[2]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{L}+\mathrm{}+\frac{p_N^{[N+1]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^k}]$$ $$O_{2(N+2)}\mathrm{\Psi }_{mx}^1(S_L^1\times ^n)$$ and also $`(4.16)`$ $$P_m(,I_1^e,\mathrm{},I_{2c}^{ch,Im})\underset{k=m}{\overset{N}{}}\frac{p_k^{[km]}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^k}O_{2(N+1m)}\mathrm{\Psi }_{mx}^m(S_L^1\times ^n).$$ Hence the expansion (4.10) is also asymptotic in the sense of $`.`$ For the statement of Theorem I in the introduction, it is only necessary to conjugate under $`\mu (\stackrel{~}{W}).`$ The rest proceeds as in the elliptic case.∎ ## 5. Wave invariants and residue trace: Proof of Theorem B The purpose of this section is to show that the wave invariants have precisely the same relation to the coefficients of the quantum normal form in the non-degenerate case that they have in the elliptic one. The characterization of the wave invariants in Theorem I will then follow from Theorem A of \[Z.1\]. We will need to use some further notation and results from \[Z.1\]: First, the kth wave invariant of a positive elliptic operator $`P`$ at a non-degenerate closed bicharacteristic $`\gamma `$ will be denoted $`\tau _{\gamma k}(P).`$ According to \[Z.1, Proposition 4.2\] we then have: $`(5.1)`$ $$\tau _{\gamma k}(P)=\tau _{\gamma k}(P_1^{2k+4}+P_o^{2k+2}+\mathrm{}+P_{k1}^o)$$ where $`P_j^k`$ denotes the first k terms in the Taylor expansion of the jth homogeneous part of the complete symbol of $`P`$ at $`\gamma `$. Thus, $`\tau _{\gamma k}(P)`$ involves the (2k+4)th jet of the principal symbol, the (2k+2)-jet of the subprincipal term, …, up to the zero-jet of term of homogeneity order (-k-1). As in \[Z.1, §4\], we will also rewrite the normal form in terms of of $`D_s`$ and $`H_{\alpha ,\lambda ,(\mu ,\nu )}`$ using that $$\frac{p_\nu (I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(L)^\nu }=\frac{p_\nu (I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(LD_s)^\nu }(I\nu \frac{H_{\alpha ,\lambda ,(\mu ,\nu )}}{LD_s}+\frac{1}{2}\nu (\nu 1)(\frac{H_{\alpha ,\lambda ,(\mu ,\nu )}}{LD_s})^2+\mathrm{}).$$ By (5.1), we can drop the $`D_s^{(k+1)+\nu }H_{\alpha ,\lambda ,(\mu ,\nu )}^{k+1\nu }`$ and higher terms, so $`𝒟_{k+1}`$ can be written in the form $`(5.2)`$ $$𝒟_{k+1}LD_s+H_{\alpha ,\lambda ,(\mu ,\nu )}+\frac{\stackrel{~}{p}_1(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{LD_s}+\frac{\stackrel{~}{p}_2(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(LD_s)^2}+\mathrm{}+\frac{\stackrel{~}{p}_{k+1}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(LD_s)^{k+1}}$$ modulo terms which make no contribution to $`\tau _{k\gamma }`$. We then use the fact (\[Z.1, (4.3)\],\[Z.2\]) that $`(5.3)`$ $$\tau _{\gamma k}(P)=resD_t^k\psi _ϵ(D_s,y,D_y)e^{itP}\|_{t=L}$$ where $`res`$ is the non-commutative residue and where $`\psi _ϵ(D_s,y,D_y)`$ denotes a microlocal cut-off to the cone (1.2.1). Note that in contrast to the elliptic case, the microlocal cut-off cannot be constructed in $`𝒜_{p,q,c}`$ since the neighboorhoods given by $`I<ϵ\sigma `$ in terms of mixed hyperbolic-elliptic actions are of infinite transverse symplectic volume. This does not pose a genuine problem, but accounts for a number of modifications to the elliptic case in \[Z.1\]. For the gauging elliptic operator we use $`LD_s`$. To simplify the notation we will put $`(5.4)`$ $$𝒫_{k+1}(D_s,I_1^e,\mathrm{},I_{2c}^{ch,Im}):=\frac{\stackrel{~}{p}_1(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{LD_s}+\frac{\stackrel{~}{p}_2(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(LD_s)^2}+\mathrm{}+\frac{\stackrel{~}{p}_{k+1}(I_1^e,\mathrm{},I_{2c}^{ch,Im})}{(LD_s)^{k+1}}.$$ so that: $`(5.5)`$ $$\tau _{\gamma k}(\sqrt{\mathrm{\Delta }})=Res_{z=0}TrD_t^k\psi _ϵ(D_s,y,D_y)e^{it[\frac{1}{L}(2\pi LD_s+H_{\alpha ,\lambda ,(\mu ,\nu )})+𝒫_{k+1}]}(LD_s)^z\|_{t=L}.$$ As in the elliptic case, a key role will be played by the (formal) trace $$T(\alpha ,\lambda ,(\mu ,\nu )):=Tre^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}.$$ Its precise definition is the following: Since $`e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}`$ is an element of the metaplectic representation $`\mu `$ of $`Mp(n,)`$, $`T(\alpha ,\lambda ,(\mu ,\nu ))`$ may be identified with the character $`Ch`$ of $`\mu `$ evaluated at the associated element $`P_\gamma =exp(\mathrm{\Xi }_{H_{\alpha ,\lambda ,(\mu ,\nu )}})Mp(n,).`$ Here, $`H_{\alpha ,\lambda ,(\mu ,\nu )}`$ denotes the quadratic function on $`^{2n}`$ which gives the complete Weyl symbol of the corresponding action operator, and as above $`exp\mathrm{\Xi }`$ denotes the flow at time 1 of its Hamilton vector field, or, more correctly, the lift to $`Mp(n,)`$ which corresponds to $`e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}`$ under $`\mu `$. Since $`Mp(n,)`$ is a semi-simple Lie group, the character $`Ch`$ is a real analytic function on the open dense subset $`Mp(n,)_{reg}`$ of regular elements of $`Mp(n,),`$ where it is given by the Harish-Chandra formula \[Kn\]. We will need below the explicit formula for $`Ch(x)`$ in terms of the eigenvalues of $`x`$. For elements of $`Mp(n,)`$ not having 1 as an eigenvalue, we recall that $`Ch(x)`$ is given by $$Ch(x)=\frac{i^\sigma }{\sqrt{\|det(Ix)\|}}$$ where $`\sigma `$ is a certain Maslov index. For non-degenerate $`x`$ with p pairs of eigenvalues $`e^{\pm i\alpha _j}`$ of modulus one, q pairs of positive real eigenvalues $`e^{\pm \lambda _j}`$ and c quadruplets of eigenvalues $`e^{\pm (\mu _j\pm i\nu _j)}`$, $`Ch(x)`$ is therefore given (up to a Maslov factor) by $`(5.6)`$ $$T(\alpha ,\lambda ,(\mu ,\nu ))=\mathrm{\Pi }_{j=1}^p\frac{e^{\frac{1}{2}i\alpha _j}}{1e^{i\alpha _j}}\mathrm{\Pi }_{j=1}^q\frac{e^{\frac{1}{2}\lambda _j}}{1e^{\lambda _j}}\mathrm{\Pi }_{j=1}^c\frac{e^{\frac{1}{2}(\mu _j+i\nu _j)}}{1e^{\mu _j+i\nu _j}}\frac{e^{\frac{1}{2}(\mu _ji\nu _j)}}{1e^{\mu _ji\nu _j}}.$$ Here we have selected one eigenvalue $`\rho `$ from each symplectic pair $`\rho ,\rho ^1`$ (see §1.1-2). The ambiguity is fixed by the Maslov factor $`i^\sigma `$, which can (and will) be ignored below for the sake of brevity. We can now give: Proof of Theorem B: Since $`e^{2\pi iLD_s}I`$ on $`L^2(S_L^1)`$ we have $$a_{k\gamma }=\tau _{\gamma k}(\sqrt{\mathrm{\Delta }})=$$ $`(5.7)`$ $$Res_{z=0}Tr\psi _ϵ(D_s,y,D_y)[\frac{1}{L}2(\pi LD_s+H_{\alpha ,\lambda ,(\mu ,\nu )})+𝒫_{k+1}]^ke^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}e^{iL𝒫_{k+1}}(LD_s)^z.$$ In view of the microlocal cutoff, the operator under the trace is of trace class for $`Rez`$ sufficiently large. Indeed, in estimating the trace we may eliminate the unitary factors and we are then left with a pseudodifferential operator whose complete symbol is a polynomial in $`(\sigma ,y,\eta )`$ times a factor of $`\sigma ^{Rez}\psi _ϵ(\sigma ,y,\eta ).`$ The integral in the transverse $`(y,\eta )`$ variables is bounded by the volume of the ball $`(y^2+\eta ^2)<\sigma `$ and hence is of order $`\sigma ^n`$. Since a pseudodifferential operator is Hilbert-Schmidt if its Weyl symbol is in $`L^2`$, the operator under the trace is Hilbert-Schmidt for $`Rez>n`$ and in particular is of trace class. Moreover, since it is the non-commutative residue of a Fourier Integral operator, one knows apriori that it admits a meromorphic continuation to $``$ with at most simple poles \[Z.2\]. Hence the residue is well defined. As in \[Z.1\], we view the trace as a function of the parameters $`(\alpha ,\lambda ,\mu ,\nu )`$ and use the explicit form of the exponential in $`e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}`$ to rewrite (5.7) in the form $$Res_{z=0}\underset{n=1}{\overset{\mathrm{}}{}}n^zTr\psi _ϵ(D_s,y,D_y)\{[\frac{1}{L}(2\pi n+\underset{j=1}{\overset{p}{}}\alpha _jD_{\alpha _j}+\underset{j=1}{\overset{q}{}}\lambda _j_{\lambda _j}+$$ $`(5.8)`$ $$+\underset{j=1}{\overset{2c}{}}(\mu _j_{\mu _j}+\nu _jD_{\nu _j})+𝒫_{k+1}(n,D_{\alpha _1},\mathrm{},D_{\nu _{2c}},L)]^ke^{iL𝒫_{k+1}(n,D_{\alpha _1},\mathrm{},D_{\nu _{2c}},L)}e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}\}.$$ Here, we have used that $`D_s`$ commutes with $`(y,D_y)`$ to replace it by its eigenvalue in the $`s`$-trace, and we have repeatedly used identities of the form $`F(D_x)e^{ixP}=F(P)e^{ixP}`$ ($`x`$). Since $`𝒫_{k+1}(n,D_{\alpha _1},\mathrm{},D_{\nu _{2c}},L)`$ is a symbol of order $`1`$ in $`n`$ with coefficients given by polynomials in the operators $`D_{\alpha _j}`$ (etc.), we can expand the kth power in (5.8) as an operator-valued polyhomogeneous function of $`n`$. At least formally, we can also expand the exponential $`e^{iL𝒫_{k+1}(n,D_{\alpha _1},\mathrm{},D_{\nu _{2c}},L)}`$ in a power series and then expand each term in the power series as a polynomial in $`n^1`$. Collecting powers of n, the right side of (5.8) may be written in the form $`(5.9)`$ $$Res_{z=0}\underset{n=1}{\overset{\mathrm{}}{}}\underset{j=o}{\overset{\mathrm{}}{}}n^{z+kj}_{k,kj}(D_{\alpha _1},\mathrm{},D_{\mu _c+i\nu _c},D_{\mu _ci\nu _c})Tr\psi _ϵ(n,y,D_y)e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}},$$ with $`_{k,kj}(D_\alpha ,\mathrm{},D_{\mu _c+i\nu _c},D_{\mu _ci\nu _c}L)`$ the coefficient of $`n^{kj}`$ in (5.8). The expansion of the exponential is justified as in the elliptic case: as in \[Z.1, (4.29)\] we may write $$e^{iL𝒫_{k+1}}:=e_N(iL𝒫_{k+1})+(iL𝒫_{k+1})^{N+1}b_N(iL𝒫_{k+1})$$ with $`e_N(ix)=1+ix+\mathrm{}+\frac{(ix)^N}{N!}`$, with $`𝒫_{k+1}`$ short for $`𝒫_{k+1}(n,D_{\alpha _1},\mathrm{},D_{\nu _{2c}})`$ and with $`b_N(ix)`$ a bounded function. The $`e_N`$ term contributes a finite number of terms of the desired form (5.9). For the remainder, we expand $`(iL𝒫_{k+1})^{N+1}`$ as a polynomial in $`n^1`$ with coefficients given by operators $`Q_{Np}(D_{\alpha _1},\mathrm{},D_{\nu _{2c}})`$ and observe that each term has a factor of $`n^{N1}`$. For each such term, we remove the coefficient operator $`Q_{Np}`$ from the sum $`_n`$, as above, leaving only the factor of $`b_N`$. Since $`b_N(ix)`$ is a bounded function, it follows that $`b_N(iL𝒫_{k+1})`$ is a bounded operator on $`L^2`$; and since each term of the resulting sum has at least the factor $`n^{zN1+k}`$ (possibly multiplied by a further negative power of $`n`$), we see that the remainder is a sum of terms of the form $`(5.10)`$ $$Res_{z=0}Q_{Np}(D_{\alpha _1},\mathrm{},D_{\nu _{2c}})\underset{kn}{}n^{z+kN1l}b_N(i𝒫_{k+1}(n,D_{\alpha _1},\mathrm{}))Tr\psi _ϵ(n,y,D_y)e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}.$$ We then observe that the sum is bounded by $`_{m=1}^{\mathrm{}}m^{RezN1+k+n}`$, hence converges absolutely and uniformly for $`Rez>N+k+n`$. It follows that for $`N>(n+k)`$ the sum in (5.10) defines a holomorphic function of $`z`$ in a half-plane containing $`z=0`$ and since the operations of taking the residue in $`z`$ and derivatives in $`\alpha `$ commute, each term (5.10) is zero. This justifies (5.10) and shows that it is actually a finite sum in j, say $`j<M`$ (in fact M=(k+1)(n+k+1)). The residue in (5.10) is therefore well-defined and independent of $`ϵ`$. Since $`Tr\psi _ϵ(n,y,D_y)e^{iH_{\alpha ,\lambda ,(\mu ,\nu )}}T(\alpha ,\lambda ,\mu ,\nu )`$ in the sense of distributions as $`ϵ\mathrm{}`$ we must have $`(5.11)`$ $$a_{\gamma k}=Res_{z=0}\underset{m=1}{\overset{\mathrm{}}{}}\underset{j=o}{\overset{M}{}}m^{z+kj}_{k,kj}(D_\alpha ,\mathrm{},D_{\mu _c+i\nu _c},D_{\mu _ci\nu _c})T(\alpha ,\lambda ,\mu ,\nu )$$ $$=Res_{z=0}\underset{j=0}{\overset{M}{}}\zeta (z+jk)_{k,kj}(D_\alpha ,\mathrm{},D_{\mu _c+i\nu _c},D_{\mu _ci\nu _c})T(\alpha ,\lambda ,\mu ,\nu ).$$ Here, $`\zeta `$ is the Riemann zeta-function, which has only a simple pole at $`s=1`$ with reside equal to one. It follows that the only term contributing to (5.11) is that with $`j=k+1`$ and hence we have $`(5.12)`$ $$a_{\gamma k}=_{k,1}(D_\alpha ,\mathrm{},D_{\mu _c+i\nu _c},D_{\mu _ci\nu _c})T(\alpha ,\lambda ,\mu ,\nu ).$$ It follows that the wave invariants consist of the geometric data contained in the coefficients of the $`_{k,1}`$’s, and hence in the normal form coefficients. But the algorithm for constructing the normal form is essentially the same as in the elliptic case and so the geometric data entering into the normal form coefficients is of precisely the same kind.∎ ## 6. Inverse Problems: Proofs of Theorem II and Corollary II.1 Our first goal in this section is to prove that the wave invariants of $`\gamma ,\gamma ^2,\mathrm{}`$ determine the quantum normal form coefficients at $`\gamma .`$ Proof of Theorem II: We recall that the quantum normal coefficients are the coefficients of the action monomials in the action polynomials $`\stackrel{~}{p}_j(I_1^e,\mathrm{},I_{2c}^{ch,Im})`$ of (5.2). These coefficients determine, and are determined by, the coefficients of the monomials in the polynomials $`p_\nu (I_1^e,\mathrm{},I_{2c}^{ch,Im})`$ in Theorem B. They also corresponding bi-uniquely to the coefficients of the classical action monomials in the complete symbols of either set of action polynomials. We also observe that the quantum normal form coefficients determine, and are determined by, the coefficients of the constant coefficient partial differential operator (PDO) $`(6.1k)`$ $$_{k,1}(D_\alpha ,D_\lambda ,D_{\mu +i\nu },D_{\mu i\nu }):=\underset{(a,b,c_1,c_2)^n:\|a\|+\|b\|+\|c_1\|+\|c_2\|k+1}{}C_{k;abc_1c_2}D_\alpha ^aD_\lambda ^bD_{\mu +i\nu }^{c_1}D_{\mu i\nu }^{c_2}$$ where $`a^p,b^q,c_1,c_2^c.`$ This can be proved easily by induction on k: In the case k=1, $`_{k,1}`$ is obtained from $`\stackrel{~}{p}_1`$ by substituting the variables $`D_{\alpha _1},`$ etc. in for the variables $`I_1^e,`$ etc. Assuming inductively that we have determined the coefficients of $`\stackrel{~}{p}_1,\mathrm{},\stackrel{~}{p}_k`$ from those of $`_{1,1},\mathrm{},_{1,k}`$, we note that $`\stackrel{~}{p}_{k+1}`$ contributes to the residue (5.7) for the first time at the $`k+1`$st stage. Since it comes with the denominator $`D_s^{k+1}`$, it only contributes to the residue when composed with $`D_s^k`$. From the form of (5.7) it is clear that only one term involving $`\stackrel{~}{p}_{k+1}`$ contributes non-trivially, and that is the one which appears in the linear term in the expansion of the exponential. Hence, its contribution to $`_{k+1,1}`$ is again just the substitution of the variables $`D_{\alpha _1},`$ etc. in for the variables $`I_1^e,`$ etc. We next observe that the quantum normal form of $`\mathrm{\Delta }`$ at any iterate $`\gamma ^N`$ of $`\gamma `$ is the same as for the primitive $`\gamma `$ itself. Hence the PDO $`_{k,1}`$ is independent of the number $`N`$ of iterations. On the other hand, under the iteration $`\gamma \gamma ^N`$, the Poincare map transforms by $`P_{\gamma ^N}P_\gamma ^N`$. Therefore the expression in (5.12) for the kth wave invariant of $`\gamma ^N`$ is given by: $$_{k,1}(D_\alpha ^{},D_\lambda ^{},D_{\mu ^{}+i\nu ^{}},D_{\mu ^{}i\nu ^{}})$$ $`(6.2k,N)`$ $$\mathrm{\Pi }_{j=1}^p\frac{e^{\frac{1}{2}i\alpha _j^{}}}{(1e^{i\alpha _j^{}})}\mathrm{\Pi }_{j=1}^q\frac{e^{\frac{1}{2}\lambda _j^{}}}{(1e^{\lambda _j^{}})}\mathrm{\Pi }_{j=1}^c\frac{e^{\frac{1}{2}(\mu _j^{}+i\nu _j^{})}}{(1e^{(\mu _j^{}+i\nu _j^{})})}\frac{e^{\frac{1}{2}(\mu _j^{}i\nu _j^{})}}{(1e^{(\mu _j^{}i\nu _j^{})})}\|_{(\alpha ^{},\lambda ^{},\mu ^{},\nu ^{})=N(\alpha ,\lambda ,\mu ,\nu )}.$$ It therefore suffices to prove that for all k the coefficients of the PDO $`_{k,1}`$ can be determined from its values (6.2 k, N) on $`T(N\alpha ,N\lambda ,N\mu ,N\nu )`$ for $`N=\pm 1,\pm 2,\mathrm{}.`$ We begin the proof by noting that (6.2 k, N) can be rewritten as: $`(6.3k,N)`$ $$\mathrm{\Pi }_{j=1}^pe^{\frac{1}{2}iN\alpha _j}\mathrm{\Pi }_{j=1}^qe^{\frac{1}{2}N\lambda _j}\mathrm{\Pi }_{j=1}^ce^{\frac{1}{2}N(\mu _j+i\nu _j)}e^{\frac{1}{2}N(\mu _ji\nu _j)}$$ $$_{k,1}(D_\alpha ^{}+\frac{1}{2},D_\lambda ^{}+\frac{1}{2},D_{\mu ^{}+i\nu ^{}}+\frac{1}{2},D_{\mu ^{}i\nu ^{}}+\frac{1}{2})[\mathrm{\Pi }_{j=1}^p\frac{1}{(1e^{i\alpha _j^{}})}$$ $$\mathrm{\Pi }_{j=1}^q\frac{1}{(1e^{\lambda _j^{}})}\mathrm{\Pi }_{j=1}^c\frac{1}{(1e^{(\mu _j^{}+i\nu _j^{})})}\frac{1}{(1e^{(\mu _j^{}i\nu _j^{})})}]\|_{(\alpha ^{},\lambda ^{},\mu ^{},\nu ^{})=N(\alpha ,\lambda ,\mu ,\nu )}.$$ Making the substitutions $`D_\alpha D_\alpha +\frac{1}{2}`$ (etc.) in (6.1 k) we obtain a new PDO whose coefficients $`C_{k;abc_1c_2}^{}`$ correspond in a bi-unique way with the original $`C_{k;abc_1c_2}`$’s. Hence it will suffice to show that we can determine the $`C_{k;abc_1c_2}`$’s from the values (6.3 k,N). To do so, we will regard (6.3 k, N) as the values at integral points $`z=N`$ of a function of $`z`$. From the fact that $$D_\alpha (1e^{i\alpha })^1=[(1e^{i\alpha })^2(1e^{i\alpha })^1]$$ we see that this function is a polynomial in $`(1e^{iz\alpha _j})^1,(1e^{z\lambda _j})^1,(1e^{z(\mu +i\nu }))^1,(1e^{z(\mu i\nu }))^1.`$ We clear the denominators to obtain the entire function $`(6.4k,z)`$ $$\underset{(a,b,c_1,c_2)^n:\|a\|+\|b\|+\|c_1\|+\|c_2\|k+1}{}C_{k;abc_1c_2}^{}[\mathrm{\Pi }_{j=1}^p(1e^{iz\alpha _j})\mathrm{\Pi }_{j=1}^q$$ $$(1e^{z\lambda _j})\mathrm{\Pi }_{j=1}^c(1e^{z(\mu _j+i\nu _j)})(1e^{z(\mu _ji\nu _j)})]^{(k+1)}\mathrm{\Pi }_{j=1}^pe^{\frac{1}{2}iz\alpha _j}\mathrm{\Pi }_{j=1}^qe^{\frac{1}{2}z\lambda _j}\mathrm{\Pi }_{j=1}^ce^{\frac{1}{2}z(\mu _j+i\nu _j)}e^{\frac{1}{2}z(\mu _ji\nu _j)}$$ $$D_\alpha ^{}^aD_\lambda ^{}^bD_{\mu ^{}+i\nu ^{}}^{c_1}D_{\mu ^{}i\nu ^{}}^{c_2}[\mathrm{\Pi }_{j=1}^p\frac{1}{(1e^{i\alpha _j^{}})}\mathrm{\Pi }_{j=1}^q\frac{1}{(1e^{\lambda _j^{}})}\mathrm{\Pi }_{j=1}^c\frac{1}{(1e^{(\mu _j^{}+i\nu _j^{})})}\frac{1}{(1e^{(\mu _j^{}i\nu _j^{})})}]\|_{(\alpha ^{},\lambda ^{},\mu ^{},\nu ^{})=z(\alpha ,\lambda ,\mu ,\nu )}$$ which is an exponential polynomial of the form $`(6.5k)`$ $$\underset{\beta ^n}{}c_{k;\beta }e^{z\beta +(\frac{1}{2},\mathrm{},\frac{1}{2}),(\alpha ,\lambda \mu ,\nu )}.$$ (6.6) Lemma 1 The coefficients $`C_{k;abc_1c_2}`$ can be determined from the coefficients $`c_{k;\beta }`$ in (6.5 k). Proof: We first show that the coefficients $`C_{k;abc_1c_2}`$ with $`\|a\|+\|b\|+\|c_1\|+\|c_2\|=k+1`$ can be determined from the $`c_{k;\beta }`$’s. Indeed, before multiplying by $$[\mathrm{\Pi }_{j=1}^p(1e^{iz\alpha _j})\mathrm{\Pi }_{j=1}^q(1e^{z\lambda _j})\mathrm{\Pi }_{j=1}^c(1e^{z(\mu _j+i\nu _j)})(1e^{z(\mu _ji\nu _j)})]^{(k+1)}$$ $$\mathrm{\Pi }_{j=1}^pe^{\frac{1}{2}iz\alpha _j}\mathrm{\Pi }_{j=1}^qe^{\frac{1}{2}z\lambda _j}\mathrm{\Pi }_{j=1}^ce^{\frac{1}{2}z(\mu _j+i\nu _j)}e^{\frac{1}{2}z(\mu _ji\nu _j)}$$ $`C_{k;abc_1c_2}^{}`$ is uniquely determined as the coefficient of the monomial $$\mathrm{\Pi }_{j=1}^p(1e^{iz\alpha _j})^{(a_j+1)}\mathrm{\Pi }_{j=1}^q(1e^{z\lambda _j})^{(b_j+1)}\mathrm{\Pi }_{j=1}(1e^{z(\mu _j+i\nu _j)})^{(c_{1j}+1)}(1e^{z(\mu _ji\nu _j)})^{(c_{2j}+1)}.$$ Since we are multiplying thru by a quantity independent of $`a,b,c_1,c_2`$, it follows that $`C_{k;abc_1c_2}^{}`$ is uniquely determined as the coefficient the of the monomial $$\mathrm{\Pi }_{j=1}^p(1e^{iz\alpha _j})^{(k+1)(a_j+1)}e^{\frac{1}{2}iz\alpha _j}\mathrm{\Pi }_{j=1}^q(1e^{z\lambda _j})^{(k+1)(b_j+1)}e^{\frac{1}{2}z\lambda _j}$$ $$\mathrm{\Pi }_{j=1}^c(1e^{z(\mu _j+i\nu _j)})^{(k+1)(c_{1j}+1)}(1e^{z(\mu _ji\nu _j)})^{(k+1)(c_{2j}+1)}e^{\frac{1}{2}z(\mu _j+i\nu _j)}e^{\frac{1}{2}z(\mu _ji\nu _j)}.$$ Expanding into an exponential polynomial, we find that $`C_{k;abc_1c_2}^{}`$ is uniquely determined as the coefficient of the monomial $$\mathrm{\Pi }_{j=1}^pe^{iz((k+1)(a_j+1)+\frac{1}{2})\alpha _j}\mathrm{\Pi }_{j=1}^qe^{z((k+1)(b_j+1)+\frac{1}{2})\lambda _j}$$ $$\mathrm{\Pi }_{j=1}^ce^{z((k+1)(c_{1j}+1)+\frac{1}{2})(\mu _j+i\nu _j)}e^{z((k+1)(c_{2j})+\frac{1}{2})(\mu _ji\nu _j)}).$$ Uniqueness follows from the fact that the vector $`\beta +\frac{1}{2}`$ is a minimal element of the set of exponent vectors occuring in (6.5 k). Since $`C_{k;abc_1c_2}^{}=C_{k;abc_1c_2}`$ when $`\|a\|+\|b\|+\|c_1\|+\|c_2\|=k+1`$, we have determined $`C_{k;abc_1c_2}.`$ We then remove the $`C_{k;abc_1c_2}(D_\alpha +\frac{1}{2})^a(D_\lambda +\frac{1}{2})^b(D_{\mu +i\nu }+\frac{1}{2})^{c_1}(D_{\mu i\nu }+\frac{1}{2})^{c_2}`$ terms with $`\|a\|+\|b\|+\|c_1\|+\|c_2\|=k+1`$ in (6.5 k). This leaves only terms with coefficients $`C_{k;abc_1c_2}`$ with $`\|a\|+\|b\|+\|c_1\|+\|c_2\|k.`$ Hence we can continue the process of recovering coefficients until the end. ∎ Let us now rewrite $`(6.6a)`$ $$\underset{\beta ^n}{}c_{k;\beta }e^{z\beta +(\frac{1}{2},\mathrm{},\frac{1}{2}),(\alpha ,\lambda ,\mu ,\nu )}$$ in the form $`(6.6b)`$ $$\underset{j=1}{\overset{M}{}}a_{jk}e^{z\omega _j}.$$ (6.7) Lemma The complex exponents $`\omega _j`$ in (6.6b), together with $`\pi `$, are independent over the rationals. Moreover, the coefficients $`c_{k:\beta }`$ can be determined from the coefficients $`a_{jk}`$ Proof: The $`\omega _j`$’s are rational linear combinations of the exponents $`\alpha _j,\lambda _j,\mu _j,\nu _j`$, which by assumption are independent, with $`\pi `$, over the rationals. This independence also implies that the exponents $`\beta +(\frac{1}{2},\mathrm{},\frac{1}{2}),(\alpha ,\lambda ,\mu ,\nu )`$ are all distinct. Hence the coefficients in (6.6a)-(6.6b) are the same. ∎ The proof of Theorem II is thus reduced to the following general statement about exponential polynomials. (6.8) Lemma Suppose that the exponents $`\omega _j`$ of an exponential polynomial (6.6b) are independent (with $`\pi `$) over the rationals. Then the coefficients $`a_{jk}`$ of can be determined from the values of this polynomial at $`z=N.`$ Proof: If not, there would exist a polynomial with the given complex frequencies which vanished at all integers $`z=N`$. But the different terms $`e^{z\omega _j},e^{z\omega _k}`$ have different exponential growth rates along $`z=N`$ or $`z=N`$ ($`N`$) unless $`Re\omega _j=Re\omega _k.`$ Let us write the large sum as a sum of smaller sums with a common $`Re\omega .`$ Each of the smaller sums must separately vanish for $`z.`$ Multiply each one by the relevant factor of $`e^{Re\omega }`$. Each then turns into an exponential polynomial with imaginary exponents, which vanishes for all $`z=N.`$ Since the frequencies are independent(with $`\pi `$) over $`𝐐`$, each of these polynomials must vanish identically if it vanishes at integral points. Hence the coefficients $`a_{jk}`$ are uniquely determined by values of the large sum at integral points. ∎
warning/0002/gr-qc0002006.html	ar5iv	text	# Active controls in interferometric detectors of gravitational waves: inertial damping of the VIRGO superattenuator11footnote 1Lecture given at the International Summer School on Experimental Physics of Gravitational Waves - Urbino (Italy), September 6-18, 1999 ## 1 Introduction Operating an interferometer for gravitational waves detection requires the implementation of many active controls on the different parts of the apparatus. For instance, feedback controls are needed to reduce the laser frequency and power fluctuations, to keep the optical cavities in resonance and to maintain the interferometer output on a dark fringe. The basic idea is that the apparatus works at its best strictly around a well defined working point (laser fluctuations and shot noise at a minimum, optical cavities in resonance and interferometer output on a dark fringe). Internal and external disturbances overload the dynamic range of the interferometer. Therefore, the system has to be forced to remain in the correct working position. This is achieved via feedback controls. In this lecture we examine in detail one of the controls needed to reach the final goal of operating the interferometer: the inertial damping of the VIRGO superattenuator (SA) . The SA can be described as a chain of mechanical “filters”, each one acting as a “spring” in 6 degrees of freedom. The normal mode frequencies of the SA range between 0.04 and about 2 Hz. At frequencies $`f>>2`$ Hz the SA acts as a steep filter of the seismic vibrations of the ground (an attenuation factor of $`10^{15}`$@10 Hz is expected). Therefore, in the interferometer detection band (10 Hz-few kHz), the suspended mirror is “disconnected” from the ground. Beside being a low pass filter for ground vibrations the SA is a tool for actively controlling the mirror position. Forces that move and steer the mirror can be exerted in 3 points of the SA chain: the inverted pendulum (IP) , the marionetta (a special mechanical tool designed to steer the mirror) and the mirror itself (from a suspended reference mass). Keeping the interferometer in the operating position requires the mirror to have a maximum RMS relative motion of less the $`10^{12}`$ m. In the frequency range where the SA is fully effective ($`f>>4`$ Hz) the residual motion of the mirror is negligible. On the other hand, the mirror free motion in the region of the SA normal modes is $`100`$ $`\mu `$m . Feedback forces acting on the SA must reduce the mirror motion from $`100`$ $`\mu `$m to $`10^{12}`$ m. The dynamic range of a feedback system able to perform this control has to be huge: the control is performed in three steps (hierarchical control ). The first step is a damping of the SA normal modes in order to reduce the mirror residual motion to less than 10 $`\mu `$m . This is necessary in order to control the mirror position acting on the lower stages without reinjecting noise into the detection band. We describe in the following an implementation of a high gain and wideband damping of the SA resonances. ## 2 Experimental setup The setup (fig. 1) of the experiment is composed of a full scale superattenuator, provided with 3 accelerometers (placed on the top of the IP), 3 LVDT position sensors (measuring the relative motion of the IP with respect to an external frame), 3 coil-magnet actuators. The accelerometers work in the range DC-400 Hz and have acceleration spectral sensitivity $`10^9\mathrm{m}\mathrm{s}^2\mathrm{Hz}^{1/2}`$ below 3 Hz . The sensors and actuators are all placed in pin-wheel configuration. The sensor and actuator signals are computer controlled by a ADC (16 bit)-DSP-DAC (20 bit) system. The DSP handles the signals of all the sensors and actuators. It can combine them by means of matrices, create complex feedback filters (like the one of fig. 8) with high precision poles/zeroes placement and perform all the calculations at a high sampling rate (10 kHz). The suspended mirror is also provided with LVDT position sensors to measure its displacement with respect to ground. ## 3 The approach to the problem of control The IP, where all the sensors are placed and on which the forces are exerted, has 3 main resonant modes: two translation ($`X,Y`$) and the rotation around the vertical axis ($`\mathrm{\Theta }`$). Each sensor is sensitive to all the modes and each actuator can excite all the modes. In control theory language, such a system is defined MIMO (multiple in-multiple out). Controlling a MIMO system can be very difficult. Our approach has been different: the signals of the 3 sensors are digitally mixed (using proper transformation matrices) to build up virtual sensors, sensitive to one mode only and “blind” to the others. At the same time, we build up virtual actuators able to excite each mode separately. The system is thus uncoupled into 3 SISO (single in-single out) subsystems. In terms of analytical mechanics this means the system is described in the normal modes basis. The equations of motions take the form: $$\ddot{x}_k+\omega _k^2x_k=q_k,k=1\mathrm{}3$$ (1) where $`x_k`$ is the $`k`$th normal mode, $`\omega _k/2\pi `$ is the corresponding resonant frequency and $`q_k`$ the generalized force on that mode. Let $`𝐮=(u_1,\mathrm{},u_3)`$ the vector made by the outputs of the 3 sensors and let $`𝐱=(x_1,\mathrm{},x_3)`$ the vector made by the 3 virtual sensors. Analogously, let $`𝐯`$ the vector made by the 3 currents driving the actuators and $`𝐪`$ the vector of the 3 generalized forces. The transformation from the system of the physical sensors/actuators to that of the virtual ones is operated by two matrices, such that: $`𝐯`$ $`=`$ $`\mathrm{𝐃𝐪}`$ (2) $`𝐱`$ $`=`$ $`\mathrm{𝐒𝐮}`$ (3) The sensing ($`𝐒`$) and driving ($`𝐃`$) matrices are experimentally measured. The measurement procedure is described in details in ref. 9. The result of the measurement is determined by the mechanical characteristics of the system (resonant frequencies, quality factors), the geometry of the sensors and actuators and their calibration, but it is not needed to know them in order to measure the matrices. The logic of diagonalisation is explained by fig. 2. After the diagonalisation the system can be considered as composed of 3 uncoupled oscillators, and the control strategy for each of them has to be defined independently (see fig. 3). ## 4 Inertial damping: principle The control we describe here is called inertial damping because it is performed by using (mostly) inertial sensors (accelerometers). In the following , with the help of a simple model, we explain why this is the best choice to achieve a high performance damping. Let us consider a simple pendulum of mass $`m`$ and length $`l`$. Let $`x`$ be the abscissa of the suspended mass, $`x_0`$ that of the suspension point. Let $`F_{\mathrm{fb}}`$ the external force on the pendulum (i.e. the feedback force to control it). The equation of motion is then: $$F_{\mathrm{fb}}=m\ddot{x}+\gamma \dot{x}+k(xx_0)$$ (4) where $`\gamma `$ is the viscous dissipation factor and $`k=mg/l`$. The control loop of such a system is sketched in fig. 4, where $`H(s)`$ is the mechanical transfer function, $`G(s)`$ is the compensator and out is the output of the sensor used. The goal of the control is to damp the pendulum resonance. This can be done easily with a viscous (theoretical) feedback force: $$F_{\mathrm{fb}}=\gamma ^{}\dot{x}$$ (5) Our sensors do not measure $`x`$. Their output is: $$out=\{\begin{array}{cc}xx_0\hfill & \text{for displacement sensors}\hfill \\ \ddot{x}\hfill & \text{for accelerometers}\hfill \end{array}$$ (6) Therefore, the actual “viscous” force that can be built if position sensors are used has the form: $$F_{\mathrm{fb}}^p=\gamma ^{}\frac{\mathrm{d}}{\mathrm{d}t}(xx_0)$$ (7) It can be easily shown that with such a feedback force the closed loop equation of motion (in Laplace space) reduces to: $$x(s)=\frac{\omega _0^2+G_0s}{s^2+\omega _0^2+(\omega _0/Q+G_0)s}x_0(s)$$ (8) where $`G_0=\gamma ^{}/m`$ is a gain parameter<sup>2</sup><sup>2</sup>2In a real feedback system a frequency dependent gain function $`G(s)`$ rather than a gain parameter has to be considered. measuring the intensity of the viscous feedback force, and $`Q`$ is the open loop quality factor. When the loop is closed a damping of the resonance is achieved: $$Q^{}\stackrel{G>>1}{}\frac{\omega _0}{G}$$ (9) Nevertheless, as the gain is increased, a larger amount of noise is reinjected off-resonance. This is associated to the term “$`G_0s`$” in the numerator of (8) and depends on the fact that the sensor used to build up the feedback force measures the position of the pendulum with respect to ground. Therefore, an infinitely efficient feedback would “freeze” the pendulum to ground (which is seismic noisy), reducing its motion at the resonance, with the drawback of bypassing its attenuation properties above resonance. The situation is fairly different when an inertial sensor is used. In this case the viscous feedback force is obtained by integrating the accelerometer output, and the output does not depend on $`x_0`$: $$F_{\mathrm{fb}}^a=\gamma ^{}\ddot{x}dt$$ (10) The closed loop equation of motion is then: $$x(s)=\frac{\omega _0^2}{s^2+\omega _0^2+(\omega _0/Q+G_0)s}x_0(s)$$ (11) A damping of the resonance is obtained (exactly as in the previous case) but without reinjection of off-resonance noise. In fig. 5 a simulation of the closed loop transfer function $`x(s)/x_0(s)`$ is shown in the two cases. Up to now we have considered a simple viscous damping. It is possible to increase the bandwidth of the control if the feedback force contains a term proportional to $`x`$ (the double integral of the accelerometer signal). The result obtained in this case is shown in fig. 7. ## 5 Control strategy In this section we extend the principles of the previous section and describe the strategy to control the SA. The basic idea of inertial damping is to use the accelerometer signal to build up the feedback force. As a matter of fact, an infinitely efficient feedback using only the inertial sensor information, would null the acceleration of the pendulum, but it would not do anything if the pendulum moves at constant velocity. Such a control would be unstable with respect to drifts. Therefore, if the control band is to be extended down to DC, a position signal is necessary. Our solution was a merging of the two sensors: the virtual LVDT (position) and accelerometer signals are combined in such a way that the LVDT signal ($`l(s)`$) dominates below a chosen cross frequency $`f_{\mathrm{merge}}`$ while the accelerometer signal ($`a(s)`$) dominates above it (see fig. 7 and ref. 7). The feedback force has the form<sup>3</sup><sup>3</sup>3Actually, the LVDT signal $`l(s)`$ is properly filtered in order to preserve feedback stability at the crossover frequency and in order to reduce the amount of reinjected noise at $`f>f_{\mathrm{merge}}`$.: $$F_{\mathrm{fb}}=G(s)\left[a(s)+ϵl(s)\right]$$ (12) where $`G(s)`$ is the digital filter transfer function (see fig. 8) and $`ϵ`$ is the parameter whose value determines $`f_{\mathrm{merge}}`$. We have chosen $`f_{\mathrm{merge}}10`$ mHz (corresponding to $`ϵ510^3`$). This approach stabilizes the system with respect to low frequency drifts at the cost of reinjecting a fraction $`ϵ`$ of the seismic noise via the feedback. We describe in the following the feedback design for the 3 d.o.f., starting from the the translational ones. The virtual $`X`$ and $`Y`$ sensors show many resonant peaks (the modes of a chain of pendulums) and this requires a more sophisticated feedback strategy. The digital filter used to control the translation modes ($`G(s)`$) is shown in fig. 8 (LEFT). It shows three main features: * for $`0.01<f<2`$ Hz the gain is proportional to $`f^2`$. This corresponds to the case of fig. 7: the accelerometer signal is integrated twice and the feedback force is proportional to $`x`$; * for $`f>2`$ Hz the gain is proportional to $`f^1`$. The accelerometer signal is integrated once: the feedback force is proportional to the velocity and a viscous damping is achieved; * the peaks visible in the filter are necessary to compensate the corresponding dips in the mechanical transfer function ($`H(s)`$) of fig. 3, in order to make the feedback stable. Fig. 8 (RIGHT) shows the open loop gain transfer function $`G(s)H(s)`$. The damping strategy for the $`\mathrm{\Theta }`$ mode is simpler: the $`\mathrm{\Theta }`$ virtual sensor (fig. 3, RIGHT) shows one resonance peak only and no dips: no compensation is necessary. Apart from this, the feedback strategy is similar to the ones used for the translational modes. ## 6 Inertial damping: experimental results The result of the inertial control (on 3 d.o.f.) is shown in figure 9. The measurement has been performed in air. The noise on the top of the IP is reduced over a wide band (10 mHz \- 4 Hz). A gain $`>1000`$ is obtained at the main SA resonance (0.3 Hz). The RMS translational motion of the IP (calculated as $`x_{\mathrm{RMS}}(f)=\sqrt{_f^{\mathrm{}}\stackrel{~}{x}^2(\nu )d\nu }`$) in 10 sec. is reduced from more than 30 to 0.3 $`\mu `$m. The closed loop floor noise corresponds to the fraction of seismic noise reinjected by using the position sensors for the DC control and can, in principle, be reduced by a steeper low pass filtering of the LVDT signal at $`f>f_{\mathrm{merge}}`$ and by lowering $`f_{\mathrm{merge}}`$: both this solution have drawbacks and need a careful implementation. Preliminary measurements of the displacement of the mirror with respect to ground have been performed in air, using an LVDT position sensor. The residual RMS mirror motion in 10 sec. is<sup>4</sup><sup>4</sup>4This number has been obtained with a simpler feedback design, less aggressive than the one of fig. 8: the gain raised as $`1/f`$ (pure viscous damping force), the crossover frequency was 30 mHz and no compensation of the dips was needed.: $$x_{\mathrm{RMS}}(0.1\mathrm{Hz})3\mu \mathrm{m}.$$ (13) When the damping is on such a measurement can provide only an upper bound because the LVDT output is dominated by the seismic motion of the ground. ## 7 Conclusions In this lecture we have tried to outline how to face the problem of reducing the free motion of the suspended optical components of the VIRGO interferometer, associated to the resonances of the suspension. This is only the beginning: once the motion of the mirrors is reduced to a few microns, the lower control stages can operate to lock the interferometer in the correct operation state. ## Acknowledgments The author wishes to thank all the people of the Pisa and Florence VIRGO Groups. Among them, special thanks to Diego Passuello, Alberto Gennai and Andrea Marin.
warning/0002/hep-th0002002.html	ar5iv	text	# Untitled Document Brown Het-1211 Duality and Combinatorics of Long Strings in ADS3 Mihail Mihailescu and Sanjaye Ramgoolam Brown University Providence, RI 02912 mm,ramgosk@het.brown.edu The counting of long strings in ADS3, in the context of Type IIB string theory on $`ADS_3\times S^3\times T^4`$, is used to exhibit the action of the duality group $`O(5,5;Z)`$, and in particular its Weyl Subgroup $`S_5Z_2`$, in the non-perturbative phenomena associated with continuous spectra of states in these backgrounds. The counting functions are related to states in Fock spaces. The symmetry groups also appear in the structure of compactifications of instanton moduli spaces on $`T^4`$. 12/99 1. Introduction We study counting problems related to the phenomena of long strings in ADS3, in the context of type IIB string theory on $`ADS_3\times S^3\times T^4`$, which is one of the examples which enters the Maldacena conjecture . This example is of special interest because the CFT dual is a tractable 2D CFT based on an orbifold $`S^N(T^4)`$. Various aspects of the operator algebra and correlation functions have been studied for example in . An important issue that has to be understood better in order to make the orbifold CFT more useful is the precise map between the moduli space of the CFT and that of the spacetime theory. Significant steps in this direction have been made in . Closely related to the issue of moduli is the issue of dualities, since the the duality group $`O(5,5;Z)`$ appears in the description of the moduli space $`O(5,5;Z)\backslash O(5,5;R)/O(5)\times O(5)`$ of $`R^6\times T^4`$ compactifications of type IIB. The ADS3 background is obtained by choosing a string in six dimensions and going to the near-horizon limit. The string can be a bound state of one or more of the following : D-string, D5-brane wrapped on the $`T^4`$, D3-branes wrapped along the two cycles of the torus, NS one-brane ( elementary string) or NS 5-brane wrapped on $`T^4`$. The allowed charges of the string live are vectors in a lattice $`\mathrm{\Gamma }^{5,5}`$ ( described more explicitly in section 2 ). After choosing a vector in $`\mathrm{\Gamma }^{(5,5)}`$ describing the background string it is often useful to focus on subgroup $`O(4,5;Z)`$ of $`O(5,5:Z)`$, see for example . For the kinds of questions we will be asking it will be interesting to look at subgroups of $`O(5,5:Z)`$ which may or may not belong to $`O(4,5;Z)`$. A very interesting class of phenomena in ADS3, not directly accessible from the orbifold model at its free point, was studied in . This involves, in the simplest case, an ADS3 background obtained from the near horizon limit of $`Q_1`$ D1 branes and $`Q_5`$ D5-branes. This system allows a D1 ( or D5 ) brane to split off and grow to infinity at finite cost in energy. A semiclassical calculation (valid for large $`Q_1,Q_5`$ ) of the split string worldvolume action shows that the string worldvolume theory is a Liouville theory, which is known to have a continuous spectrum above a threshold. The threshold is at $`Q_5/4`$ if we have a split D1 string, and $`Q_1/4`$ if we have a split D5-string, The spectrum of the Hamiltonian in ADS3 therefore has a continuum starting at $`Q/4`$ where $`Q=min(Q_1,Q_5)`$. This formula has the very interesting symmetry under exchange of $`Q_1`$ and $`Q_5`$, which is in $`O(5,5;Z)`$ but not in $`O(4,5;Z)`$. These long string phenomena correspond to the splitting $`(Q_1,Q_5)=(Q_11,Q_5)+(1,0)`$ or $`(Q_1,Q_5)=(Q_1,Q_51)+(0,1)`$. The conditions for such splittings $`Q=q^{(i)}`$ to be BPS were described in in terms of the geometry of lattices and 5-planes and $`R^{(5,5)}`$ ( we review this in section 2 ). For generic moduli there are no BPS splittings. An interesting class of splittings happens when the NS sector B-field moduli and the RR sector C-field moduli are set to zero ( For the bulk of this paper we will work with the $`B=C=0`$ condition, with off-diagonal components of the metric set to zero, and we get back to non-zero $`B,C`$ in section 7 ). The same symmetry under exchange of $`Q_1`$ and $`Q_5`$ can be seen in the simple exercise of counting of the number of distinct splittings of the $`(Q_1,Q_5)`$ system. The counting is not symmetric under exchange of $`(Q_1,Q_5)`$ with $`(Q_1Q_5,1)`$. In this paper we study generalizations of the counting of BPS splittings, and we find that the Weyl group of $`O(5,5;Z)`$, which is generated by the symmetric group $`S_5`$ and a $`Z_2`$, indicated by writing the Weyl group as $`S_5Z_2`$, has an interesting action on the counting functions associated with BPS splittings. It is natural to expect that it is also a symmetry of the more detailed dynamics of continuous spectra associated with such splittings. Note that we are discussing splittings for systems where the charges have no common factor. As mentioned in when the charge vector is non-primitive there are BPS splittings for arbitrary moduli. Section 2 is a review of relevant background. Sections 3-5 deal with splittings of different kinds of charges. Section 6 discusses the connections of these BPS brane separation problems and the associated symmetries with compactifications of instanton moduli spaces. Section 7 discusses some aspects of splittings beyond $`B=C=0`$. Section 8 discusses the description from the gauge theory point of view of splittings involving NS charges from D-brane systems. 2. Review 2.1. Backgrounds as 5-planes Choosing a background for IIB theory on $`R^6\times T^4`$ requires choosing a 5-plane in $`R^{(5,5)}`$, modulo discrete identifications ( for related discussions see ). The 5-plane is spanned by vectors $`E^{(0)},E^{(1)},E^{(2)},E^{(3)},E^{(4)}`$. We use the notation of ( with a minor reshuffling of entries ) for the vectors spanning the positive 5-plane $`\mathrm{\Theta }`$ in $`R^{(5,5)}`$. $$\begin{array}{cc}& E^\mu =(v^\mu ;C.v^\mu ,0)R^{4,4}\times R^{(1,1)}\hfill \\ & E^4=(0;\beta ,1)R^{4,4}\times R^{(1,1)}\hfill \\ & v^i=(B.\omega ^i,0;\omega ^i)R^{1,1}\times R^{(3,3)}\hfill \\ & v^0=(\alpha ,1;0)R^{1,1}\times R^{(3,3)}\hfill \end{array}$$ Here $`\mu `$ runs from $`0`$ to $`4`$ and $`i`$ runs from $`0`$ to $`3`$. We also have $`\beta =\frac{1}{g_6^2}\frac{1}{2}C.C`$, and $`\alpha =V\frac{1}{2}B.B`$. The charges of strings live in a lattice $`\mathrm{\Gamma }^{(5,5)}R^{(5,5)}`$. The charge $`Q`$ has components $`(Q_1,Q_5;Q_{12},Q_{34},Q_{13},Q_{42},Q_{14},Q_{23};N_1,N_5)`$. Here $`(Q_1,Q_5)`$ are the charges from D1 strings and wrapped D5 strings, and $`(N_1,N_5)`$ are the charges from NS strings and wrapped NS 5-branes. The $`Q_{ij}`$ are the charges of strings obtained by wrapping D3 branes on the $`(ij)`$ cycle. The bilinear form on $`\mathrm{\Gamma }^{(5,5)}`$ evaluated on two vectors $`Q^{(1)}=(Q_1^{(1)},Q_5^{(1)};Q_{12}^{(1)},Q_{34}^{(1)},Q_{13}^{(1)},Q_{42}^{(1)},Q_{14}^{(1)},Q_{23}^{(1)};N_1^{(1)},N_5^{(1)})`$ and $`Q^{(2)}=(Q_1^{(2)},Q_5^{(2)};Q_{12}^{(2)},Q_{34}^{(2)},Q_{13}^{(2)},Q_{42}^{(2)},Q_{14}^{(2)},Q_{23}^{(2)};N_1^{(2)},N_5^{(2)})`$, is $$\begin{array}{cc}& (Q^{(1)},Q^{(2)})=Q_1^{(1)}Q_5^{(2)}+Q_1^{(2)}Q_5^{(1)}+Q_{12}^{(2)}Q_{34}^{(1)}+Q_{12}^{(1)}Q_{34}^{(2)}\hfill \\ & +Q_{13}^{(1)}Q_{42}^{(2)}+Q_{13}^{(2)}Q_{42}^{(1)}+Q_{14}^{(1)}Q_{23}^{(2)}+Q_{14}^{(2)}Q_{23}^{(1)}+N_1^{(1)}N_5^{(2)}+N_1^{(2)}N_5^{(1)}\hfill \end{array}$$ For a rectangular torus ( i.e off-diagonal components of the metric set to zero ) and with vanishing $`B`$ and $`C`$ fields, we have $$\begin{array}{cc}\hfill E^0& =(V,1;0;0,0)\hfill \\ \hfill E^i& =(0,0;\omega ^i;0,0)\hfill \\ \hfill E^4& =(0,0;0;\frac{1}{g_6^2},1)\hfill \end{array}$$ We can choose : $$\begin{array}{cc}& \omega ^1=R_1R_2dx_1dx_2+R_3R_4dx_3dx_4\hfill \\ & \omega ^2=R_1R_3dx_1dx_3R_2R_4dx_2dx_4\hfill \\ & \omega ^3=R_1R_4dx_1dx_4+R_2R_3dx_2dx_3\hfill \end{array}$$ $`R_i`$ are the circumferences of the circles. $`x_i`$ are variables with periodicity $`1`$. With these formulae the tension $`T(Q)`$ ( in units of $`1/g_6`$, where $`g_6`$ is the six-dimensional string coupling ), is given by $`T(Q)=Q_+`$, where $`Q_+`$ is the projection of $`Q`$ to the positive 5-plane. With these expressions, we can recover the mass formulae of . Since physical quantities are given by projections of vectors in $`\mathrm{\Gamma }^{(5,5)}`$ to the 5-plane, the space of physically inequivalent vacua is obtained by modding out by the symmetries of $`\mathrm{\Gamma }^{(5,5)}`$ : $`O(5,5;Z)\backslash O(5,5;R)/O(5,R)\times O(5,R)`$. When we choose a charge for the string in 6 dimensions living in $`\mathrm{\Gamma }^{(5,5)}`$, having the property $`Q^2>0`$, the attractor equations imply that the near horizon moduli satisfy the condition that $`Q`$ is parallel to the 5-plane . 2.2. BPS splitting into multiple parts The condition for a splitting $`Q=q^{(1)}+q^{(2)}`$ to be BPS can be expressed by saying that the projection of $`q^{(1)}`$ or $`q^{(2)}`$ to $`\mathrm{\Theta }`$ is proportional to $`Q`$ . The projection has to be non-negative i.e $`q.Q0`$. The density of states in the continuum associated with the splittings into multiple summands will be larger than in the case of two summands. The condition for these splittings to be BPS can be obtained as in . Consider a vector $`Q\mathrm{\Gamma }^{(5,5)}`$ of charges which can be written as $`Q=q^{(1)}+q^{(2)}+\mathrm{}q^{(l)}`$, satisfying the condition that the projections to the positive 5-plane obey $`\|Q_+\|=\|q_+^{(1)}\|+\|q_+^{(2)}\|+\mathrm{}\|q_+^{(l)}\|.`$ The near horizon geometry satisfies $`Q_+=Q`$. These conditions can be satisfied if each vector $`q^{(i)}`$ in the decomposition has a projection to the 5-plane, $`q_+^{(i)}`$ which is parallel to $`Q`$. If we start with a system of charges $`(Q_1,Q_5)`$ and study its splittings at $`B=C=0`$ with generic radii, we get a set of splittings sitting in a $`\mathrm{\Gamma }^{1,1}`$ lattice. The symmetry of this lattice includes exchanging $`(Q_1,Q_5)(Q_5,Q_1)`$, but does not include $`(Q_1,Q_5)(Q_1Q_5,1)`$. The duality group $`O(\mathrm{\Gamma }^{5,5})`$ does allow us to map the system $`(Q_1,Q_5)`$ to both $`(Q_5,Q_1)`$ and to $`(Q_1Q_5,1)`$. The first kind of map allows us to start from $`B=C=0`$ and end with $`B=C=0`$. The second does not . There is also S-duality symmetry which is often discussed and exploited. It allows us to relate the physics of backgrounds containing NS-NS fields, to the physics of backgrounds with RR charge. It preserves the $`B=C=0`$ condition. We will study splittings with the condition $`B=C=0`$, torus rectangular, and identify counting functions describing the splittings. These functions exhibit symmetries in the Weyl group, not too surprisingly since this is known to preserve these conditions on the moduli . In section three we look at splittings of the charge vector $`(Q_1,Q_5;\stackrel{}{0};0,0)`$ system. Section 4 deals with the splittings of $`(Q_1,Q_5;Q_{12},Q_{34},0,0,0,0;0,0)`$ system. Section 5 describes splittings of $`(Q_1,Q_5;Q_{12},Q_{34},Q_{13},Q_{42},Q_{14},Q_{23};N_1,N_5)`$. 2.3. Liouville theory on long string A Liouville description of long strings was given in . Further discussion from a 2D gauge theory point of view appeared in . We have a pair of branes, one with charge $`q`$ and the other with charge $`Qq`$ splitting from one with charge $`Q`$. The tensions of the split strings also add up to the tension of the string of charge $`Q`$. A Liouville theory on the long string was derived $$S=Tr_0^2\sqrt{g}((\mathrm{\Phi })^2+\mathrm{\Phi }R)$$ where $`T`$ is the tension of the brane and $`r_0`$ is the radius of $`AdS_3`$. The Liouville field $`\mathrm{\Phi }`$ is related to the radial direction. After rescaling the field, the above action is a Liouville action having the background charge $`𝒬=\sqrt{4\pi Tr_0^2}`$. The conformal field theory has a central charge $`c_{Liouville}=3𝒬^2`$ and a massgap $`\mathrm{\Delta }_0=\frac{𝒬^2}{8}`$. We assume now that $`Q^2>>q^2`$ and that the $`AdS`$ geometry is given by $`Q`$. Using the same arguments leading to $`(2.1)`$ we obtain using $`T(q)=qQ/\|Q\|`$ and $`r_0^2=2\pi \|Q\|`$ that $`𝒬=2qQ`$. The previous expression reduces to the expected one in the case of $`q`$ being a $`D1`$ string, i.e giving a central charge of $`6Q_5=6(Q_1Q_5)6(Q_11)(Q_5)`$. In general, we may expect the string of charge $`q`$ to see a geometry given by $`Qq`$. Interpreting the Liouville theory as living on an “interaction string” whose central charge is the difference between $`3Q^2`$ and $`3(Qq)^2+3q^2`$, we have $$𝒬=2q(Qq),c_{Liouville}=6q(Qq)$$ Using the finite $`Q_5`$ result in we would expect by U-duality that a split string of charge $`q`$ would lead to central charges, and thresholds given by : $$\begin{array}{cc}& 𝒬=2q(Qq),\hfill \\ & c=3𝒬^2,\hfill \\ & \mathrm{\Delta }_0=\frac{(𝒬1)^2}{4𝒬}.\hfill \end{array}$$ 3. Counting of splittings of $`(Q_1,Q_5;\stackrel{}{0};0,0)`$ We describe the BPS splittings of the system $`(Q_1,Q_5;0;0,0)`$. We are always working with $`B=C=0`$ and rectangular tori, unless explicitly specified otherwise. We first consider generic radii, and then rational radii, where we further specialize to $`R_i=g_6^2=1`$. 3.1. Splittings of $`(Q_1,Q_5;\stackrel{}{0};0,0)`$ for generic radii. We would like to perform a detailed counting of the splittings for a given charge. Consider the charge $`(N,1)`$. We can split it as follows : $$(N,1)=(q,1)+(q^{},0)$$ where $`q+q^{}=N`$. The possible choices of $`q`$ range from $`0`$ to $`N1`$, so there are $`N`$ of them. For small $`q^{}`$ such BPS splittings lead to a continuum of states as explained in . There are also splittings where we decompose $`(N,1)`$ into a sum of more than two vectors. These will also contribute to the Hamiltonian for string theory on $`ADS3`$ a continuum of states ( with the higher density of states associated with a multistring system as opposed to a two-string system ). To formulate physically sensible counting rules we have to decide whether a subsystem like $`(1,0)+(1,0)`$ should be counted as identical to $`(2,0)`$ or as different. The $`(1,0)`$ system with unit D1-charge can be mapped to an elementary string. Since we understand the perturbative multi-string Hilbert space, we know that in the sector with winding number two, we can have two singly wound states or a doubly wound state. So $`(1,0)+(1,0)`$ should be counted differently. Similarly we can argue that a string with charges $`(2,2)`$ should be counted as distinct from as string $`(1,1)+(1,1)`$. By duality this is related to the fact that with winding number $`2`$ and momentum $`2`$, we can have a single string with these quantum numbers or two separate strings with these quantum numbers. We have two kinds of splittings of the charge vector $`(N,1)`$. The first has the form $`(N,0)=(0,1)+(N_1,0)+\mathrm{}(N_l,0)`$, where $`N`$ is being partitioned into $`l`$ non-zero parts $`N=N_1+N_2+\mathrm{}N_l`$. The second kind of splitting takes the form $`(N,0)=(N_1,1)+(N_2,0)+\mathrm{}(N_l,0)`$. For each partition of $`N`$ we can place the $`1`$ in a number of different ways, but if the partition contains some integer more than once, placing the $`1`$ next to different copies of the integer should not be counted more than once. As we see in the next paragraph this is automatically taken into account by a Fock space description in terms of bosonic oscillators. Equivalently for each partition of $`N`$, there are $`k`$ different ways of placing the $`1`$ where $`k`$ is the number of distinct integers in the partition. So the total number of splittings of the second kind is equal to the sum of $`k`$ over all the partitions of $`N`$. This counting can be written in terms of a function $$\stackrel{~}{P}(x,t)=\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{(1x^l)}\frac{1}{(1tx^l)}$$ Let $`P(x,t)=_{N,m=0}^{\mathrm{}}P(N,m)x^Nt^m`$. The number of splittings is $$N_s(N,1)=P(N,0)+P(N,1)$$ Note that $`P(N,0)`$ is just the number of partitions of $`N`$, and it counts the number of splittings of the form $`(N,1)=(0,1)+(N_1,0)+(N_2,0)+\mathrm{}(N_k,0)`$. The second term $`P(N,1)`$ counts the number of splittings of the form $`(N,1)=(N_1,1)+(N_2,0)+\mathrm{}(N_k,0)`$. This counting can be described in Fock space language. Consider the oscillators $`\alpha _{k,l}`$ where $`k`$ can take values from $`1`$ to $`N`$, and $`l`$ is a discrete index which can take values $`0`$ or $`1`$. $`P(N,0)`$ counts the number of states in a Fock space where all the $`l`$ take the value $`0`$, and the oscillator indices $`k`$ add up to $`N`$. $$\alpha _{N_1,0}\mathrm{}\alpha _{N_k,0}\|0>$$ $`P(N,1)`$ counts the number of states where the $`l`$ values add up to $`1`$, which is the same as saying that one of the $`l`$ values takes the value $`1`$. One of these states is : $$\alpha _{N_1,1}\mathrm{}\alpha _{N_k,0}\|0>$$ The bosonic statistics of the $`\alpha `$ oscillators guarantees that associating the $`1`$ with different copies of the same $`N_i`$ overcounted. A more economical way of writing the number of splittings is to use the generating function : $$P_1(x_1,x_2)=\underset{k,l}{}\frac{1}{(1x_1^kx_2^l)}$$ where $`k`$ and $`l`$ extend from $`0`$ to $`\mathrm{}`$, except that $`k=l=0`$ is not allowed. Define the coefficient $`P_1(n,m)`$ as the coefficients appearing in the expansion : $$P_1(x_1,x_2)=\underset{n,m}{}P_1(n,m)x_1^nx_2^m$$ The number of splittings is just $`P_1(N,1)`$. It is clear that $`P_1(N,1)`$ is equal to $`P_1(1,N)`$. We can also count the number of ways of splitting a more general charge $`(Q_1,Q_5)`$. We can write the number of distinct splittings $`N_s(Q_1,Q_5)`$ in terms of the coefficients of this quantity as follows: $$N_s(Q_1,Q_5)=P_1(Q_1,Q_5)$$ This number is equal to the number of states in a Fock space of the form $$\alpha _{(k_1,l_1)}\alpha _{(k_2,l_2)}\mathrm{}\alpha _{(k_m,l_m)}\|0>$$ where the integers $`k_1`$ through $`k_m`$ add up to $`Q_1`$ and the integers $`l_1`$ through $`l_m`$ add up to $`Q_5`$. Oscillators of the form $`\alpha _{(k,0)}`$ or $`\alpha _{(0,l)}`$ are allowed but $`\alpha _{(0,0)}`$ is not allowed. It is clear that we have the symmetry : $$N_s(Q_1,Q_5)=N_s(Q_5,Q_1).$$ To make the counting even more explicit, we can restrict to splittings into two parts only, to get for the system $`(Q_1Q_5,1)`$ a number $`Q_1Q_5`$. For the system $`(Q_1,Q_5)`$ we get $`\frac{(Q_1+1)(Q_5+1)}{2}1`$. These formulae clearly show that there is a symmetry under exchange of $`Q_1`$ and $`Q_5`$, but no symmetry under exchange of $`(Q_1,Q_5)`$ into $`(Q_1Q_5,1)`$. Note that in all of the above counting we have not included splittings of the form $`(Q_1,Q_5)=(q_1^{},q_5^{})+(q_1+q_1^{},q_5+q_1^{})`$, with $`q_1^{},q_5^{}`$ positive. From the connection to instanton moduli spaces discussed in a later section, this is a natural restriction to consider. For a splitting of this form, when $`q_1^{},q_5^{}`$ are positive and small, the arguments of show that the energy required to create such a long string is infinite and the system is not BPS. A duality invariant way of stating this restriction is that we are considering splittings $`Q=q^{(1)}+\mathrm{}q^{(l)}`$ where the projections of $`q^{(i)}`$ to the 5-plane are not only proportional to $`Q`$ but that the constant of proportionality is positive i.e $`q^{(i)}.Q0`$. 3.2. Refinements of counting Note that to keep the formulae simple, we have counted splittings by merely counting the number of distinct string charges that can appear in the splittings of a given charge. One may also weight the splittings by the number of distinct states associated with the ground states of the split strings. If we choose to do such a counting, we should have, for a charge $`(k,l)`$, a multiplicity of $`p_{16}(kl)`$, where $`p_{16}(kl)`$ is the number of states at level $`kl`$ in a Fock space with $`16`$ distinct oscillators. Such a counting can be done by modifying the above partition functions as follows : We replace in (3.1) the expression $`\frac{1}{(1x_1^kx_2^l)}`$ by $`\frac{1}{(1x_1^kx_2^l)^{p_{16}(kl)}}`$. A further refinement would involve defining a partition function which sums the lowest energy states coming from each splitting by weighting with an energy dependent exponential. This requires further dynamical information on the energy associated with each splitting. The simplest counting problem we have considered, which we extend to systems with more charges in the subsequent sections, suffices for the purpose at hand which is to exhibit the appropriate subgroups of the duality group. 3.3. Splittings of $`(Q_1,Q_5;\stackrel{}{0};0,0)`$ at special radii. Consider the splitting of a $$\begin{array}{cc}\hfill (Q_1,Q_5;0,0,0,0,0,0;0,0)& q+(Qq);\hfill \\ \hfill q& =(q_1,q_5,q_{12},q_{34},0,0,0,0;0,0)\hfill \end{array}$$ According to , this is a BPS splitting if the projection of $`q`$ to the 5-plane is proportional to $`(Q_1,Q_5)`$. Vanishing of the projection along $`E^1`$ leads to the condition $$\frac{q_{12}}{q_{34}}=\frac{R_1R_2}{R_3R_4}$$ This means that such splitting can only happen when $`\frac{R_1R_2}{R_3R_4}`$ is a rational number. While the small instanton singularity discussed in requires tuning $`B`$-fields, these singularities require tuning geometrical moduli. The fact that we need to tune the B-field to zero could be understood from properties of instantons because the instantons are no longer point-like for non-zero generic B-fields. It is interesting that these conditions on the geometry can be understood from properties of instanton moduli spaces ( anti-self dual connections in our conventions). The splitting should correspond to instanton solutions where the connection takes block diagonal form : $$A=\left(\begin{array}{ccc}& ..& ..\\ & A^{(1)}& ..\\ & ..& ..\\ & ..& A^{(2)}\\ & ..& ..\end{array}\right)$$ Here $`A^{(1)}`$ is a connection in $`U(q_5)`$ with instanton number $`q_1`$, and fluxes $$\frac{1}{2\pi }trF_{12}=q_{34},\frac{1}{2\pi }trF_{34}=q_{12}$$ $`A^{(2)}`$ is a $`U(Q_5q_5)`$ connection with instanton number $`Q_1q_1`$, and fluxes $`q_{12},q_{34}`$. Consider first the case of special cases where $`q_1q_5+q_{12}q_{34}=0`$ ( where $`q_1`$ is the number of anti-instantons, proportional to $`trFF`$ ). We can realize this as a $`U(q_5)U(Q_5)`$ instanton configuration built from constant field strengths $`F_{12},F_{34}`$. The flux quantization conditions combined anti-self duality, lead to the rationality conditions on the radii : $$\begin{array}{cc}& F_{12}R_1R_2=q_{34}\hfill \\ & F_{34}R_3R_4=q_{12}\hfill \\ & F_{12}=F_{34}\hfill \\ & \frac{R_1R_2}{R_3R_4}=\frac{q_{34}}{q_{12}}\hfill \end{array}$$ The anti-self duality is required for supersymmetry of $`SU(Q_5)`$ configurations. The diagonal constant $`U(Q_5)`$ part is not constrained. This follows from the SUSY variation of the world-volume fermions : $$\delta _{\xi ,\stackrel{~}{\xi }}\lambda =\xi \mathrm{\Gamma }_{\mu \nu }F^{\mu \nu }+\stackrel{~}{\xi }$$ The diagonal part of the $`U(q_5)`$ configuration is part of the $`SU(Q_5)`$. Such constant field strength solutions (torons ) were written down in and their connection to D-brane bound states studied in . Torons are in this sense closely related to splittings involving $`q^2=0`$. To get configurations with $`q^2>0`$, we have to turn on anti-self dual configurations inside the $`SU(q_5)`$ which can be either point-like or smooth. More generally we can have $$(Q_1,Q_5;0,0,0,0,0,0;0,0)(q_1,q_5;q_{12},q_{34},q_{13},q_{42},q_{14},q_{23};n_1,n_5)+\mathrm{}$$ If all the radii are adjusted to be equal to each other, and the six-dimensional coupling also adjusted to $`1`$ : we now have four independent parameters which can be arbitrary integers, $`(q_{12},q_{13},q_{14},n_1)`$. The remaining parameters are determined by $$\begin{array}{cc}& q_{34}=q_{12}\hfill \\ & q_{42}=q_{13}\hfill \\ & q_{23}=q_{14}\hfill \\ & n_5=n_1\hfill \end{array}$$ The norm of the first vector is equal to $`2(q_1q_5q_{12}^2q_{13}^2q_{14}^2n_1^2)`$. The BPS condition requires that this is greater or equal to zero. When we realize these configurations in gauge theory, we have a condition $`q_5>0`$. The condition $`q^20`$ means that $`q_10`$. Since the $`q_1`$ sum to $`Q_1`$ and the $`q_5`$ sum to $`Q_5`$, we have $`0q_1Q_1`$ and $`0q_5Q_5`$. The condition $`q^20`$ then only allows a finite number of solutions. The splittings of this form can be counted using generating functions : $$P_2(x_1,x_2,x_3,x_4,x_5,x_6)=\underset{k,l,m,n,p,r}{}\frac{1}{(1x_1^kx_2^lx_3^mx_4^nx_5^px_6^r)}$$ The product is constrained by $$\begin{array}{cc}& k0,l0\hfill \\ & klm^2n^2p^2r^20\hfill \end{array}$$ Defining $`P(n_1,n_2,n_3,n_4,n_5,n_6)`$ as the coefficient of $`x_1^{n_1}x_2^{n_2}x_3^{n_3}x_4^{n_4}x_5^{n_5}x_6^{n_6}`$ in $`P(x_1,x_2,x_3,x_4,x_5,x_6)`$, the desired counting function is $`P(Q_1,Q_5,0,0,0,0)`$. At the special radii and $`g_6`$ we have a large class of splittings involving strings with extra 3-anti3 charges and NS1-5 charges $`(q_1,q_5;q_{12},q_{34},q_{13},q_{42},q_{14},q_{23};n_1,n_5)`$. We saw that the charges $`q_{34},q_{42},q_{23},n_5`$ are determined by the charges $`q_{12},q_{13},q_{14},n_1`$ respectively. There is an $`S_4Z_2`$ subgroup of the Weyl group which acts on the charges of these split strings. The $`S_4`$ just permutes the four independent charges and the $`Z_2`$ acts by a reflection $`q_{12}q_{12}`$. The counting function $`P(Q_1,Q_5)`$ of course has the symmetry of exchanging $`Q_1`$ and $`Q_5`$. A $`Z_2\times (S_4Z_2)`$ subgroup of $`(S_5Z_2)`$ is therefore manifest here. The first factor is a symmetry acting on the charge of the string we start with i.e exchanging $`Q_1,Q_5`$. The second is a symmetry acting on the strings that can appear in the splitting. By choosing $`R_i=g_6^2=1`$ we made sure that the $`S_4Z_2`$ acts as a symmetry of the Hamiltonian for the corresponding ADS background. For more general radii satisfying the rationality conditions, a split string for one gemoetrical modulus is mapped to a split string at another geometrical modulus. 4. Splittings of a system $`(Q_1,Q_5;Q_{12},Q_{34},\stackrel{}{0};0,0)`$ Consider first the splittings of $$(Q_1,Q_5,Q_{12},Q_{34})=(q_1,q_5;q_{12},q_{34},0,0,0,0;0,0)+\mathrm{}$$ We again consider rectangular tori with $`B=C=0`$. The condition on the near horizon geometry takes the form : $$\begin{array}{cc}& \frac{Q_1}{Q_5}=V\hfill \\ & \frac{Q_{12}R_1R_2}{Q_{34}R_3R_4}=1\hfill \end{array}$$ The condition that the projection of the vector $`q`$ to the five-plane is parallel to $`Q`$ takes the form : $$\frac{(q_1Q_5+q_5Q_1)}{Q_1Q_5}=\frac{(q_{12}Q_{34}+q_{34}Q_{12})}{(Q_{12}Q_{34})}$$ If we consider splittings of the more general form : $$(Q_1,Q_5,Q_{12},Q_{34})=(q_1,q_5,q_{12},q_{34},q_{13},q_{42},q_{14},q_{23},n_1,n_5)+\mathrm{}$$ we still have the above conditions and we have restritions on the moduli of the form : $$\begin{array}{cc}& \frac{q_{13}R_1R_3}{q_{42}R_4R_2}=1\hfill \\ & \frac{q_{14}R_1R_4}{q_{23}R_3R_2}=1\hfill \\ & \frac{n_1g_6^2}{n_5}=1\hfill \end{array}$$ Let us first consider splittings of the first kind. These can be counted as the coefficient of $`x_1^{Q_1}x_2^{Q_5}x_3^{Q_{12}}x_4^{Q_{34}}`$ in the series $$P(x_1,x_2,x_3,x_4)=\underset{q_1,q_5,q_{12},q_{34}}{}\frac{1}{(1x_1^{q_1}x_2^{q_5}x_3^{q_{12}}x_4^{q_{34}})}$$ Here the $`q`$’s obey the conditions in (4.1), in addition to the usual $`q^20`$ and $`q.Q0`$. A class of solutions of these conditions ( we haven’t proved that they give a complete set of solutions ) can be written as $$q=\lambda _1q_{(1)}+\lambda _2q_{(2)}$$ where $`q_{(1)}=(\frac{Q_1}{\lambda _L},0;\frac{Q_{12}}{\lambda _L},0,0,0,0,0;0,0)`$ and $`q_{(2)}=(0,\frac{Q_5}{\lambda _R},0;0,\frac{Q_{34}}{\lambda _R},0,0,0,0,0;0,0)`$. Here $$\lambda _L=\mathrm{gcd}(Q_1,Q_{12});\lambda _R=\mathrm{gcd}(Q_5,Q_{34})$$ and $$0\lambda _1\lambda _L;0\lambda _2\lambda _R.$$ Splittings coming from this class of solutions can be counted as the coefficient of $`x_1^{\lambda _L}x_2^{\lambda _R}`$ in the series : $$\underset{\lambda _1,\lambda _2}{}\frac{1}{(1x_1^{\lambda _1}x_2^{\lambda _2})}$$ This counting function $`N_s(Q_1,Q_5,Q_{12},Q_{34})`$ is symmetric under a group $`S_2Z_2`$. The $`S_2`$ exchanges the $`(Q_1,Q_5)`$ pair with the $`(Q_{12},Q_{34})`$ pair. The $`Z_2`$ generator can be taken to be the exchange of $`(Q_1,Q_5)`$ to $`(Q_5,Q_1)`$. As in the previous section we can proceed to the case where the other ratios of radii and the coupling are adjusted according to (4.1) and we can further specialize to the case where all $`R_i=1=g_6^2`$, to find a group $`S_3Z_2`$ acting on the allowed extra charges. 5. Splittings of a system $`(Q_1,Q_5,Q_{12},Q_{34},Q_{13},Q_{42},Q_{14},Q_{23},N_1,N_5)`$ The condition that this charge be parallel to the 5-plane gives : $$\begin{array}{cc}& \frac{Q_1}{Q_5V}=1\hfill \\ & \frac{Q_{12}R_1R_2}{Q_{34}R_3R_4}=1\hfill \\ & \frac{Q_{13}R_1R_3}{Q_{42}R_2R_4}=1\hfill \\ & \frac{Q_{14}R_1R_4}{Q_{23}R_2R_3}=1\hfill \\ & \frac{N_1g_6^2}{N_5}=1\hfill \end{array}$$ Consider splittings of the form $`Qq+\mathrm{}`$, where the charge $`q`$ has components $`(q_1,q_5;q_{12},q_{34},q_{13},q_{42},q_{14},q_{23};n_1,n_5)`$. We have conditions $`q^20`$, and requiring that the projection of $`q`$ to the 5-plane be parallel to $`Q`$ leads to conditions : $$\begin{array}{cc}& \frac{(q_1Q_5+q_5Q_1)}{Q_1Q_5}=\frac{(q_{12}Q_{34}+q_{34}Q_{12})}{Q_{12}Q_{34}}=\frac{(q_{13}Q_{42}+q_{42}Q_{13})}{Q_{13}Q_{42}}\hfill \\ & =\frac{(q_{14}Q_{23}+q_{23}Q_{14})}{Q_{14}Q_{23}}=\frac{(n_1N_5+n_5N_1)}{N_1N_5}\hfill \end{array}$$ A class of solutions to these equations takes the form $$\begin{array}{cc}& \lambda _1q^{(1)}+\lambda _2q^{(2)},\hfill \\ & 0\lambda _1\lambda _L\hfill \\ & 0\lambda _2\lambda _R\hfill \end{array}$$ where $`\lambda _L=\mathrm{gcd}(Q_1,Q_{12},Q_{13},Q_{14},N_1)`$ and $`\lambda _R=\mathrm{gcd}(Q_5,Q_{34},Q_{42},Q_{23},N_5)`$, and $`q^{(1)}=\frac{1}{\lambda _L}(Q_1,0,Q_{12},0,Q_{13},0,Q_{14},0,N_1,0)`$ and $`q^{(2)}=\frac{1}{\lambda _R}(0,Q_5,0,Q_{34},0,Q_{42},0,Q_{23},0,N_5)`$. The number of these splittings $`N_s(Q_1,Q_5,Q_{12},Q_{34},Q_{13},Q_{42},Q_{14},Q_{23};N_1,N_5)`$ can be counted using the coefficient of $`x_1^{\lambda _L}x_2^{\lambda _R}`$ in the following series: $$\underset{\lambda _1,\lambda _2}{}\frac{1}{(1x_1^{\lambda _1}x_2^{\lambda _2})},$$ where $`\lambda _10,\lambda _20`$, but $`(\lambda _1,\lambda _2)=(0,0)`$ is excluded. The function $`N_s(Q_1,Q_5,Q_{12},Q_{34},Q_{13},Q_{42},Q_{14},Q_{23};N_1,N_5)`$ is invariant under $`S_5Z_2`$. The $`Z_2`$ can be taken to be the exchange of $`(Q_1,Q_5)`$ to $`(Q_5,Q_1)`$. The $`S_5`$ permutes the five blocks : $`(Q_1,Q_5)`$, $`(Q_{12},Q_{34})`$, $`(Q_{13},Q_{42})`$, $`(Q_{13},Q_{42})`$ , $`(Q_{14},Q_{23})`$ and $`(N_1,N_5)`$. 6. Beyond $`B=C=0`$. The duality group $`O(5,5:Z)`$ can be used to map the charges $`(Q_1,Q_5,0,0)`$ to the charge $`(Q_1Q_5,1,0,0)`$. An $`O(2,2)`$ subgroup of the $`O(5,5)`$ will suffice to do that, as pointed out by : $$\left(\begin{array}{ccccccccc}& adQ_5^2& & bcQ_1^2& & acQ_1Q_5& & bdQ_1Q_5& \\ & bc& & ad& & ac& & bd& \\ & abQ_5& & abQ_1& & a^2Q_5& & b^2Q_1& \\ & cdQ_5& & cdQ_1& & c^2Q_1& & d^2Q_5& \end{array}\right)$$ Here $`a,b,c,d`$ have been chosen to satisfy the condition $`adQ_5bcQ_1=1`$. In fact there exists a choice of $`c=d=1`$ which leads to $$\left(\begin{array}{ccccccccc}& aQ_5^2& & bQ_1^2& & aQ_1Q_5& & bQ_1Q_5& \\ & b& & a& & a& & b& \\ & abQ_5& & abQ_1& & a^2Q_5& & b^2Q_1& \\ & Q_5& & Q_1& & Q_1& & Q_5& \end{array}\right)$$ . These matrices allow us to see explicitly that these transformations which violate the $`B=C=0`$ condition. We can use this map start with the splittings of a $`(Q_1,Q_5)`$ system to get splittings of a $`(Q_1Q_5,1)`$ system. The splittings of $`(Q_1,Q_5)=(4,3)`$ are studied here. This can be mapped by $`O(5,5;Z)`$ to the charge $`(12,1)`$. A first class of splittings of $`(4,3)`$ involving no extra charges is : $$\begin{array}{cc}\hfill (4,3)& (4,0)+(0,3)\hfill \\ & (4,1)+(0,2)\hfill \\ & (4,2)+(0,1)\hfill \\ & (3,0)+(1,3)\hfill \\ & (3,1)+(1,2)\hfill \\ & (3,2)+(1,1)\hfill \\ & (3,0)+(1,0)\hfill \\ & (2,0)+(2,3)\hfill \\ & (2,1)+(2,2)\hfill \end{array}$$ We also have $`12`$ splittings of the charge $`(12,1)`$ into different charges of the form $`(k,1)+(12k,0)`$ where $`k`$ ranges from $`11`$ to $`0`$. When these are mapped to splittings of the $`(4,3)`$ system using $`O(5,5;Z)`$ ( in fact $`O(2,2;Z)`$ suffices ), we get a bunch of splittings of $$(4,3)(16k,k9,k12,k12)+(k12,12k,12k,12k)$$ This shows that the class of counting functions etc. that we have described can be used to obtain information about splittings at moduli which go beyond the $`B=C=0`$ condition. In general the $`B,C`$ values have explicit dependence on $`Q_1`$ and $`Q_5`$, and such dependence will show up in the coefficients of marginal operators necessary to go from the free point of the orbifold description in terms of $`S^{Q_1Q_5}(T)`$ to the point which has a simple description in terms of $`(Q_1,Q_5)`$ system . The data associated with splittings and thresholds etc. seems best viewed as living on the space $$\mathrm{\Gamma }^{5,5}\times O(5,5;R)/O(5)\times O(5),$$ with a manifest $`O(5,5;Z)`$ symmetry which acts on both charges and moduli. It will be interesting to see if the class of splittings at $`B=C=0`$, a large class of which we have considered in this paper, exhausts all the physics of long strings that one encounters as one moves in the space (6.1). 7. Instantons Properties of instanton moduli spaces have come up in the discussions of the splittings above. Some aspects are developed in more detail here, and some puzzling features discussed. First consider the splittings of $`(Q_1,Q_5)`$ which do not involve extra 3-brane charges. They correspond to a decomposition of the compactified moduli space of instantons. We will discuss some features <sup>1</sup> We thank George Daskalopoulos for instructive discussions on this subject of the Uhlenbeck compactification . The compactified moduli space of instantons on $`T^4`$ for gauge group $`U(r)`$ and instanton number $`k`$, $`\overline{}_{(r,k)}`$, admits a double stratification, labelled by two integers $`(f,p)`$ which count the number of $`U(1)`$ flat connections and the number of point-like instantons respectively. T-duality inverts the volume, exchanges rank and instanton number, and flips the integers $`f`$ and $`p`$. $$\overline{}_{r,k}=\underset{(f,p)}{}(_{rf,kp}^{(0)}\times S^p(T)\times S^f(T^{}))(S^k(T)\times S^r(T^{}))$$ The Moduli spaces $`^{(0)}`$ appearing in the above involve no flat connections and no point-like instantons. The integer $`f`$ extends from $`0`$ to $`r2`$ and $`p`$ extends from $`0`$ to $`k2`$. Naively one may have allowed them to extend to $`r1`$ or $`k1`$, but $`_{1,l}`$ and $`_{l,1}^{(0)}`$ are empty for any $`l`$. In fact such strata would be expected if we want a compactification which knows about all the possible split strings, including those of charge $`(l,1)`$ as explained below. Let us consider first the case where $`r,k`$ are large, and $`f,p`$ are small. The symmetric product $`S^p(T)`$ describes $`p`$ long strings moving on $`T`$. The sigma model on this stratum contains sectors parametrized by partitions of $`p`$ which describe different ways of partitioning the $`p`$ long D1-strings into bound states. Similarly the symmetric product $`S^f(T)`$ contains sectors describing $`f`$ long D5-strings. A stratum parametrized by $`(p,f)`$ contains the physics of splittings of the type $`(Q_1,Q_5)(Q_1f,Q_5p)+(f_1,0)+\mathrm{}(f_l,0)+(0,p_1)+\mathrm{}(0,p_k)`$ where $`p=p_1+p_2+\mathrm{}p_k`$ and $`f=f_1+f_2+\mathrm{}f_l`$ are partitions of $`p`$ and $`f`$. Geometrical structures corresponding to $`(f,p)`$ bound states with both $`f`$ and $`p`$ non-zero do not seem to appear. Another puzzle appears when we allow $`p`$ to be comparable to $`Q_5`$. When $`p`$ is equal to $`Q_51`$ we do not have the corresponding stratum, whereas the algebraic calculation of allowed BPS aplittings continues to allow such splittings. Perhaps other compactifications, e.g the Gieseker compactification, do include the extra strata which would give a more detailed correspondence between strata and splittings. As discussed in section 3.2, when we consider splittings involving extra 3-brane charges we need to consider reducible connections, i.e of block diagonal form where each block contains non-zero flux, but the fluxes add up to zero. When the original system contains extra 3-brane charges, we need to start with a bundle of rank $`Q_5`$ instanton number $`Q_1`$ and magnetic fluxes $`ϵ^{ijkl}Q_{kl}`$. The BPS splitting conditions in (5.1) for example do not allow a point-like instanton of charges $`(1,0,0,0\mathrm{})`$. This means that for such a choice of moduli and charges, the instanton moduli spaces will have no strata corresponding to shrinking instantons. Rather the strata will correspond to solutions of (5.1). 8. Splitting NS charges from D-brane systems The splittings of the $`(Q_1,Q_5)`$ system which involve extra charges like three branes are easily described in the gauge theory of $`U(Q_5)`$ which contains magnetic fluxes allowing the description of the appropriate splittings. But the splittings which involve extra charges NS1-NS5 seem to be harder to describe. The NS1-charges are of course described by electric fluxes. We don’t know how to describe the NS5-charges as a flux in the D-brane gauge theory. This is the famous transverse 5-brane problem in attempts to use Yang Mills as Matrix theory for compactifications on tori of dimension larger than $`4`$. The transverse 5-brane problem is solved by appealing to little string theory but available descriptions of this theory do not allow an explicit description of these splittings which generalize in a simple way the description in terms of point-like instantons, flat connections, block-diagonal connections that we have described. A description which stays within the confines of moduli spaces of six-dimensional gauge theories, can be given at the expense of giving up the restriction to a fixed rank gauge group. The following remarks have some analogies to developments on duality in Matrix theory appearing in . An example of a splitting which involves the NS1-NS5 charges is $$(Q_1,Q_5;\stackrel{}{0};0,0)(q_1,q_5;\stackrel{}{0};1,1)+(Q_1q_1,Q_5q_5;\stackrel{}{0};1,1)$$ For such a splitting to exist at $`B=C=0`$, we need to adjust the six-dimensional coupling $`g_6^2=1`$. One way to give a gauge theoretic description of such a system is to map the charge $`(q_1,q_5;\stackrel{}{0};1,1)`$ to $`(\stackrel{~}{q}_1,\stackrel{~}{q}_5;\stackrel{}{0};0,0)`$, with $`\stackrel{~}{q}_1\stackrel{~}{q}_5=q_1q_51`$, at the cost of turning on some non-trivial RR and/or B-fields. But this is a system which we can describe using $`U(\stackrel{~}{q}_5)`$ gauge theory with instanton number $`\stackrel{~}{q}_1`$, with non-trivial couplings turned on due to the background fields. Similarly we can map $`(Q_1q_1,Q_5q_5;\stackrel{}{0};1,1)`$ to $`(q_1^{},q_5^{})`$ with $`q_1^{}q_5^{}=(Q_1q_1)(Q_5q_5)1`$, obtaining a $`U(q_5^{})`$. Generically $`q_5^{}+\stackrel{~}{q}_5Q_5`$, so we cannot think of this in terms of something happening within $`U(Q_5)`$ gauge theory. It is intriguing that whereas all the other splitting processes could be understood, at least roughly, in terms of some apprpriate compactification of instanton moduli spaces, we here have to go beyond gauge theory at a fixed rank. It should be interesting to explore the systematic use of gauge theories of arbitrary rank to understand the full physics of the $`(Q_1,Q_5)`$ system, using duality along the above lines. 9. Summary and Outlook We found that counting functions related to the number of splittings of a string in six dimensions contains interesting information on symmetries associated with the phenomena of long strings in ADS3 and the associated continuums of the spectrum. The Weyl sub-group of $`O(5,5;Z)`$ played an interesting role. The nature of the BPS splittings also gives information about the structure of the compactifications of instanton moduli spaces. The symmetries of BPS splittings are related to symmetries acting on strata of the compactifications. Some puzzles were raised regarding the details of this correspondence. It would be very interesting if these puzzles can be solved and the Fock space structures we described can be derived from appropriate compactifications of ( possibly $`\alpha ^{}`$ and $`g_s`$-corrected ) instanton moduli spaces, thus generalizing the familiar relations between Fock spaces and instanton moduli spaces . While we focused here on $`T^4`$, the discussion is easily generalized to $`K3`$ and similar connections to instanton moduli spaces on $`K3`$ should exist. When we consider splittings of the $`(Q_1,Q_5)`$ system which include vectors carrying non-trivial NS charges, we need to go beyond conventional compactifications of instanton moduli spaces to even get a rough gauge theoretic understanding of the splittings. One avenue is to use facts about the duality group $`O(5,5;Z)`$ to show that we may need to go beyond instanton moduli spaces for bundles of a fixed rank to get a full gauge theory understanding of the $`(Q_1,Q_5)`$ system. Whereas we used the counting of long strings to identify symmetries of the physics of long strings. It will be interesting to explore the consequences of these symmetries for the detailed dynamics, e.g the thresholds where the density of continuum states jumps. The counting problems themselves allow refinements mentioned in section 3. 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warning/0002/cond-mat0002342.html	ar5iv	text	# 1 Introduction. ## 1 Introduction. Mean field spin glass models are considered as a prototype of disordered, frustrated systems and, more generally, of a large class of complex systems that can be successfully analyzed using the ideas developed in the study of spin glasses . Among these, the Sherrington–Kirkpatrick model has a primary importance. This model is by now well understood in its general features, as described by Parisi with an ingenious method and the ultrametric Ansatz . This picture has been confirmed by extensive numerical simulations and some rigorous results . In particular, F. Guerra has given a rigorous motivation for the introduction of a functional order parameter of Parisi type, and has shown how in this framework a simple Ansatz allows to express the thermodynamic variables and some physical observables in terms of that order parameter . In the present paper, the Ansatz of Guerra is extended, and is developed a method to express all physical observables in terms of the functional order parameter, in a mathematically rigorous framework. The method is general, and can be applied to other mean field disordered models such as the multi-spin interaction spin glass and the neural networks. It is well known that the whole physical content of mean field spin glass models is contained in the overlap random variables. Given $`s`$ replicas there are $`s(s1)/2`$ overlaps between them, where $`s`$ ranges on the natural numbers. Therefore, the physics of the model is fully contained in a probability distribution on an infinite-dimensional space. Overlaps do not fluctuate in the hight temperature phase : the Sherrington–Kirkpatrick solution turns out to be correct and the overlap distribution is trivial. In the low temperature phase this cannot happen : overlaps do fluctuate . Fluctuations are constrained by the symmetry under permutations of replicas and by the gauge symmetry. Thermodynamical constraints are expressed by Ghirlanda–Guerra relations, in the slightly stronger case when suitable infinitesimal interactions are added to the Hamiltonian (this is also known as the stochastic stability property ). By the Ansatz that the overlap distribution is ultrametric, Parisi gave a solution of the model, in terms of a functional order parameter . Ultrametricity is a simple constraint on the support that considerably simplifies the overlap distribution : together with the previously stated constraints, it reduces the problem to the determination of the mono-dimensional, one-overlap distribution $`P_{12}`$. This is proven in the last section of this paper. A functional order parameter of Parisi type can be introduced rigorously to give a functional representation of the marginal martingale function, and therefore of the free energy . This representation is not unique : there is an infinite set of functional order parameters giving rise to the same free energy. By an Ansatz on this representation, some overlap correlation functions has been expressed through the functional order parameter . In this paper, by an extension of the Ansatz, we explicitly find the full overlap distribution in terms of the functional order parameter, and we show how ultrametricity naturally emerges. The method (and the paper) goes as follows. We introduce a generating functional of physical observables (i.e., expectations of overlap functions), derived from the marginal martingale function (sect. 3). Through the solution of a non-linear antiparabolic equation, and exploiting the Ansatz, we represent it in terms of the functional order parameter $`x`$ (sect. 4). Then, we solve the antiparabolic equation by asymptotic expansion and explicitly find the overlap probability distribution. The physical interpretation of the functional order parameter is obtained, and ultrametricity of overlaps is derived as a natural consequence of the branching diffusion process underlying the equation (sect. 5). Finally, it is shown that complete ultrametricity of overlaps results from ultrametricity of the 3-replicas overlap distribution. Moreover, it is proved that ultrametricity and the Ghirlanda–Guerra identities are complementary in order to determine the full overlap distribution, in the sense that one can hold independently of the other, but together they determine explicitly the overlap measure in terms of the one-overlap distribution $`P_{12}`$ (sect. 6). ## 2 Overlaps in the Sherrington–Kirkpatrick model. The mean field model of a spin glass is defined on sites $`i=1,2,\mathrm{},N`$. To each site is assigned the Ising spin variable $`\sigma _i=\pm 1`$, so that a configuration of the system is described by the application $`\sigma :i\sigma _iZ_2=\{1,1\}`$. The spins on two different sites $`i`$ and $`j`$ are coupled through the random variables $`J_{ij}`$, all independent from each other and equally distributed. For the sake of simplicity we assume a Gaussian distribution, with $$\mathrm{E}(J_{ij})=0,\mathrm{E}(J_{ij}^2)=1,$$ (1) where $`\mathrm{E}`$ denotes averages on the $`J`$ variables. The $`J_{ij}`$’s are called quenched variables, because they do not participate to thermalisation. The Hamiltonian of the Sherrington–Kirkpatrick model is $$H_N(\sigma ,J)=\frac{1}{\sqrt{N}}\underset{(i,j)}{}J_{ij}\sigma _i\sigma _j,$$ (2) where the sum extends over all the $`N(N1)/2`$ couples of sites. The normalization factor $`1/\sqrt{N}`$ is needed to have the correct behavior of the thermodynamic variables in the limit $`N\mathrm{}`$. Denoting with $`\beta `$ the inverse temperature (in proper units), we introduce the partition function $`Z_N(\beta ,J)`$ and the free energy density $`f_N(\beta ,J)`$ : $`Z_N(\beta ,J)`$ $`=`$ $`{\displaystyle \underset{\sigma _1\mathrm{}\sigma _N}{}}e^{\beta H_N(\sigma ,J)},`$ (3) $`\beta f_N(\beta ,J)`$ $`=`$ $`{\displaystyle \frac{1}{N}}\mathrm{log}Z_N(\beta ,J).`$ (4) The associated Boltzmann state $`\omega _{N,\beta ,J}`$ is defined by $$\omega _{N,\beta ,J}(A)=\frac{1}{Z_N(\beta ,J)}\underset{\sigma _1\mathrm{}\sigma _N}{}A(\sigma )e^{\beta H_N(\sigma ,J)},$$ (5) for a generic function $`A`$ of the spin variables. Another relevant quantity is the average of internal energy density $`u_N(\beta )`$ $$u_N(\beta )=\frac{1}{N}\mathrm{E}\omega _{N,\beta ,J}(H_N(\sigma ,J))=\mathrm{E}\frac{}{\beta }(\beta f_N(\beta ,J)).$$ (6) In the thermodynamic limit the free energy density is self-averaging in quadratic mean . For the internal energy density the same property has been proven for almost all values of $`\beta `$, but is believed to hold without restrictions . One of the main features of the mean field spin glass model is the existence of observables that do not self-average in the thermodynamic limit. This is one of the fundamental intuitions contained in the Parisi Ansatz of replica symmetry breaking. Indeed Pastur and Shcherbina have proven that if a suitably chosen order parameter (coming from the response of the system to an external random field) is self-averaging in the thermodynamic limit, then the solution of the model has the Sherrington–Kirkpatrick form : this is unphysical at high $`\beta `$, because it gives negative entropy. Moreover, self-averaging of the Edward–Anderson order parameter implies that the overlap distribution is the trivial one corresponding to the replica symmetric Ansatz of S.–K. . Let us consider $`s`$ copies (replicas) of the system, whose configurations are given by the Ising spin variables $`\sigma _i^{(1)},\mathrm{},\sigma _i^{(s)}`$, and denote with $`\omega _J^{(a)},a=1,2,\mathrm{},s`$ the relative Boltzmann states, the dependence on $`\beta `$ and $`N`$ being understood. We introduce the product state $`\mathrm{\Omega }_J`$ by $$\mathrm{\Omega }_J=\omega _J^{(1)}\omega _J^{(2)}\mathrm{}\omega _J^{(s)},$$ (7) where all the states $`\omega _J^{(a)}`$ are subject to the same values of the quenched variables $`J`$, and the same temperature $`\beta `$. The overlap between the two replicas $`a`$ and $`b`$, $`Q_{ab}`$, is defined by $$Q_{ab}=\frac{1}{N}\underset{i}{}\sigma _i^{(a)}\sigma _i^{(b)},$$ (8) with the obvious bounds $`1Q_{ab}1`$. The importance of overlaps lies in the fact that all physical observables can be expressed in the form $$\mathrm{E}\mathrm{\Omega }_J[F(Q_{12},Q_{13},\mathrm{})],$$ (9) for some function $`F`$. For $`F`$ smooth, we can introduce the random variables $`q_{12},q_{13},\mathrm{}`$, through the definition of their averages $$F(q_{12},q_{13},\mathrm{})=\mathrm{E}\mathrm{\Omega }_J[F(Q_{12},Q_{13},\mathrm{})].$$ (10) Notice that the expectation $``$ includes both the thermal average and the average $`\mathrm{E}`$ over disorder. The overlap distribution carries the whole physical content of the model . Let us recall some considerations about the overlap distribution. The average $`\mathrm{E}`$ over quenched variables introduces correlation between different groups of replicas. For example we have, in general, $$q_{12}^2q_{34}^2q_{12}^2q_{34}^2.$$ (11) The $``$ average is obviously invariant under permutations of replica indices (e.g. $`q_{12}^2q_{13}^2=q_{23}^2q_{13}^2`$, $`q_{12}^2=q_{34}^2`$). Moreover, it is invariant under the gauge transformations defined by $$q_{ab}\epsilon _aq_{ab}\epsilon _b,$$ (12) where $`\epsilon _a=\pm 1`$. This is an easy consequence of the fact that each of the $`\omega _J^{(a)}`$ is an even state on the respective $`\sigma ^{(a)}`$. It follows, for instance, that polynomials in the overlaps which are not gauge invariant have zero mean. These symmetries furnish important restrictions on the the overlap distribution, but even more important constraints have been given by , using simple arguments based on convexity properties and positivity of fluctuations. Consider $`s`$ replicas, and the $`s(s1)/2`$ overlaps between them. Let us denote by $`𝒜_s`$ the associated algebra of observables. Introduce the overlap $`q_{a,s+1}`$, between replica $`a`$ and an additional replica $`s+1`$, and consider the conditional probability distribution $`\stackrel{~}{P}_{(a,s+1)}(q_{a,s+1}\|𝒜_s)`$ of $`q_{a,s+1}`$ given the overlaps among the first $`s`$ replicas. By adding to the Hamiltonian suitable infinitesimal external fields, and taking the thermodynamic limit with a careful procedure, Guerra and Ghirlanda have demonstrated that the following theorem holds for a very general class of probability measures, including short range models . ###### Theorem 2.1 Given the overlaps among $`s`$ replicas, the overlap between one of these, let say $`a`$, and an additional replica $`s+1`$ is either independent of the former overlaps, or it is identical to one of the overlaps $`q_{ab}`$, with $`b`$ running from 1 to $`s`$, excluding $`a`$. Each of these cases have probability $`1/s`$: $`\stackrel{~}{P}_{(a,s+1)}(q_{a,s+1}\|𝒜_s)={\displaystyle \frac{1}{s}}P_{12}(q_{a,s+1})+{\displaystyle \frac{1}{s}}{\displaystyle \underset{ba}{}}\delta (q_{a,s+1}q_{ab}).`$ (13) Results of this kind have been obtained by Parisi in the frame of replica method , and by Aizenmann and Contucci . ## 3 A generator of overlap distributions. Let $`\omega `$ be a generic even state on the Ising spins $`\sigma _1,\mathrm{},\sigma _N`$, possibly depending on the quenched variables $`J_{ij}`$ and let $`f__1:\mathrm{IR}\mathrm{IR}`$ be an even, convex function, such that $`\|f__1(y)\|c\|y\|`$ asymptotically with $`\|y\|\mathrm{}`$ for some positive $`c`$. Let the generating functional $`\psi _N(\omega ,f__1)`$ be given by $$\psi _N(\omega ,f__1)=\mathrm{E}\mathrm{log}\omega (\mathrm{exp}f__1(h_N(\sigma ,J))),$$ (14) where $`h_N(\sigma ,J)=N^{1/2}_iJ_i\sigma _i`$ is the cavity field, and the $`J_i`$’s are fresh noise with the same properties of $`J_{ij}`$. The functional $`\psi _N(\omega ,f_1)`$ contains all informations on the distribution of the replicated cavity fields $`h^{(a)}h_N(\sigma ^{(a)},J)`$. That, in turn, is related to the overlap distribution through the well known formula $$\mathrm{E}\mathrm{\Omega }_J\left(\mathrm{exp}\left(i_ak_ah^{(a)}\right)\right)=\mathrm{exp}\left(_{a,b}k_ak_bq_{ab}/2\right).$$ (15) We expand the logarithm in power series and we introduce replicas: $$\begin{array}{ccc}\psi _N(\omega ,f__1)\hfill & =& \mathrm{E}\left[\mathrm{ln}\left(1\omega \left(1\mathrm{exp}\left(f__1(h)\right)\right)\right)\right]\hfill \\ & =& \underset{s=1}{\overset{\mathrm{}}{}}\frac{1}{s}\mathrm{E}\left[\omega \left(1\mathrm{exp}\left(f__1(h)\right)\right)\right]^s\hfill \\ & =& \underset{s=1}{\overset{\mathrm{}}{}}\frac{1}{s}\mathrm{E}\mathrm{\Omega }_J\left[\underset{a=1}{\overset{s}{}}\left(1\mathrm{e}^{f__1(h^{(a)})}\right)\right]\hfill \end{array}$$ (16) where $`h`$ denotes the cavity field and $`h^{(a)}`$ its replicas. Let us introduce the generalized Fourier transform $`\varphi `$, which is a well defined even generalized function : $$1\mathrm{exp}\left(f__1(y)\right)=_{\mathrm{}}^{\mathrm{}}dk\varphi (k)e^{iky}$$ (17) By the convenient replacement $`\phi (k)\varphi (k)\mathrm{exp}(k^2/2)`$, we finally have $$\psi _N(\omega ,f__1)=\underset{s=1}{\overset{\mathrm{}}{}}\frac{1}{s}\mathrm{d}^sk\underset{a=1}{\overset{s}{}}\phi (k_a)\mathrm{exp}\left(_{(a,b)}k_ak_b\mathrm{q}_{ab}\right)$$ (18) where the sum in the exponential is over the couples $`(a,b)`$, for $`1a<bs`$. The dependence of the r.h.s. on $`f_1`$ is through the function $`\phi `$. Notice that terms $`s=2,3`$ contain the characteristic functions of the distributions of overlaps among 2 and 3 replicas, respectively. It important to notice that the thermodynamic functions can be represented through the functional $`\psi _N(\omega ,f_1)`$. Consider the case $`f_1(y)=\mathrm{log}\mathrm{cosh}\beta y`$, and the corresponding function $`\psi _N^{}(\beta )\psi _N(\omega ,\mathrm{log}\mathrm{cosh}\beta )`$, where $`\omega `$ is the Boltzman state of SK model. Then the following holds . ###### Proposition 3.1 Assume the existance of the limit $`lim_N\mathrm{}\psi _N^{}(\beta )=\psi ^{}(\beta )`$, uniformly on a compact region $`0\beta \stackrel{~}{\beta }`$, with $`\psi ^{}`$ continuous in $`\beta `$, as a consequence. Let us define $$\alpha (\beta )=\mathrm{log}2+_0^1\psi ^{}(\beta \sqrt{1q})𝑑q,$$ (19) so that the $`\beta `$ derivative $`\alpha ^{}(\beta )`$ exists and the following holds $$\alpha (\beta )+\beta \alpha ^{}(\beta )/2=\mathrm{log}2+\psi ^{}(\beta ).$$ (20) Then, we have, for $`0\beta \stackrel{~}{\beta }`$, $$\underset{N\mathrm{}}{lim}\frac{1}{N}\mathrm{E}\left(\mathrm{log}Z_N(\beta ,J)\right)=\alpha (\beta ),\underset{N\mathrm{}}{lim}\frac{1}{N}_\beta \mathrm{E}\left(\mathrm{log}Z_N(\beta ,J)\right)=\alpha ^{}(\beta ).$$ (21) ## 4 The functional order parameter. In the frame of the cavity method, a functional order parameter of Parisi type was introduced by Guerra as a functional representation of the marginal martingale function . Then, he showed that by a simple Ansatz some overlap correlation functions can be expressed in terms of the functional order parameter . In this section we give an extension of the representation Theorem, thus obtaining a functional representation of the physical observables. Exploiting the Ansatz, the generating functional $`\psi _N(\omega ,f_1)`$ is expressed in terms of the functional order parameter. Therefore, the explicit form of the overlap distribution can be extracted. Let us introduce the convex set $`𝒳`$ of functional order parameters of the type $$x:[0,1]qx(q)[0,1],$$ (22) with the $`L^1(dq)`$ distance norm. We induce on $`𝒳`$ a partial ordering, by defining $`x\overline{x}`$ if $`x(q)\overline{x}(q)`$ for all $`0q1`$, and introduce the extremal order parameters $`x_0(q)0`$ and $`x_1(q)1`$, such that for any $`x`$ we have $`x_0(q)x(q)x_1(q)`$. For each $`x`$ in $`𝒳`$, and for suitable $`f__1`$ (see the previuos section), let us define the function with values $`f(q,y;x,f__1),0q1,y\mathrm{IR}`$, as the solution of the nonlinear antiparabolic equation $$_qf+\frac{1}{2}(f^{\prime \prime }+x(q)f^2)=0,$$ (23) with final condition $$f(1,y;x,f__1)=f__1(y)$$ (24) In (23), $`f^{}=_yf`$ and $`f^{\prime \prime }=_y^2f`$. With these definitions, the following representation theorem holds . ###### Theorem 4.1 There exists a nonempty hyper-surface $`\mathrm{\Sigma }_N(\omega ,f__1)`$ in $`𝒳`$ such that, for any $`x\mathrm{\Sigma }_N(\omega ,f__1)`$ and $`f`$ solution of (23,24), we have the following representation $$\psi _N(\omega ,f_1)=f(0,0;x,f_1).$$ (25) Any family of functional order parameters, $`x_ϵ`$, depending continuously in the $`L^^1`$ norm on the variable $`ϵ`$, $`0ϵ1`$, with $`x_00`$, and $`x__11`$, and nondecreasing in $`ϵ`$, must necessarily cross $`\mathrm{\Sigma }_N(\omega ,f__1)`$ for some value of the variable $`ϵ`$ (we say that $`\mathrm{\Sigma }_N(\omega ,f__1)`$ has the monotone intersection property). A similar representation holds also in the infinite volume limit. Of particular interest are those states $`\omega `$ such that the representation (25) holds with some $`x`$, depending on $`\omega `$, but independent on $`f_1`$, with some possible error vanishing in the limit $`N\mathrm{}`$. We call such states $`x`$-representable. Some examples of $`x`$-representable states are shown in . An attractive conjecture is that the Boltzmann state of mean field spin glass models is $`x`$-representable. Indeed, this must be the case if $`x`$ is the correct order parameter. We will refer to this as the tomographic Ansatz : in the $`𝒳`$ space the hyper-surfaces $`\{\mathrm{\Sigma }_{\mathrm{}}(\omega ,f_1),f_1F_1\}`$ have a common point $`x`$, which gives the physical content of the theory. By this Ansatz, we can express the full probability distribution of overlaps in terms of the functional order parameter. Let us state the following theorem, one of the main results of this paper, leaving the proof to the next sections. ###### Theorem 4.2 Let $`\omega `$ be an even state on the Ising spin variables $`\sigma _i`$, depending on the quenched variables $`J`$, and suppose it is $`x`$-representable, with $`x(0)=0`$ and $`x(1)=1`$. Then the following holds. 1. The probability distributions of overlaps among $`s=2,3`$ replicas are given in terms of the functional order parameter $`x`$ by the following expressions : $`P_{12}(q)P(q)`$ $`=`$ $`{\displaystyle \frac{d}{dq}}x(q),`$ (26) $`P_{12,23,13}(q_{12},q_{23},q_{13})`$ $`=`$ $`{\displaystyle \frac{1}{2}}x(q_{12})P(q_{12})\delta (q_{12}q_{23})\delta (q_{12}q_{13})+`$ (27) $`+{\displaystyle \frac{1}{2}}(P(q_{12})P(q_{23})\theta (q_{12}q_{23})\delta (q_{13}q_{23})+\mathrm{cyclic}\mathrm{perm}.).`$ 2. Assume in addition the hypothesis of Theorem 2.1. Then the overlap distribution is uniquely determined in terms of the functional order parameter $`x`$, and the $`s`$-replicas marginals (i.e., the distribution of overlaps among $`s`$ replicas) can be given explicitly for any $`s`$ (see section 6). We have used Dirac’s $`\delta `$ function and the step function $`\theta `$. Extension to regions of negative $`q`$’s is made by gauge symmetry, as shown in the next section. Equation (26) gives the physical meaning of the functional order parameter; equation (27) corresponds to ultrametricity of the overlap distribution, as is proven in the following. For other values of $`x(0)`$ and $`x(1)`$, slightly different results can be obtained. As is shown extensively in the next section, ultrametricity arises naturally as a consequence of the branching diffusion process underlying equation (23, 24). All results are in full agreement with those found in the frame of replica symmetry breaking method with Parisi Ansatz . ## 5 Asymptotic solution of the antiparabolic equation. The results (26,27) of Theorem 4.2 are obtained by equation ( 25), and the tomographic Ansatz. Both members of equation (25) are expressed as asymptotic series, which are then compared term by term. The first one is given by equation (18), the second is obtained in this section. Let us transform equation (23, 24) $$\{\begin{array}{c}_qf_q+\frac{1}{2}\left(f_q^{\prime \prime }+x_qf_{q}^{}{}_{}{}^{2}\right)=0\hfill \\ f(1,y;x,f__1)=f__1(y)\hfill \end{array}$$ into an equivalent form. When $`x(0)=0`$ and $`x(1)=1`$, satisfied by physical order parameters, it is convenient to make the substitution $$g_q(y)=\left[1\mathrm{exp}\left(x_qf_q(y)\right)\right]/x_q$$ (28) the $`x`$ and $`f__1`$ dependence of $`f`$ being understood. The resulting equation is $$\{\begin{array}{c}_qg_q+\frac{1}{2}g_q^{\prime \prime }=\rho _q\left[x_qg_q+(1x_qg_q)\mathrm{ln}(1x_qg_q)\right]/x_q^2\hfill \\ g(1,y;x,f__1)g__1(y)=1\mathrm{exp}\left(f__1(y)\right)\hfill \end{array}$$ (29) where $`x_qx(q)`$ and $`\rho _q\mathrm{d}x_q/\mathrm{d}q`$. Notice that the final condition for $`g`$ is equal to the function used in the expansion of $`\psi _N(\omega ,f_1)`$ (eq. 16), and that $`g(0,y)=f(0,y)`$. Let us re-write equation (29) in integral form: $$g_q=N_{1q}g_1+_q^1dq^{}\frac{\rho _q^{}}{x_q^{}^2}N_{q^{}q}\left[x_q^{}g_q^{}+(1x_q^{}g_q^{})\mathrm{ln}(1x_q^{}g_q^{})\right]$$ (30) as one can straightforwardly see by simple inspection. Here $`N_qN(q,y)=\mathrm{exp}\left(y^2/2q\right)/\sqrt{2\pi q}`$ is the usual heat kernel and the symbol $``$ is the convolution operation on $`y`$ variable <sup>1</sup><sup>1</sup>1$`(fg)(y)dy^{}f(yy^{})g(y^{})`$. Equation (30) can be handled by asymptotic expansion of the r.h.s. term under square brackets : $$g_q=N_{1q}g_1+\underset{i=2}{\overset{\mathrm{}}{}}\frac{1}{i(i1)}_q^1dq^{}\rho _q^{}x_q^{}^{(i2)}N_{q^{}q}\left[\left(g_q^{}\right)^i\right]$$ (31) We write the above equation in the “moment space” : let $`\eta _q`$ be the Fourier transform of $`N_qg_q`$ in the $`y`$ variable and $`\phi `$ that of $`N_1g__1`$. Thus we have, after simple algebraic manipulation $$\eta _z=\phi _z+\mathrm{F}_z\left[\eta \right]$$ (32) where $`z`$ is a collective variable for $`(q,k)`$, $`\phi _z\phi (k)`$ is the same function appearing in eq. (18), $`\mathrm{F}_z\left[\eta \right]`$ is a function of $`z`$ and a functional of $`\eta `$ $$\mathrm{F}_z\left[\eta \right]\underset{i=2}{\overset{\mathrm{}}{}}\frac{1}{i!}\widehat{\mathrm{O}}_z^{^{(i)}}[\eta ,\mathrm{},\eta ]$$ (33) and the $`\widehat{\mathrm{O}}_z^{^{(i)}}`$ are well defined multi-linear integral operators $$\begin{array}{cc}\hfill \widehat{\mathrm{O}}_z^{^{(i)}}[\phi __1,\mathrm{},\phi _i]& \\ \hfill (i2)!& \mathrm{d}^ik\delta (k_1+\mathrm{}+k_ik)\underset{a=1}{\overset{i}{}}\phi _a(k_a)\hfill \\ & _q^1dq^{}\rho _q^{}x_q^{}{}_{}{}^{i2}\mathrm{exp}(q^{}\underset{(a,b)}{}k_ak_b)\hfill \end{array}$$ Every term in the asymptotic expansion is well defined. Notice that the representation Theorem 4.1 can be rephrased as $$\psi _N(\omega ,f_1)=g(0,0)=dk\eta (0,k)$$ (34) and this is the form that we will use in the sequel. The $`i`$–th functional derivative of F$`{}_{z}{}^{}[\eta ]`$ w.r.t. $`\eta `$ calculated in zero gives the integral kernel of the $`\widehat{\mathrm{O}}_z^{^{(i)}}`$ operator. In particular $$\mathrm{F}_z\left[\eta \right]\|_{\eta =0}=0;\frac{\delta \mathrm{F}_z\left[\eta \right]}{\delta \eta _w}\|_{\eta =0}=0$$ (35) Replacing $`\eta =\mathrm{L}[\phi ]`$ in (32) we have $$\mathrm{L}_z\left[\phi \right]\phi _z+\mathrm{F}_z\left[\mathrm{L}[\phi ]\right],$$ (36) which defines iteratively the inverse functional $`L_z[.]`$ : $$\mathrm{L}_z\left[\phi \right]=\underset{s=1}{\overset{\mathrm{}}{}}\frac{1}{s!}\mathrm{d}^ik\underset{a=1}{\overset{s}{}}\phi (k_a)L_z^{(s)}(k_1,\mathrm{},k_s)$$ (37) It is easy to check that $$\mathrm{L}_z\left[\phi \right]\|_{\phi =0}=0;\frac{\delta \mathrm{L}_z\left[\phi \right]}{\delta \phi _w}\|_{\phi =0}=\delta _{zw}$$ (38) where $`\delta _{zw}`$ is the usual Dirac’s function, and that, by derivating (36) w.r.t. $`\phi `$, $$\frac{\delta \mathrm{L}_z\left[\phi \right]}{\delta \phi _w}\delta _{zw}+dw^{}\frac{\delta \mathrm{F}_z}{\delta \eta _w^{}}\left[\mathrm{L}[\phi ]\right]\frac{\delta \mathrm{L}_w^{}\left[\phi \right]}{\delta \phi _w}$$ (39) Subsequent functional derivatives w.r.t. $`\phi `$, calculated in $`\phi 0`$, and the properties (35) and (38) allow us to obtain straightforwardly all the integral kernels $`L_z^{(s)}(k_1,\mathrm{},k_s)`$ for any $`s`$, in terms of $`\widehat{\mathrm{O}}_z^{^{(i)}}`$ operators. We thus obtain $$dk\eta (0,k)=\underset{s=1}{\overset{\mathrm{}}{}}\frac{1}{s}\mathrm{d}^sk\underset{a=1}{\overset{s}{}}\phi (k_a)\frac{1}{(s1)!}\mathrm{L}_o^{(s)}(k_1,\mathrm{},k_s)$$ (40) where $`\mathrm{L}_o^{(s)}(k_1,\mathrm{},k_s)`$ are the integral kernels, calculated for $`q=0`$ and integrated on the overall delta dependence in the $`k`$ variable. They are of the form: $$\frac{1}{(s1)!}\mathrm{L}_o^{(s)}(k_1,\mathrm{},k_s)=\underset{(a,b)}{}\mathrm{dy}_{ab}\rho _s^{^{(+)}}(\{\mathrm{y}_{ab}\})\mathrm{exp}\left(_{(a,b)}k_ak_b\mathrm{y}_{ab}\right)$$ (41) where $`\rho _s^{^{(+)}}`$ has support on a subset of $`[0,1]^{s(s1)/2}`$. In appendix we report the explicit expressions of $`\rho _s^{^{(+)}}`$, for $`s=2,3,4`$ and the recipe to construct it for a generic $`s`$. By eq. (34) we can compare the asymptotic expansions in eq. (18) and eq. (40). As the function $`\phi (k)`$ is even, only the even part in the $`k`$’s of integral kernels from both sides can be equated. Let us define on $`[1,1]^{s(s1)/2}`$ the function $`\rho _s`$, extending by ’gauge symmetry’ the function $`\rho _s^{^{(+)}}`$ $`\rho _s(\{\mathrm{y}_{ab}\})`$ $`=`$ $`2^s{\displaystyle \underset{\{\epsilon \}}{}}\rho _s^{^{(+)}}(\{\epsilon _a\mathrm{y}_{ab}\epsilon _b\})`$ (42) $`=`$ $`2^{(s1)}{\displaystyle \underset{\{\epsilon \}:\epsilon _1=1}{}}\rho _s^{^{(+)}}(\{\epsilon _a\mathrm{y}_{ab}\epsilon _b\})`$ where the sums run over all the $`\epsilon _a=\pm 1,a=2,\mathrm{},s`$ and $`\rho _s^{^{(+)}}=0`$ if its argument is outside $`[0,1]^{s(s1)/2}`$; we have used the invariance $`\epsilon _a\epsilon _a`$ to fix $`\epsilon _1=1`$ in the second line. We finally have $$\mathrm{exp}\left(_{(a,b)}k_ak_b\mathrm{q}_{ab}\right)=\underset{(a,b)}{}\mathrm{dy}_{ab}\rho _s(\{\mathrm{y}_{ab}\})\mathrm{exp}\left(_{(a,b)}k_ak_b\mathrm{y}_{ab}\right)$$ (43) which proves part $`a`$ of Theorem 4.2. For $`s>3`$ the l.h.s. of eq. (43) does not correspond to the characteristic function of the overlap distribution, as the number of the $`k`$’s parameters is not sufficient, but to its restriction on a hyper-surface of dimension $`s`$. In the next section, assuming the hypothesis of Theorem 2.1, we show that the results obtained so far allow us to construct the full overlap distribution function. The resulting $`s=4`$ overlap distribution coincides with $`\rho _s`$. This is a strong indication that $`\rho _s`$ is the correct distribution also in the case of no additional interactions, as required by Theorem 2.1. For a generic $`s`$ the distribution $`\rho _s`$ has the ultrametric form $$\rho _s(\{\mathrm{y}_{ab}\})=\underset{i:A_i^sA^s}{}p_i\rho _s^{^{(i)}}(\{\mathrm{y}_{ab}\}\|A_i^s)$$ (44) Here $`A_i^s`$ are disjoint sets, made by portions of hyper-planes in $`[1,1]^{s(s1)/2}`$, with dimension $`\|A_i^s\|s1`$; $`A^s`$ is the union set; $`p_i`$ are positive numbers, which sum to one, and $`\rho _s^{^{(i)}}(\mathrm{}\|A_i^s)`$ are probability densities, whose supports are the sets $`A_i^s`$. The r.h.s. can thus be interpreted as a composite probability formula: $`p_i`$ is the probability that the ultrametric event $`A_i^s`$ happens and $`\rho _s^{^{(i)}}(\mathrm{}\|A_i^s)`$ is the overlap probability, conditioned to $`A_i^s`$. The events $`A_i^s`$ are disjoint and each $`\rho _s^{^{(i)}}`$ effectively depends only on at most $`s1`$ variables. ## 6 Ultrametric distributions. Consider the set $`\mathrm{\Phi }`$ of random variables $`q_{a,b}[1,\mathrm{\hspace{0.17em}1}]`$ : $$\mathrm{\Phi }=\{q_{a,b},(a,b)\phi C\},$$ (45) where $`C`$ is the set of couples $`(a,b)`$ of natural numbers $`a,b\mathrm{IN},a<b`$. To the set $`\phi `$ are then associated a probability space and the probability distribution $`P_\phi (\mathrm{\Phi })`$ on it. The distribution functions $`P_\phi `$ satisfy the consistency conditions $$P_{\phi ,\phi ^{}}(\mathrm{\Phi },\mathrm{\Phi }^{})\underset{\alpha \phi ^{}}{}dq_\alpha =P_\phi (\mathrm{\Phi }),$$ (46) for all disjoint sets $`\phi ,\phi ^{}C`$. In the following we will often write $`\{A,B\}AB`$ and $`ab(a,b)`$ when not ambiguous. Let us introduce the operator $`\mathrm{\Pi }_{l,m}`$ that, acting on $`\phi `$, permutates the indices $`l`$ and $`m`$. E.g.: $`\mathrm{\Pi }_{12}\{(1,2),(2,3)\}=\{(1,2),(1,3)\}.`$ Let $`\phi ^{}=\mathrm{\Pi }_{l,m}\phi `$, and denote by $`\mathrm{\Phi }^{}`$ the associated set of $`q`$’s. According to the symmetries of the $``$ average, we ask the probability measure $`P`$ to be gauge invariant, and symmetric under permutations of indices, in the following sense: $$P_\phi ^{}(\mathrm{\Phi }^{})=P_\phi (\mathrm{\Phi })=P_\phi (\mathrm{\Pi }_{l,m}\mathrm{\Phi }^{})\left(\mathrm{\Pi }_{l,m}P_\phi \right)(\mathrm{\Phi }^{}).$$ (47) This defines the operator $`\mathrm{\Pi }_{l,m}`$ on the space of distributions. Let us consider a very particular class of distributions for the overlaps between three replicas, i.e., the ultrametric distributions : $`P_{12,23,13}(q_{12},q_{23},q_{13})`$ $`=`$ $`B(q_{12},q_{23})\theta (q_{12}q_{23})\delta (q_{13}q_{23})`$ $`+`$ $`B(q_{23},q_{12})\theta (q_{23}q_{12})\delta (q_{13}q_{12})`$ $`+`$ $`B(q_{13},q_{23})\theta (q_{13}q_{23})\delta (q_{12}q_{23}),`$ where $`B`$ is a distribution. This simply states that among the three overlaps, two are equal and the third is greater or equal. From eq. (6), a simple application of symmetries and of the consistency conditions (46), leads to ###### Proposition 6.1 If the distribution $`P_{12,23,13}`$ has the form (6), then for any tern of replicas, $`(a,b,c)`$, the operator $`_{a,b,c}`$ is defined such that $$P_{ab,ac,\phi }=_{a,b,c}(P_{ab,\phi },P_{ac,\phi }),$$ (49) where $`\phi C,(b,c)\phi `$ and $`(a,b),(a,c)\phi `$. The operator $`_{a,b,c}`$ is defined through its values $`_{a,b,c}(P_{ab,\phi },P_{ac,\phi })(q_{ab},q_{ac};\mathrm{\Phi })=`$ (50) $`=`$ $`P_{ab,\phi }(q_{ab};\mathrm{\Phi })\left[\theta (q_{ab}q_{bc})\delta (q_{ac}q_{bc})+\theta (q_{bc}q_{ab})\delta (q_{ac}q_{ab})\right]+`$ $`+`$ $`P_{ac,\phi }(q_{ac};\mathrm{\Phi })\theta (q_{ac}q_{bc})\delta (q_{ab}q_{bc})`$ $``$ $`\delta (q_{ab}q_{bc})\delta (q_{ab}q_{ac}){\displaystyle P_{ac,\phi }(q_{ac};\mathrm{\Phi })\theta (q_{ac}q_{bc})𝑑q_{ac}}.`$ Note that when the set $`\phi `$ is symmetric under permutation of indices $`b,c`$, we can introduce the operator $`\stackrel{~}{}_{a,b,c}`$, $$P_{ac,\phi ^{}}=\stackrel{~}{}_{a,b,c}(P_\phi ^{})_{a,b,c}(P_\phi ^{},\mathrm{\Pi }_{b,c}P_\phi ^{}),$$ (51) where $`\phi ^{}\{(a,b),\phi \}`$. This property of the overlap distribution corresponds to ultrametricity. In fact, eq. (50) simply states that for any triangle of overlaps in a given set $`\stackrel{~}{\phi }`$, two overlaps are equal and the third is greater or equal. The proof of the theorem and of the subsequent lemma as well, does not depend on the nature of the $`q_,`$ variables, but only on symmetries and general properties of probability spaces. ###### Lemma 6.2 In the hypothesis of theorem 6.1 we can express the probability distribution of the overlaps between $`s+1`$ replicas in terms of the distribution of the overlaps between $`s`$ replicas and $`q_{1,s+1}`$ (for $`s3`$). The proof goes as follows. Given $`s3`$ and $`ls`$, we define the set $`\phi _{s,l}`$ by $$\phi _{s,l}\{(a,b),1a<bs\}\{(c,s+1),1cl\},$$ (52) such that the following simple relations hold: $`\{(l+1,s+1),\phi _{s,l}\}`$ $`=`$ $`\phi _{s,l+1},`$ (53) $`\phi _{s,s}`$ $`=`$ $`\phi _{s+1,0}.`$ (54) Applying formula (51), with $`l+1s`$, we have $$P_{\phi _{s,l+1}}=P_{(l+1,s+1),\phi _{s,l}}=\stackrel{~}{}_{s+1,1,l+1}(P_{\phi _{s,l}}).$$ (55) By iteration we have the thesis $$P_{\phi _{s+1,0}}=\stackrel{~}{}_{s+1,1,s}\mathrm{}\stackrel{~}{}_{s+1,1,2}(P_{\phi _{s,1}}).$$ (56) Moreover, by definition of conditional probability we have $$P_{\phi _{s,1}}=\stackrel{~}{P}_{(1,s+1)}P_{\phi _{s,0}},$$ (57) where $`\stackrel{~}{P}_{(1,s+1)}`$ is given by eq.(13). Therefore we have proven the following ###### Theorem 6.3 If the 3-replicas overlap distribution $`P_{12,23,13}`$ is ultrametric (i.e., of the form (6)), and in the limits of validity of theorem 2.1, the overlap distribution is uniquely determined in terms of $`P_{12}`$. The explicit form of the distributions of overlaps among $`s`$ replicas, for any $`s`$ (i.e., the $`s`$-replicas marginals $`P_{\phi _{s,0}}`$), can be calculated by repeated applications of eqs. (56,57). Since Theorem 4.2.$`a`$ proves the hypothesis of Theorem 6.3 in the case of mean field spin glass models, this completes the proof of its part b. The explicit construction (56, 57) clearly shows that ultrametricity and the Ghirlanda–Guerra relations can be considered as complementary in order to determine the full overlap distribution, in the sense that one can hold independently of the other, but together they determine explicitly the overlap measure in terms of the one-overlap distribution. Results of this kind were obtained by Parisi in the case $`s=3`$ . ## 7 Conclusions It has been shown how mean field disordered models can be successfully analyzed in a mathematically rigorous framework, with a simple Ansatz which is completely different from the Replica Simmetry Breaking Ansatz . In the S.–K. spin glasses case, the main features of the accepted physical solution – the Parisi solution – have been obtained. The method exploited, due to F. Guerra, is based on the cavity method and general theorems, and can therefore be applied to other disordered mean field models such as the multi-spin interaction spin glasses or neural networks. The functional order parameter $`x(q)`$ has been introduced in the S.–K. model. By the Ansatz that $`x`$ is indeed the correct order parameter, all physical observables have been expressed in terms of it. The physical interpretation of the functional order parameter (i.e. $`dx(q)/dq=P(q)`$) results, and ultrametricity of overlaps is derived as a natural consequence of a branching diffusion process. It has been shown by explicit construction that ultrametricity of the 3-replicas overlap distribution together with the Ghirlanda–Guerra relations determines the distribution of overlaps among $`s`$ replicas, for any $`s`$, in terms of $`P_{12}`$. ## Acknowledgments The authors wish to express their warmest thanks to F. Guerra, for fruitful suggestions and stimulating discussions. ## 8 Appendix We report the explicit expressions of $`\rho _s^{^{(+)}}(\{\mathrm{y}_{ab}\})`$ for $`s=2,3,4`$. For two replicas we have $$\rho _2^{^{(+)}}(\mathrm{y}_{12})=_0^1𝑑q\rho (q)\delta (\mathrm{y}_{12}q)$$ (58) for three replicas $`\rho _3^{^{(+)}}(\{\mathrm{y}_{ab}\})`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑q\rho (q)x(q){\displaystyle \underset{(a,b)G_3}{}}\delta (\mathrm{y}_{ab}q)+`$ $`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{\pi }{}}{\displaystyle {}_{}{}^{(3)}_{0}^{1}}𝑑q{\displaystyle _q^1}𝑑q^{}\rho (q)\rho (q^{})\delta (\mathrm{y}_{\pi _1\pi _2}q^{}){\displaystyle \underset{(a,b)G_3(\pi _1,\pi _2)}{}}\delta (\mathrm{y}_{ab}q)`$ and for four replicas $`\rho _4^{^{(+)}}(\{\mathrm{y}_{ab}\})=`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle _0^1}𝑑q\rho (q)x^2(q){\displaystyle \underset{(a,b)G_4}{}}\delta (\mathrm{y}_{ab}q)+`$ $`+`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \underset{\pi }{}}{\displaystyle {}_{}{}^{(6)}_{0}^{1}}𝑑q{\displaystyle _q^1}𝑑q^{}\rho (q)x(q)\rho (q^{})\delta (\mathrm{y}_{\pi _1\pi _2}q^{}){\displaystyle \underset{(a,b)G_4(\pi _1,\pi _2)}{}}\delta (\mathrm{y}_{ab}q)+`$ $`+`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \underset{\pi }{}}{\displaystyle {}_{}{}^{(4)}_{0}^{1}}dq{\displaystyle _q^1}dq^{}\rho (q)\rho (q^{})x(q^{}){\displaystyle \underset{(a,b)G_3(\pi _1,\pi _2,\pi _3)}{}}\delta (\mathrm{y}_{ab}q^{})\times `$ $`\times `$ $`{\displaystyle \underset{(a,b)G_4G_3(\pi _1,\pi _2,\pi _3)}{}}\delta (\mathrm{y}_{ab}q)+`$ $`+`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \underset{\pi }{}}{\displaystyle {}_{}{}^{(3)}_{0}^{1}}dq{\displaystyle _q^1}dq^{}{\displaystyle _q^1}dq^{\prime \prime }\rho (q)\rho (q^{})\rho (q^{\prime \prime })\delta (\mathrm{y}_{\pi _1\pi _2}q^{})\delta (\mathrm{y}_{\pi _3\pi _4}q^{\prime \prime })\times `$ $`\times `$ $`{\displaystyle \underset{(a,b)G_4\{(\pi _1,\pi _2),(\pi _3,\pi _4)\}}{}}\delta (\mathrm{y}_{ab}q)+`$ $`+`$ $`{\displaystyle \frac{1}{6}}{\displaystyle \underset{\pi }{}}{\displaystyle {}_{}{}^{(12)}_{0}^{1}}dq{\displaystyle _q^1}dq^{}{\displaystyle _q^{}^1}dq^{\prime \prime }\rho (q)\rho (q^{})\rho (q^{\prime \prime })\delta (\mathrm{y}_{\pi _1\pi _2}q^{\prime \prime })\times `$ $`\times `$ $`{\displaystyle \underset{(a,b)G_3(\pi _1,\pi _2,\pi _3)(\pi _1,\pi _2)}{}}\delta (\mathrm{y}_{ab}q^{}){\displaystyle \underset{(a,b)G_4G_3(\pi _1,\pi _2,\pi _3)}{}}\delta (\mathrm{y}_{ab}q)`$ Here $`G_r(i_1,\mathrm{},i_r)`$ is the complete graph with vertices $`(i_1,\mathrm{},i_r)\{1,\mathrm{},s\}`$ <sup>2</sup><sup>2</sup>2clearly $`G_sG_s(1,\mathrm{},s)=\phi _{s,0}`$; $`_\pi ^{(n)}`$ indicates the sum on all different $`n`$ permutations $`\pi `$ on $`G_r`$ vertices’ indexes, which render permutation invariant the associated measure. The numbers $`p_i`$, the probabilities of different ultrametric events, are obtained by normalizing the corresponding measures; counting together the permutations of variables they are, for three replicas, $`(1/4,3/4)`$ and for four replicas $`(1/9,1/6,2/9,1/6,1/3)`$. The recipe to construct $`\rho _s^{^{(+)}}(\{\mathrm{y}_{ab}\})`$ is based on the costruction of abstact trees with a root and $`s`$ “leaves”, which carry the indices $`\mathrm{y}_{ab}`$ The $`\rho _s^{^{(+)}}`$ is given by a sum on all such trees constructed by elementary branchings: each element in the sum is an integral on at most $`s1`$ variables of the weight $`w_T(.)`$ associated with the tree $`T`$. For $`T`$ given, $`w_T`$ is the product of the combinatorial factor $`[(s1)!]^1`$ times the weights of the branchings forming the tree <sup>3</sup><sup>3</sup>3a branching formed by an input and, say, $`i`$ outputs, has a weight $`w_i(q)=(i2)!\rho _qx_q^{i2}`$ and suitable $`\theta `$ and $`\delta `$ functions on the integral variables and the the output variables $`\mathrm{y}_{ab}`$, according to the tree structure. A simple way to deduce the number of structurally equivalent graphs, goes as follow: we use a scale transformation in (36) to obtain the generic term $`\mathrm{L}_z^{(s)}[\phi ,\mathrm{},\phi ]`$ in terms of $`\{\mathrm{L}_.^{(s^{})}[\phi ,\mathrm{},\phi ]\}`$, for $`1s^{}<s`$ in the expansion of $`\mathrm{L}_z\left[\phi \right]\mathrm{L}_z^{(s)}[\phi ,\mathrm{},\phi ]/s!`$ . Let $`\phi \lambda \phi `$ be this scale transformation: it is $`\mathrm{L}_z^^1\left[\phi \right]=\phi _z`$ and, for $`s2`$ $$\underset{s=2}{\overset{\mathrm{}}{}}\frac{\lambda ^s}{s!}\mathrm{L}_z^{(s)}[\phi ,\mathrm{},\phi ]=\underset{i=2}{\overset{\mathrm{}}{}}\frac{\lambda ^i}{i!}\widehat{\mathrm{O}}_z^{^{(i)}}[\underset{s__1=2}{\overset{\mathrm{}}{}}\frac{\lambda ^{s__11}}{s__1!}\mathrm{L}_.^{(s__1)}\left[\phi \right],\mathrm{},\underset{s_i=2}{\overset{\mathrm{}}{}}\frac{\lambda ^{s_i1}}{s_i!}\mathrm{L}_.^{(s_i)}\left[\phi \right]]$$ (61) By multilinearity of the operators, equating terms with equal powers of $`\lambda `$, we have $$\mathrm{L}_z^{(s)}[\phi ,\mathrm{},\phi ]=\underset{i=2}{\overset{s}{}}\underset{\{m_j\}}{}{}_{}{}^{}\left(\widehat{\mathrm{O}}_z^{^{(i)}}[(\mathrm{L}^{^{(1)}}\left[\phi \right])^{m__1},\mathrm{},(\mathrm{L}^{(s1)}\left[\phi \right])^{m_{s1}}]\right)_{(S)}^{}$$ (62) the sum $`_{\{m_j\}}^{}`$ is on all $`m_j0`$ with the bounds $`_jm_j=i`$ and $`_jjm_j=s`$; $`(\mathrm{L}_j[\phi ])^{m_j}`$ is briefly for $`m_j`$ repetitions of $`\mathrm{L}_j[\phi ]`$ operator as argument of $`\widehat{\mathrm{O}}_z^{^{(i)}}`$ and finally $`(S)`$ is the symmetrical factor in the $`\phi `$’s given by $$S=\frac{s!}{m_1!\mathrm{}m_{s1}!2!^{m_2}\mathrm{}(s1)!^{m_{s1}}}$$ (63) which counts all structurally equivalent graphs.
warning/0002/astro-ph0002097.html	ar5iv	text	# A phase-reference study of the quasar pair 1038+528A,B ## 1 Introduction The quasar pair 1038+528 A,B (Owen et al. owen78 (1978)) consists of two flat-spectrum radio sources, with redshifts 0.678 and 2.296 (Owen et al. owen80 (1980)), separated on the sky by only 33″. This system provides a unique opportunity to carry out high precision, relative astrometric studies using the full precision of VLBI relative phase measurements, since most sources of phase errors are common for the 2 sources (Marcaide & Shapiro marca83 (1983)). VLBI studies of the mas-scale structure of flat-spectrum quasars show that they typically have “core-jet” morphologies, consisting of a highly compact feature (the “core”) located at the base of an extended linear feature or line of lower brightness components (the “jet”). Both 1038+528 A and B exhibit such structures. In standard models of extragalactic radio sources, these radio-emitting features arise from a collimated beam of plasma which is ejected with a highly relativistic bulk velocity from a region close to a central massive object such as a black hole (see eg. Blandford & Königl blandford86 (1986)). Whilst jet features may correspond to shocks in the moving plasma, and can give rise to the observed “superluminal” component motions in some sources (Porcas porcas87 (1986)), the “core” emission is thought to arise from a more-or-less permanent location close to the origin of the beam, where the ambient conditions correspond to a transition from optically thick to optically thin emission at the observed frequency. Although the “core” position may thus be frequency-dependent, for a fixed observing frequency the core should provide a stable marker, anchored to the central mass of the quasar, whose location can be used to define a precise position for the object as a whole. Although short time-scale variations in physical conditions may cause small changes in the “core” location, over long time-scales it may be used to track any systematic proper motion of the quasar. The results from a near decade-long VLBI monitoring program on 1038+52A,B at $`\lambda `$ 3.6 and 13 cm (from 1981.2 to 1990.5) are reported by Rioja et al. (rioja97 (1997)), whose main conclusions can be summarized as follows: 1. There is no evidence of any relative proper motion between the quasars A and B. The uncertainties in the astrometric parameters result in an upper bound to any systematic relative motion between the cores of 10 $`\mu `$as yr<sup>-1</sup>, consistent with zero. 2. A compact feature within the jet of quasar B, chosen as the reference point for the structure, expands away from the core at a steady, slow rate of $`18\pm 5\mu `$as yr<sup>-1</sup>, corresponding to v=$`(0.8\pm 0.2)h^1`$ c for a Hubble constant, H<sub>0</sub> = 100 h km s<sup>-1</sup> Mpc<sup>-1</sup>; q$`{}_{0}{}^{}=0.5`$. These values are used here throughout. 3. The accuracy of the relative separation measurement is limited by noise and source structure, with estimated precisions of about $`50\mu `$as at $`\lambda `$ 3.6 cm at any epoch. 4. Confirmation of the consistently large offset (about 0.7 mas) between the positions of the peak of brightness (“core”) at $`\lambda `$ 3.6 and 13 cm in quasar A. New VLBI observations of this pair were made in November 1995 (1995.9) at $`\lambda `$ 2, 3.6 and 13 cm. In this paper we report on results from our analysis of the 3.6 cm observations and investigate the temporal evolution of the source structures and relative separation from all four epochs spanning $`15`$ years. Investigations of frequency-dependent source structure have also been made from a comparison of the astrometric measurements of the separations between A and B at all 3 wavelengths observed in 1995; these will be presented elsewhere (Rioja & Porcas in preparation). Our new observations are described in Sect. 2. In Sect. 3 we describe the data reduction and mapping techniques used, and in Sect. 4 an analysis of the measurements in the maps. In Sect. 5 we compare the astrometric results from these observations with those from previous epochs and analyse the changes in separation. Conclusions are presented in Sect. 6. ## 2 Observations The pair of radio sources 1038+528 A and B was observed with the NRAO Very Long Baseline Array (VLBA) on November 10, 1995, for a total of 13 hours, alternating every 13 minutes between observations in dual 3.6/13 cm mode and observations at 2 cm. The 100m telescope at Effelsberg was also included in the array for the 3.6/13 cm scans. The primary beamwidths of all the antennas were sufficiently large that both sources could be observed simultaneously at all wavelengths. Each 10 minute observation of 1038+52A,B was preceded by a 3 minute observation of the compact calibration source 0917+624, to monitor the behaviour of the array. All stations used VLBA terminals to record an aggregate of 64 MHz bandwidth for each scan, using 1-bit sampling, subdivided into 8 channels (mode 128-8-1). For the dual 3.6/13 cm scans, four 8-MHz channels were recorded for each band (2254.5–2286.5 MHz; 8404.5–8436.5 MHz), using RHC polarisation. At 2 cm, eight 8-MHz channels (15 331.5–15 395.5 MHz) were recorded in LHC polarisation. The correlation was made at the VLBA correlator in Socorro (New Mexico). As for previous epochs, two separate “passes” were needed, using different field centres for the two sources, to recover data for both the A and B quasars from the single observation. Output data sets were generated for the two sources, consisting of the visibility functions averaged to 2 s, with samples every 1 MHz in frequency across the bands. ## 3 Data reduction We used the NRAO AIPS package for the data reduction. We applied standard fringe-fitting, amplitude and phase (self-) calibration techniques and produced hybrid maps of each quasar. The astrometric analysis was done using two different mapping methods: a “standard” phase-referencing approach, transferring phase solutions from one quasar to the other (see e.g. Alef alef88 (1988); Beasley & Conway beasley95 (1995)) and a novel mapping method for astrometry of close pairs of sources, hybrid double mapping (HDM) (Porcas & Rioja porcas96 (1996)). Both routes preserve the signature of the relative separation of the source pair present in the calibrated phases. These analysis paths are described in Sects. 3.1 to 3.3 below. ### 3.1 Hybrid mapping in AIPS We applied standard VLBI hybrid mapping techniques in AIPS for the analysis of the observations of both quasars A and B. We used the information on system temperature, gain curves and telescope gains measured at the individual array elements, to calibrate the raw correlation coefficients. We used the AIPS task FRING to estimate residual antenna-based phases and phase derivatives (delay and rate) at intervals of a few minutes. It is important to realise that FRING is a global self-calibration algorithm, and performs an initial phase self-calibration also. We ran FRING on the A quasar data set, with a point-source input model. Anticipating our phase-referencing scheme (Sect. 3.2) we applied the antenna phase, delay and rate solutions from A to both the A and B data sets, and averaged them in time to 60 s, and over the total observed bandwidth of 32 MHz. After suitable editing of the data, we made hybrid maps of both quasars, using a number of iterations of a cycle including the mapping task MX and further phase self-calibration with CALIB. Fig. 1a and b show the hybrid maps for both sources at 3.6 cm in 1995.9. The maps are made using uniform weighting of the visibilities, a map cell size of 0.15 mas and a circular CLEAN restoring beam of 0.5 mas (these same mapping parameters are used throughout this work). The “dirty” beam has a central peak of 0.57 x 0.47 mas in PA -29°(PA = position angle, defined starting at North, increasing through East). The root-mean-square (rms) levels in the A and B maps, in regions away from the source structures (estimated using AIPS task IMSTAT) are 1.0 and 0.12 mJy/beam respectively, an indication that dynamic range considerations dominate over thermal noise in determining the map noise levels. ### 3.2 Phase referencing in AIPS In order to make an astrometric estimate of the separation between quasars A and B at this 4th epoch, we first used a ”conventional” phase-reference technique to make maps of the quasars which preserve the relative phase information. In practice this consists of using the antenna-based residual terms derived from the analysis of the data of one “reference” source (A), to calibrate the data from simultaneous observations of the other ”target” source (B). The reference quasar source structure must first be estimated from a hybrid map, and then fed back into the phase self-calibration process to produce estimates of the antenna-based residuals, free from contamination by source structure. Phase referencing techniques work under the assumption that the angular separation between the reference and target sources is smaller than the isoplanatic patch size (i.e. the effects of unmodelled perturbations, introduced by the propagation medium, on the observed phases of both sources are not very different) and that any instrumental terms are common. Geometric errors in the correlator model must also be negligible. Assuming that the antenna residuals have been “cleanly” estimated using the reference source data, the calibrated phases of the target source should be free from the errors mentioned above, but still retain the desired signature of the source structure and relative position contributions. The Fourier Transformation of the calibrated visibility function of the target source produces a “phase referenced” map. The offset of the brightness distribution from the centre of this map reflects any error in the assumed relative separation in the correlator model. If the reference source has a true “point” structure and is at the centre of its hybrid map, this offset will be equal to the error; more generally, one should also measure the offset of a reference point in the reference source map, and estimate the error in the source separation used in the correlator model from the difference between the target and reference source offsets. In general, the success of the phase-referencing technique is critically dependent on the angular separation of the target and reference sources. Simultaneous observation of the sources, as was possible here, significantly simplifies the procedure, eliminates the need for temporal interpolation, and reduces the propagation of errors introduced in the analysis. While random errors increase the noise level in the phase referenced map, systematic errors may bias the estimated angular separation. For our implementation of phase-referencing using AIPS, we chose to re-FRING the (calibrated) A data set, using our hybrid map of quasar A as an input model, and applied the adjustments to the antenna phase, delay and rate solutions to both the A and B data sets before re-averaging. We then made maps of both A and B using MX, performing no further phase self-calibration. These are our “phase-reference astrometry” (PRA) maps (shown in Fig. 2a and b) on which we performed astrometric measurements (see Sect. 4.2). Although the rms noise levels in the PRA maps are slightly higher than in the corresponding hybrid maps (2.0 mJy beam<sup>-1</sup> for A and 0.24 mJy beam<sup>-1</sup> for B), our procedure ensures that the A and B visibility functions from which they are derived have been calibrated identically. ### 3.3 New mapping method for astrometry of close source pairs While the conventional phase-referencing approach worked well for our November 1995 observations of 1038+52A and B, the method relies on making a good estimate of the antenna residuals from just one of the sources - the reference. We have devised an alternative method which extends the standard VLBI self-calibration procedure to work on both sources together, for cases where they have been observed simultaneously, and when either could be used as the reference (see Appendix A). The basis of the new method is to recognise that, since the visibility functions for both sources are corrupted by the same (antenna-based) phase and phase derivative errors, the sum of the two visibilities also suffers the same errors. We form the point-by-point sum of the two data sets, creating a new one which represents the visibility function of a “compound source” consisting of a superposition of the two structures, corrupted by the common antenna phase errors. If the source separation is close enough, the (summed) data as a function of the (averaged) uv-coordinates can be Fourier Transformed to form a map of the compound source structure, and (iterative) self-calibration in FRING or CALIB yields the antenna-based residuals. The advantage of this approach is that the antenna-based residuals are determined using both source structures simultaneously, and may thus reduce the chance that reference source structural phase terms contaminate the residuals. We term this process “Hybrid Double Mapping” (HDM); a detailed description is given in Porcas & Rioja (porcas96 (1996)). It is convenient to shift the source position in one of the data sets (by introducing artificial phase corrections) prior to the combination into a compound-source data set, to avoid superposition of the images in the map. The phase self-calibration steps which are then applied to the combined data set are identical to the case of a single source. In HDM the information on the angular separation between the sources is preserved in the process of self-calibration of the combined visibilities, and can be measured directly from the compound-image map; the relative positions between the individual source images in the compound map, taken together with any artificial position shift introduced, give the error in the assumed angular separation in the correlator model. In this approach one must be careful to use the same number of visibility measurements in each time interval from the two data sets, in order to avoid the predominance of data from a particular source. Fig. 3 shows the HDM map of quasars A and B in 1995.9 at $`\lambda `$ 3.6 cm; the B source is artificially offset by -4 mas in declination. The rms noise in the map is 0.82 mJy beam<sup>-1</sup> \- higher than that in the hybrid map of B but lower than in that of A. ## 4 Analysis of the maps ### 4.1 Source structures The 1995.9 hybrid maps of quasars A and B at $`\lambda `$ 3.6 cm (Fig. 1) show the core-jet structures typical of quasars at mas scales. They may be compared with maps from previous epochs given in Rioja et al. (rioja97 (1997)). The structure of A in the new map shows no major changes with respect to previous epochs. There is a prominent peak at the SW end of the structure (the “core”) and a jet extending in PA 15–25° containing at least two “knot” components (k1 and k2). The new map of B at 3.6 cm is qualitatively similar to those from previous epochs. It shows 2 point-like components separated by just under 2 mas in PA 127°. Spectral arguments support the identification of the NW component as a “core” (Marcaide & Shapiro marca84 (1984)); the SE component, corresponding to a knot in the jet, has been used as a reference component in previous astrometric studies. The separation between these 2 components in 1995.9 has increased, continuing the expansion along the axis of the source, as discovered from previous epochs of observations at this wavelength (Rioja et al. rioja97 (1997)). There is no trace of the third, extreme SE component, seen in maps of this source at 13 cm. This feature is evidently of lower surface brightness at 3.6 cm and is resolved out at the resolution of these observations. We used AIPS task IMSTAT to estimate total flux densities for quasars A and B (within windows surrounding the sources in the hybrid maps). The values are given in Table 1. Table 1 also lists the fluxes and relative positions of the most prominent features in the maps of A and B, obtained using task JMFIT to find parameters of elliptical Gaussian functions which best fit the various source sub-components. The formal errors from the fits, however, do not give realisitic values for the parameter uncertainties. The distribution of flux between the core and k1 in quasar A, and their relative separation, are quite uncertain, for example. ### 4.2 Estimating positions of reference features The astrometric measurement of a separation between two non-point sources must always refer to the measured positions of reference points within maps (or other representations) of the source structures. The selection of suitable reference points is crucial in monitoring programs, where the results from the analysis of a multi-epoch series of observations are compared. Ideally, a reference point should correspond to the peak of a strong, unresolved component, which is well separated from other radio emission within the source structure. For the 1995.9 epoch observations of 1038+528 A,B we selected the same reference features as those used for the analyses of previous observing epochs. These are the ”core” component for quasar A, and the prominent SE component for quasar B. These features are labelled with a cross in Fig. 2a and b. The core of A is indeed strong and compact, but has the disadvantage that it merges with knot k1. Although the SE component of B is no longer the strongest feature at 3.6 cm, it has always been strong at both 3.6 and 13 cm wavelengths, is reasonably compact and is easily distinguishable in maps made at longer wavelengths, thus facilitating spectral studies. Our astrometric analysis refers to the measured positions of the peaks of these components in A and B. We used the AIPS task MAXFIT to measure the position of these peaks in the PRA and HDM maps. MAXFIT defines the location of a peak in a given map region by fitting a quadratic function to the peak pixel value and those of the adjacent pixels. A comparison of this method of defining the peak position with that used for earlier epochs is described in the next section. ### 4.3 Position error analysis An analysis of errors presented in Rioja et al. (rioja97 (1997)) shows that the dominant uncertainty in the astrometric measurements of the separation between this close pair of quasars comes from the limited reproducibility of the reference point positions in the VLBI maps, from epoch to epoch. The magnitude of this effect is hard to quantify, however, since it depends on the nature of the source brightness distribution surrounding the reference point, and the method used to define the position of the peak, in addition to the resolution of the array and the signal-to-noise ratio of the peak in the map. A rough estimate of the error due solely to finite signal-to-noise in the maps is given by dividing the beam size by the ratio of the component peak to the rms noise level in the maps (see e.g. Thompson et al. thompson86 (1986)). This yields values of 3.3, 3.5 and 1.4 $`\mu `$as for the A and B PRA maps and the HDM map. These may be taken to represent lower limits to the reference point position errors; realistic errors will be larger, and will depend on the nature of the reference features and the manner in which the position is estimated. It is important to choose a definition of the reference point position such that it can be reproduced reliably from epoch to epoch, and is as independent as possible from the parameters used in making the map (e.g. cell size and beam width). The AIPS task JMFIT can be used to fit an elliptical Gaussian to a component in a CLEAN map, for example. However, the position of the peak of the Gaussian depends on how asymmetric the component brightness distribution is, and the area of the map to which the fit is restricted. MAXFIT fits just to the local maximum around the peak map value, and is thus less sensitive to the rest of the distribution. We have attempted to quantify some limits to reproducibility arising from the use of MAXFIT for defining the peak position in CLEAN maps. We investigated the effect of changing the true position of a point-like source with respect to the pixel sampling (here 3.3 pixels per beam) by offsetting the source position in 10 increments of 1/10 of a pixel in the visibility domain, mapping and CLEANing the new data sets, and estimating the new positions in the CLEAN maps using MAXFIT. The maximum discrepancy found between the values of the artificial offset and the shift derived by MAXFIT was 1/20 of a pixel. This corresponds to 8 $`\mu `$as in our 3.6 cm maps. For the analysis of previous observing epochs, the reference points were defined to be the centroid of the most prominent delta functions from which the CLEAN source map was derived (Rioja et al. rioja97 (1997)). We examined possible systematic differences resulting from these different definitions of reference points. One might expect the largest discrepancies to arise when the underlying source structure near the reference point is asymmetric, as in quasar A. We investigated such differences by determining “centroid” positions for both A and B reference components, using various criteria for excluding clean components from the calculation; this included the “25 percent of the value at the peak” threshold used for earlier epochs. For A the difference between this centroid position and the MAXFIT value was 0.12 pixel (18 $`\mu `$as). For B the difference was less than 0.1 pixel. These are probably the largest potential sources of error arising from using different methodologies at different epochs. Our use of two different mapping procedures - phase-reference mapping and HDM - also gives some insight into the size of position errors resulting from standard CLEAN + phase self-cal mapping algorithms. The differences between the separation estimates from the PRA maps and the HDM map are 27 and 28 $`\mu `$as in RA and Dec respectively. This would suggest that differences in the positions of peaks in maps reconstructed in different ways may vary at the 14 $`\mu `$as level. After considering the various possible effects which can limit the accuracy of postion estimates, we adopt a “conservative” value for the error in estimating the peak position in our $`\lambda `$ 3.6 cm maps, embracing all the effects detailed above, of $`18\mu `$as (this corresponds to a thirtieth of the CLEAN beam). The associated estimated error for a separation measurement between the two sources is $`25\mu `$as. ### 4.4 Astrometry results Table 2 lists the results of our astrometric measurements of reference point positions in the maps. They are presented as changes in measured separation between the reference points in A and B in 1995.9, with respect to their separation in 1981.3. The values given from the phase-reference technique correspond to the difference between the A and B reference feature position offsets in their respective PRA map. The values derived from HDM have been corrected for the artificial offset introduced before adding the A and B source visibilities. All these values have been corrected for a small error in the AIPS calculation of the u,v coordinates in the frequency-averaged data set. (Distances measured within the maps must be adjusted by a small correction factor of $`1\mathrm{\Delta }\nu (2\nu )^1=0.998`$.) Table 3 lists the coordinates of the reference source (A) adopted in the analysis and the measured coordinate separation between quasars A and B in 1995.9. ## 5 Comparison of astrometry at all 4 epochs In this section we make a comparison of the astrometric measurements from the series of 4 epochs of observations. Any increase of the temporal baseline in the program of monitoring the separation between A and B should result in a more precise identification of any systematic trends, with an improved elimination of random contributions. In Sect. 5.1 we justify comparing the astrometric values measured at the various epochs, even though non-identical observing, post-processing and analysis procedures were involved. In Sect. 5.2 we present the astrometric results from the 4 epochs. In Sects. 5.3, 5.4, 5.5 and 5.6 we present various analyses of these results, and attempt to quantify, or put upper-limits to, proper motions within and between the A and B quasars. ### 5.1 Comparison between the techniques used at different epochs Before attempting a comparison of the astrometric results from the 4 observing epochs, we need to show that any bias in the astrometric estimates introduced by the use of different procedures is small compared with other errors in the measurements for the individual epochs. The consistency between the results from previous epochs of observations has been exhaustively tested (Marcaide et al. marca94 (1994); Rioja et al. rioja97 (1997)). We outline here the largest changes involved in the fourth epoch, 1995.9, with respect to previous ones: 1. The observing array and frequency set-up used in the fourth epoch was different from previous epochs of observations (frequency range 8404.5 to 8436.5 MHz instead of 8402.99 to 8430.99 MHz at first 3 epochs). This results in a different coverage of the UV plane, leading to changes in the reconstruction of the source images. Investigations of such effects by Marcaide et al. (marca94 (1994)) show that the effect on the astrometric anaylsis is only a few $`\mu `$as. It is important to note that the observations at all 4 epochs have comparable resolutions and sample the same range of structural scales in the sources. 2. The processing of the fourth epoch was done using the VLBA correlator, which uses a theoretical model derived from CALC 8.2; we used AIPS to analyse the data with visibility phases residual to that model. For previous epochs the correlation was done at the MPIfR (Bonn) MK3 correlator, and an analysis of the data using total phases was made with VLBI3 (Robertson robertson75 (1975)). The differences between CALC 8.2 and the one implemented in VLBI3 propagate into changes of only 1-2 $`\mu `$as in the astrometric analysis of the 1038+528 A-B separation (Rioja rioja93 (1993), Rioja et al. rioja97 (1997)). This is because any such differences are “diluted” by the source separation expressed in radians - $`10^4`$ in the case of this very close source pair. 3. The values used in the analysis of previous epochs for Earth Orientation Parameters (EOP), stations and reference source coordinates were consistently derived from a single global solution provided by Goddard Space Flight Center (GSFC). For the correlation of the fourth epoch, the values used for EOP were derived from IERS solutions, and the station coordinates from USNO catalogs. We have made a comparison of the values derived for all the parameters at the 4 epochs from a single global solution from IERS (namely IERS eopc04), with the actual values used in the individual epoch analysis. The difference between the corresponding EOP values is always less than 4 mas. Such discrepancies propagate into errors in the relative position estimates at each epoch of only a few $`\mu `$as. 4. Our astrometric analysis in AIPS using a phase-referencing approach and HDM differs from the phase difference method used in VLBI3 analysis. Comparisons show that these procedures are equivalent (Porcas & Rioja porcas96 (1996); Thompson et al. thompson86 (1986)). Both involve the definition of reference points in source maps; uncertainties in the reference point positions (as described in Sect. 4.3) arise in the same way. 5. Finally, a minor VLBA correlator error (Romney priv. comm.) caused incorrect time labels to be attached to the visibility records, resulting in incorrect (u,v) values. The effect on the relative visibility phases for our source pair is small ($`0.004`$°) and can be neglected. The magnitudes of all of the effects reported in this section are much smaller than our estimate in Sect. 4.3 of the uncertainity in reproducing the reference point in the source, from epoch to epoch, and we are thus justified in comparing the astrometric results from all 4 epochs. ### 5.2 Astrometric separations at the 4 epochs The astrometric measurements of the separations between the reference points in A and B at $`\lambda `$ 3.6 cm from 4 epochs are presented in Fig. 4. It includes our new 1995.9 measurement and those from three earlier epochs, in 1981.3, 1983.4 and 1990.5, reported in Marcaide & Shapiro (marca84 (1984)), Marcaide et al. (1994) and Rioja et al. (rioja97 (1997)), respectively. The origin of the plot represents the separation at epoch 1. Changes with time in Fig. 4 represent the vector difference between any motions of the reference points in quasars A and B. The near-orthogonal nature of the source axes in 1038+52 A,B (along which one might expect any motion to occur) simplifies the interpretation of any trends seen. The new 1995.9 value follows the same steady trend towards the NW shown by the three previous epochs. Rioja et al. (rioja97 (1997)) interpreted this as an outward expansion of the reference component in quasar B at a rate of $`18\pm 5\mu `$as yr<sup>-1</sup>, and quoted an upper bound on any proper motion of quasar A of $`10\mu `$as yr<sup>-1</sup>. ### 5.3 Vector decomposition In this section we attempt to separate the individual contributions from the 2 quasars in the astrometric separation measurements presented in Fig. 4. We make no assumption about the stability of either component, but assume that any displacements of the A or B reference points from their positions at epoch 1 are along the corresponding source axis directions. This is a plausible assumption if the reference point coincides with a non-stationary component moving along a ballistic trajectory, or with the location of the peak of brightness within an active core or near the base of jet, where changes during episodes of activity are likely to occur along the jet direction. This approach is closely related to that used previously by Rioja et al. (rioja97 (1997)). For fixed assumed source axes for A and B, it results in a unique decomposition of the changes in the A-B separation into separate A and B displacements, from 1981 to 1995. It is clear that the dominant contribution to the separation changes seen in Fig. 4 comes from quasar B, in which the source axis is well defined by the 127°PA of the separation between core and reference components. For quasar A the source axis bends, from the inner “core” region (PA = 15°) to the outer jet components, and it is not so clear which direction should be chosen. In our analysis we tried a range of values for fixing the A source axis (0 to 45° in steps of 5°). For each, we calculated A and B reference-point displacements at epochs 2, 3 and 4 with respect to epoch 1. Then we performed a least-squares fit to the B displacements with time to estimate a linear expansion rate for the B reference feature along PA 127°. In Fig. 5 we plot the deconvolved B reference point displacements from the analysis with the A source axis fixed at PA 25°(the value adopted by Rioja et al. rioja97 (1997)). The fitted expansion rate is $`16.9\pm 0.6\mu `$as yr<sup>-1</sup>; the error and associated rms values take account of the small number of points and 2 degrees of freedom. This rate agrees well with the value of $`18\pm 5\mu `$as yr<sup>-1</sup> deduced by Rioja et al. (rioja97 (1997)). The rms residual from the fit (7 $`\mu `$as) is low, and vindicates our use of measurements derived from differing techniques for investigating the relative proper motion between A and B. ### 5.4 Structural evolution within 1038+528 B Our deconvolution analysis of the changes in separation measured between all 4 epochs supports the finding, previously proposed, that the B reference component moves along the source axis, away from the B core. In this section we make an independent determination of the separation rate between the core and reference component in B from measurements within the maps at the 4 epochs. Fig. 6 shows the separation between the core and reference component in B at the four epochs plotted against time. For epochs 1–3 we used the values given in Rioja et al. (rioja97 (1997)). For 1995.9 we used AIPS task UVFIT to estimate a separation from the B visibility data directly, in order to follow the methodology used for the other epochs as closely as possible; the value obtained was 1.895 mas. The slope from a least-squares fit corresponds to an expansion rate of $`13.0\pm 0.7\mu `$as yr<sup>-1</sup>. In the standard picture of extragalactic radio sources, the “core” is stationary, so this corresponds to an outward expansion of the reference component along PA 127°. The rms of the fit (8 $`\mu `$as) is again surprisingly low, implying typical errors in the separation measurements at each epoch (both within the B structure and between the reference points) of only about 10–12$`\mu `$as along the direction of the B source axis. This is considerably less than the estimate of position separation errors given in Sect. 4.3. ### 5.5 Relative proper motion The analysis presented in the previous sections demonstrate clearly that the chosen reference component within quasar B is unsuitable for use as a marker for tracing any relative proper motion between quasars A and B. The value of its expansion velocity derived in Sect. 5.4 appears to differ significantly from that deduced by vector-decomposition in Sect. 5.3. Although the difference between these estimates, if real, could be interpreted as motion of the core of B at a rate of $`4\mu `$as yr<sup>-1</sup>, this is not a conclusive result since differences of this order arise from choosing different values of PA for the motion in A in the vector decomposition method. A more suitable tracer of relative proper motion between the quasars is the variation of the separation between the cores of A and B. We have used the separations between the core and reference component measured in the B map at each epoch, and the astrometric separations between A and B, to calculate the separations between the A and B cores at each epoch; these are plotted in Fig. 7. The area occupied by the points defines an upper limit of $`10\mu `$as yr<sup>-1</sup> for any relative proper motion between the A and B cores, and hence between the quasars themselves, during the period of nearly 15 years for which the separation has been monitored with VLBI. The limit seems to be set by the relatively large deviation of the 1995.9 epoch point in the direction of the A source axis, presumably arising from the difficulty in defining the reference point at the A “core” from epoch to epoch. ### 5.6 Possible “core” motions ? Finally, we investigate any possible residual motions of the “cores” in A and B. The most likely causes of any such apparent motions are changes in the relative brightness or positions of features in the source structures at a resolution below that of the maps. One might expect that these, too, would produce effects predominantly along the source axis directions. We therefore used the vector deconvolution method on the plot of core-core separation with time to study displacements of the cores along their source axis directions. Fig. 8a and b show plots of the separated contributions from B and A, for an assumed A source axis PA 25°. The displacements for the B core seem to increase systematically. The fitted rate is $`3.8\pm 0.3\mu `$as yr<sup>-1</sup>, indicating a possible slow outward motion. The displacements for the A core do not seem to vary systematically - the fitted slope is $`5.5\pm 3.6\mu `$as yr<sup>-1</sup>. Here the scatter is considerably larger, reflecting both the difficulties of defining the reference point along the A core-jet axis, and also, perhaps, real “jitter” of the position of the peak due to variations in the “core” substructure. These plots indicate the level of stability of the individual core positions; the fits represent realistic upper limits to any possible systematic core motion in the A and B quasars along their source axis directions. ## 6 Conclusions The series of astrometric VLBI measurements of the separation between the quasar pair 1038+528 A and B, spanning nearly 15 years, provides excellent material for investigating the relative proper motion of two extragalactic radio sources and the positional stability of their cores. The changes measured in the separations between quasars A and B at 3.6 cm are dominated by the motion of the reference feature in quasar B. These astrometric results, and measurements in the hybrid maps of B are compatible with an expansion rate for the B reference component of 13–17$`\mu `$as yr<sup>-1</sup>. At a redshift of 2.296 this translates to an apparent transverse velocity of 0.55–0.70 c h<sup>-1</sup>. We note that this is an order of magnitude smaller than the more typical superluminal velocities seen in many quasars; it is a rare example of a subluminal velocity measured for a knot in a quasar jet. After correcting for the motion of the reference component in B, we can put a conservative upper bound to any relative proper motion between the quasars of $`10\mu `$as yr<sup>-1</sup>. Despite the increase in temporal baseline, this upper bound is no better than that given by Rioja et al. (rioja97 (1997)). Its value is related to the difficulty in reproducing a stable reference position along the A source axis near its “core”. Theories in which the redshifts of quasars do not indicate cosmological distances, and in which quasars are “local” and have high Doppler redshifts (e.g. Narlikar & Subramanian narlikar83 (1983)) are incompatible with our measured upper-limit to relative proper motion. Quasars at 100 Mpc distance moving at relativistic speeds would have proper motions of the order of $`600\mu `$as yr<sup>-1</sup>, nearly 2 orders of magnitude greater than our limit. Assuming cosmological distances, our limit corresponds to apparent transverse velocities of 0.43 c h<sup>-1</sup> and 0.22 c h<sup>-1</sup> at the redshifts of the B and A quasars, respectively. We have investigated the way in which the definition of reference points in a map may be only loosely “fixed” to the radio source structure, especially when the latter is strongly asymmetric. We have also developed an alternative analysis route - Hybrid Double Mapping - for imaging both sources of a close pair simultaneously, and at the same time preserving their relative astrometric information in a single map. The surprisingly low rms from the fits of linear expansion in quasar B, and the discrepancy between the two estimates ($`16.9\pm 0.6\mu `$as yr<sup>-1</sup> from the astrometric measurements and $`13.0\pm 0.7\mu `$as yr<sup>-1</sup> from the hybrid maps) are suggestive of (but do not prove) a residual motion of the core in quasar B; our decomposition along the source axis direction gives a fit of $`3.8\pm 0.3\mu `$as yr<sup>-1</sup>, corresponding to an apparent transverse velocity of 0.17 c h<sup>-1</sup>. If real, this might indicate a steady change in physical conditions at the base of the jet, or perhaps the emergence of a new knot component moving outwards with a velocity similar to the reference component, but as yet unresolved by our 0.5 mas beam. In this regard, it is interesting to note the slight extension of the core of quasar B in PA 132°given by the Gaussian model fit. The low rms derived from fits to the expansion of the B reference component indicate that we have been overly conservative in our estimate of $`18\mu `$as for the error in reference point positions. Errors at least 2 times smaller are implied, corresponding to a sixtieth of the beamwidth. It is interesting to note that such small errors are also implied in the work of Owsianik & Conway (owsianik98 (1998)), where the low scatter in the plot of expansion of the CSO source 0710+439 allows an expansion rate of $`14.1\pm 1.6\mu `$as yr<sup>-1</sup> to be determined. There are no obvious systematic motions within quasar A, but the “noise” in the estimates of position along its axis are much larger. This noise, along with any associated underlying changes in source substructure, provides a fundamental limit to estimates of any systematic core motion in A. Improvements on the estimates of (or upper bounds to) the relative motion between the quasars, or of the individual motion of the A core, will require a considerable increase in the temporal baseline of VLBI monitoring. ## Appendix A Hybrid Double Mappping (HDM) A.1 Principle of Hybrid Double Mapping The visibility function, $`V`$, measured at time $`t`$, on a baseline between antennas $`i`$ and $`j`$, is represented by a complex function with amplitude $`A`$, phase $`P`$: $$V(i,j)=Ae^{i[\varphi ]}(t)$$ For this analysis it is convenient to indentify 3 contributions to the visibility phase: $$\varphi (i,j)=\varphi _s(u,v)+\varphi _p(u,v)+\varphi _m(t)$$ where is due to source structure, evaluated w.r.t. a reference position for the source. is due to any offset of the true source position from the reference position. Both $`\varphi _s`$ and $`\varphi _p`$ are functions of the resolution coordinates, $`u`$ and $`v`$, at time $`t`$. is due to inaccuracies in the correlator model calculation of the interferometer geometry and the signal propagation delays in the ionosphere, troposphere and receiving system; it is an unknown function of time. This term can be represented by the difference of two ”antenna-based” phases, $`\theta _i`$ and $`\theta _j`$, since it can be related to the difference in signal arrival times at the two sites. (This analysis is a simplification which ignores possible ”non-closing” instrumental baseline phase terms arising from e.g. un-matched bandpasses and polarisation impurities.) $$\varphi (i,j)=\varphi _s(u,v)+\varphi _p(u,v)+\theta _i(t)\theta _j(t)$$ In conventional hybrid mapping, an iterative procedure is used to separate out the antenna-based phase terms from the ”source” terms; the latter must produce a consistent and physically plausible source structure after Fourier transformation of the corrected visibility: $$\underset{correctedvisibility}{\underset{}{Ae^{i[\varphi _s+\varphi _p]}(u,v)}}=V(i,j)\underset{antennaphaseterms}{\underset{}{e^{i[\theta _i(t)\theta _j(t)]}}}$$ However, the position offset term, $`\varphi _p`$, can also be expressed as a difference in wavefront arrival times at the 2 antennas and so it is also ”absorbed” in antenna phase terms $`\theta _i^{}`$, $`\theta _j^{}`$; the ”absolute” position information is lost: $$\underset{correctedvisibility}{\underset{}{Ae^{i[\varphi _s]}(u,v)}}=V(i,j)\underset{antennaphaseterms}{\underset{}{e^{i[\theta _i^{}(t)\theta _j^{}(t)]}}}$$ In Hybrid Double Mapping (HDM), the visibility functions of two sources observed simultaneously are added. For a close source pair, we make the same assumption as for conventional phase-referencing - that the model error phase terms are essentially the same for both sources. We make a further assumption that the $`u,v`$ coordinates are also essentially the same for both sources, for each baseline and time. The visibility sum, $`V^1+V^2`$, can then be re-written: $$\underset{correctedvisibilitysum}{\underset{}{(A^1e^{i[\varphi _s^1]}+A^2e^{i[\varphi _s^2+(\varphi _p^2\varphi _p^1)]})(u,v)}}=$$ $$=V^{sum}(i,j)\underset{antennaphaseterms}{\underset{}{e^{i[\theta _i^{}(t)\theta _j^{}(t)]}}}$$ This may be recognised as the visibility function of a ”composite” source consisting of the sum of the brightness distributions of sources 1 and 2, with antenna-based phase error terms $`\theta _i^{}`$, $`\theta _j^{}`$, as before. The HDM method consists of performing the normal hybrid mapping procedure with the visibility sum, resulting in the separation of the antenna-based errors, and a physically plausible map of the sum of the two source brightness distrubutions. An important point is that, whereas the origin of the map of the composite source is arbitrary (as it depends on the position of the starting model), the separation of the two source brightness distributions within the composite map (determined by $`\varphi _p^2\varphi _p^1`$) is fixed during the phase separation procedure, and is equal to the difference of the errors in the two source positions used for correlation. We call this the ”residual separation”. A.2 Practical aspects There are some practical aspects to be considered. If the source coordinates used in the correlator model are very precise, then the residual separation may be less than the interferometer beamwidth, and the two source distributions will lie on top of each other. In this case it is desireable to introduce an artificial position offset into one of the source visibility functions before forming the visibility sum, to ensure that the two source reference features are well separated in the HDM map. One should also arrange that the peak of one source does not lie on the sidelobe response of the other in the ”dirty” map, as this may degrade the CLEAN deconvolution process in the mapping step. Another important consideration is that the time-averaged samples of the summed visibility function contain equal contributions from both source visibility functions. When both sources are observed simultaneously this will normally be the case, except when different amounts of data are lost in the two separate correlator passes needed for the two source positions. It is important to edit the data sets carefully to fulfil this condition. The range of validity of the assumption that the $`(u,v)`$ coordinates for the two sources are the same depends on the ”dilution factor”, i.e. the reciprocal of the source separation, measured in radians. The $`(u,v)`$ value assigned to the summed visibility will be incorrect for either source by roughly 1 part in the dilution factor (roughly 1 in 6000 for 1038+528A,B). This is equivalent to having source visibility phase errors of this order, and thus limits the size of an HDM map to be less than the beamwidth times the dilution factor; the residual separation should be much smaller than this value. In the actual analysis used in this work, we first made a rough correction to the phase of the summed visibility of 1038+528 A + B, using the antenna phase and phase derivative errors from fringe-fitting 1038+528A using a point source model. However, there is no reason why one should not fringe-fit the summed visibility function directly. A.3 Applications The HDM method can in principle be applied whenever two (or more !) radio sources are observed simultaneously, but are correlated at separate field centres; however, they must be close enough so that the conditions of same $`(u,v)`$ coverage and same correlator model errors apply. The method uses the structures of BOTH sources simultaneously to separate out the antenna phase errors, as opposed to a single source in simple hybrid mapping. If both sources are strong (as with 1038+528 A and B), constraining the (single) antenna phase solutions with two structures should lead to a more rigorous and robust separation between the source and antenna phase terms. One field of application is in high resolution VLBI imaging of gravitational lens systems with wide image separations (e.g. images A and B of QSO 0957+561 with 6.1 arcsec separation) where preserving the necessary wide field-of-view from a single correlation may result in inconveniently large data sets. When one source is very weak, however, there is probably little to be gained over normal hybrid mapping. For relative astrometry studies (as described in this paper), the HDM method has some advantages over conventional phase-reference mapping and explicit phase-differencing methods. In phase-differencing astrometry, separate hybrid maps must be made of both sources to correct for source structural phase terms and the antenna phase errors are NOT constrained to be the same. Imperfect separation between source and antenna phase terms can increase the noise on the differenced phase, as well as lead to possible systematic errors. In phase-reference astrometry, only one of the sources is used to solve for the antenna phase terms; imperfect separation can lead to extra phase noise in the phase-referenced visibility of the ”target” source. In HDM we use both source structures simultaneously to separate the (common) antenna error phases from that of a single ”structure” in which the reference points of the two sources are spaced by the residual separation. When the separation between the two sources of a pair exceeds the telescope primary beamwidths, astrometric and phase-reference observations must involve switching between the sources, and the visibility phase of at least one of the sources must normally be interpolated in the observing gap. The condition that must be fulfilled for HDM to work in this case is that an equal number of observations of both sources must be added to form an average visibility function for the length of the ”solution interval” in the phase self-calibration step of HDM. This length is generally limited by the coherence time of the atmosphere, and would imply a very fast switching cycle in most cases. Another application for HDM could be in the the analysis of ”cluster-cluster” VLBI (see e.g. Rioja et al. rioja97b (1997)), in which two or more sources are observed simultaneously on VLBI baselines by using more than one telescope at each site. ###### Acknowledgements. We thank Dave Graham for help with the observations in Effelsberg, Tony Beasley for asistance at the VLBA, and Mark Reid for useful comments on the text. M.J.R. wishes to acknowledge support for this research by the European Union under contract CHGECT920011. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc.
warning/0002/cond-mat0002022.html	ar5iv	text	# Untitled Document 1 Introduction The travelling salesman problem (TSP) is a simple example of a multivariable combinatorial optmization problem and perhaps the most famous one. Given a certain set of cities and the intercity distance metric, a travelling salesman must find the shortest tour in which he visits all the cities and comes back to his starting point. It is a non-deterministic polynomial complete (NP- complete) problem \[1-3\]. In the standard formulation of TSP, we have $`N`$ number of cities distributed randomly on a continuum plane and we determine the average optimal travel distance per city $`\overline{l}_E`$ in the Euclidean metric (with $`\mathrm{\Delta }r_E=\sqrt{\mathrm{\Delta }x^2+\mathrm{\Delta }y^2}`$), or $`\overline{l}_M`$ in the Manhattan metric (with $`\mathrm{\Delta }r_M=\|\mathrm{\Delta }x\|+\|\mathrm{\Delta }y\|`$). Since the average distance per city scales (for fixed area) with the number of cities $`N`$ as $`1/\sqrt{N}`$, we find that the normalized travel distance per city $`\mathrm{\Omega }_E=\overline{l}_E\sqrt{N}`$ or $`\mathrm{\Omega }_M=\overline{l}_M\sqrt{N}`$ become the optimized constants and their values depend on the method used to optimize the travel distance. Extending the analytic estimates of the average nearest neighbour distances, in particular within a strip and varying the width of the strip to extremize (single parameter optimization approximation), one gets $`\frac{5}{8}<\mathrm{\Omega }_E<0.92`$ and $`\frac{5}{2\pi }<\mathrm{\Omega }_M<1.17`$ . Careful (scaling, etc.) analysis of the numerical results obtained so far indicates that $`\mathrm{\Omega }_E0.72`$ . Similar to many of the statistical physics problems redefined on the lattices, e.g., the statistics of self-avoiding walks on lattices (for investigating the linear polymer conformational statistics), the TSP can also be defined on randomly dilute lattices. The (percolation) cluster statistics of such dilute lattices is now extensively studied . The salesman’s optimized path on a dilute lattice is necessarily a self-avoiding one; for optimized tour the salesman cannot afford to visit any city more than once and obviously it is one where the path is non-intersecting. The statistics of self-avoiding walks on dilute lattices has also been studied quite a bit (see e.g., ). However, this knowledge is not sufficient to understand the TSP on similar lattices. The TSP on dilute lattices is a very intriguing one, but has not been studied intensively so far. The lattice version of the TSP was first studied by Chakrabarti . In the lattice version of the TSP, the $`N`$ cities are represented by randomly occupied lattice sites of a two- dimensional square lattice ($`L\times L`$), the fraction of sites occupied being $`p`$ ($`=N/L^2`$, the lattice occupation concentration). One must then find the shortest tour in which the salesman visits each city only once and comes back to its starting point. The average optimal travel distance in the Euclidean metric $`\overline{l}_E`$, and in the Manhattan metric $`\overline{l}_M`$, are functions of the lattice occupation concentration $`p`$ . We intend to study here the variation of the normalised travel distance per city, $`\mathrm{\Omega }_E=\overline{l}_E\sqrt{p}`$ and $`\mathrm{\Omega }_M=\overline{l}_M\sqrt{p}`$, with the lattice concentration $`p`$ for different system sizes. It is obvious that at $`p=1`$, all the self-avoiding walks passing through all the occupied sites will satisfy the requirements of TSP and $`\mathrm{\Omega }_E=1=\mathrm{\Omega }_M`$ (the distance between the neighbouring cities is equal to the unit lattice constant and the path between neighbouring sites makes discrete angles of of $`\pi /2`$ or its multiples with the Cartesian axes). The problem becomes nontrivial as $`p`$ decreases from unity: isolated occupied cities and branching configurations of occupied cities are found here with finite probabilities and self-avoiding walks through all the occupied cities, and only through the occupied cities, become impossible. As $`p`$ decreases from unity, the discreteness of the distance of the path connecting the two cities and of the angle which the path makes with the Cartesian axes, tend to disappear. The problem reduces to the standard TSP on the continuum in the $`p0`$ limit when all the continuous sets of distances and angles become possible. We study here the TSP on dilute lattice employing a computer algorithm which gives the exact optimized tours for small system sizes ($`N100`$) and near-optimal tours for bigger system sizes ($`100<N256`$). Our study indeed indicates that $`\mathrm{\Omega }_E`$ and $`\mathrm{\Omega }_M`$ vary with $`p`$ and $`\mathrm{\Omega }_E0.73`$ and $`\mathrm{\Omega }_M0.93`$ as $`p0`$. 2 Computer Simulation and Results We generate the randomly diluted lattice configurations following the standard Monte Carlo procedure for different system sizes. For each system size $`N`$, we vary the lattice size $`L`$ so that the lattice concentration $`p`$ varies. For each such lattice configuration, the optimum tour with open boundary conditions, is obtained with the help of the GNU tsp\_ solve developed using a branch and bound algorithm (see Fig. 1). It claims to give exact results for $`N100`$ and near-optimal solutions for $`100<N256`$. It may be noted that the program works essentially with the Euclidean distance. However there exists a geometric relationship between the Euclidean distance and the Manhattan distance. We may write $`l_E=_{i=1}^Nr_i`$, and $`l_M=_{i=1}^Nr_i\alpha _i`$, where $`r_i`$ is the magnitude of the Euclidean path vector between two neighbouring cities and $`r_i\alpha _i=r_i(\|\mathrm{sin}\theta _i\|+\|\mathrm{cos}\theta _i\|)`$ is the sum of the components of the Euclidean path projected along the Cartesian axes. Naturally, $`1\alpha _i\sqrt{2}`$. If $`l_E`$ corresponds to the shortest Euclidean path, then $`_{i=1}^Nr_i^{}>_{i=1}^Nr_i`$ , for any other path denoted by the primed set. If the optimized Euclidean path does not correspond to the optimized Manhattan path, then one will have $`_{i=1}^Nr_i^{}\alpha _i^{}<_{i=1}^Nr_i\alpha _i`$, where all the $`\alpha _i`$ and $`\alpha _i^{}`$ satisfy the previous bounds. Additionally, for random orientation of the Euclidean distance with respect to the Cartesian axes, $`\alpha _i=\alpha _i^{}=(2/\pi )_0^{\pi /2}(\mathrm{sin}\theta +\mathrm{cos}\theta )𝑑\theta =4/\pi `$. It seems, with all these constraints on $`\alpha `$’s and $`\alpha ^{}`$’s, it would be impossible to satisfy the above inequalities on $`r_i`$, and $`r_i\alpha _i`$. In fact, we checked for a set of $`50`$ random optimized Euclidean tours for small $`N`$ ($`<10`$), obtained using the algorithm, whether the optimized Manhattan tours correspond to different sequence (of visiting the cities), and did not find any. We believe that the optimized Euclidean tour necessarily corresponds to the optimized Manhattan tour. We then calculate $`l_E`$ and $`l_M`$ for each such optimized tour. At each lattice concentration $`p`$, we take about $`100`$ lattice configurations (about 150 configurations at some special points near $`p0`$) and then obtain the averages $`\overline{l}_E`$ and $`\overline{l}_M`$. We then determine $`\mathrm{\Omega }_E=\overline{l}_E\sqrt{p}`$ and $`\mathrm{\Omega }_M=\overline{l}_M\sqrt{p}`$ and study the variations of $`\mathrm{\Omega }_E`$ and $`\mathrm{\Omega }_M`$, and of the ratio $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ with $`p`$. We find that $`\mathrm{\Omega }_E`$ and $`\mathrm{\Omega }_M`$ both have variations starting from the exact result of unity for $`p=1`$ to the respective constants in the $`p0`$ limit. In fact we noted that although $`\mathrm{\Omega }_M`$ continuously decreases as $`p0`$, it remains close to unity for all values of $`p`$. We studied the numerical results for $`N=64,81,100,121,144,169,196,225\mathrm{and}256`$. The results for $`N=64\mathrm{and}100`$ have been shown in Figs. 2 and 3 respectively. We have studied the variations in the values of $`\mathrm{\Omega }_E`$ and $`\mathrm{\Omega }_M`$ against $`1/N`$ for $`p0`$, to extrapolate its value in the $`N\mathrm{}`$ limit. It appears that for the large $`N`$ limit (see Fig. 4), $`\mathrm{\Omega }_E(p0)`$ and $`\mathrm{\Omega }_M(p0)`$ eventually extrapolate to $`0.73\pm 0.01`$ (as in continuum TSP) and to $`0.93\pm 0.02`$, respectively. This result for $`\mathrm{\Omega }_E`$ (at $`p0`$) compares very well with the previous estimates . As $`p`$ changes from $`1`$ to $`0`$, the ratio $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ changes continuously from $`1`$ to about $`1.27(4/\pi )`$ (see Fig. 4), which is the average ratio of the Manhattan distance between two random points in a plane and the Euclidean distance between them . 3 Conclusions We note that the TSP on randomly diluted lattice is certainly a trivial problem when $`p=1`$ (lattice limit) as it reduces to the one-dimensional TSP (the connections in the optimal tour are between the nearest neighbours along the lattice). Here $`\mathrm{\Omega }_E(p)=\mathrm{\Omega }_M(p)=1`$. However, it is certainly NP- hard at the $`p0`$ (continuum) limit, where $`\mathrm{\Omega }_E0.73`$ and $`\mathrm{\Omega }_M0.93`$ (extrapolated for large system sizes $`N`$). We note that $`\mathrm{\Omega }_M`$ remains practically close to unity for all values of $`p<1`$. Our numerical results also suggest that $`\mathrm{\Omega }_M/\mathrm{\Omega }_E4/\pi `$ as $`p0`$. It is clear that the problem crosses from triviality (for $`p=1`$) to the NP- hard problem (for $`p0`$) at a certain value of $`p`$. We did not find any irregularity in the variation of $`\mathrm{\Omega }`$ at any $`p`$. A naive expectation might be that around the percolation point, beyond which the marginally connected lattice spanning path is snapped off , the $`\mathrm{\Omega }_E`$ or $`\mathrm{\Omega }_M`$ suffers some irregularity. The absence of any such irregularity can also be justified easily: the travelling salesman has to visit all the occupied lattice sites (cities), not necessarily those on the spanning cluster. Also, the TSP on dilute lattices has got to accomodate the same kind of frustration as the (compact) self-avoiding chains on dilute (percolating) lattices, although there the (collapsed) polymer is confined only to the spanning cluster. This indicates that the transition occurs either at $`p=1_{}`$ or at $`p=0_+`$. From the consideration of frustration for the TSP even at $`p=1_{}`$, it is almost certain that the transition occurs at $`p=1`$. However, this point requires further investigations. Acknowledgement : We are grateful to O. C. Martin and A. Percus for very useful comments and suggestions. References e-mail addresses : <sup>(1)</sup>anirban@cmp.saha.ernet.in <sup>(2)</sup>bikas@cmp.saha.ernet.in 1. M. R. Garey and D. S. Johnson, Computers and Intractability: A Guide to the Theory of NP- Completeness (Freeman; San Franscisco) (1979). 2. S. Kirkpatrick, C. D. Gelatt, Jr., and M. P. Vecchi, Science, 220, 671 (1983). 3. M. Mezard, G. Parisi and M. A. Virasoro, Spin Glass Theory and Beyond (World Scientific; Singapore) (1987). 4. J. Beardwood, J. H. Halton and J. M. Hammersley, Proc. Camb. Phil. Soc. 55, 299 (1959); R. S. Armour and J. A. Wheeler, Am. J. Phys. 51, 405 (1983). 5. A. Chakraborti and B. K. Chakrabarti, cond-mat/0001069 (2000). 6. A. Percus and O. C. Martin, Phys. Rev. Lett., 76, 1188 (1996). 7. D. Stauffer and A. Aharony, Introduction to Percolation Theory (Taylor and Francis; London) (1985). 8. K. Barat and B. K. Chakrabarti, Phys. Rep., 258, 377 (1995). 9. B. K. Chakrabarti, J. Phys. A: Math. Gen., 19, 1273 (1986). 10. D. Dhar, M. Barma, B. K. Chakrabarti and A. Tarapder, J. Phys. A: Math. Gen., 20, 5289 (1987); M. Ghosh, S. S. Manna and B. K. Chakrabarti, J. Phys. A: Math. Gen., 21, 1483 (1988); P. Sen and B. K. Chakrabarti, J. Phys. (Paris), 50, 255, 1581 (1989). 11. C. Hurtwitz, GNU tsp\_ solve, available at: http://www.cs.sunysb.edu/\̃\ algorith/implement/tsp/implement.shtml Figure captions Fig. 1 : A typical TSP for ($`N=)64`$ cities on a dilute lattice of size $`L=30`$. The cities are represented by black dots which are randomly occupied sites of the lattice with concentration $`p=N/L^20.07`$. The optimized Euclidean path is indicated. Fig. 2 : Plot of $`\mathrm{\Omega }_E`$, $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ against $`p`$ for $`N=64`$ cities, obtained using the optimization programs (exact). The error bars are due to configurational fluctuations. The extrapolated values of $`\mathrm{\Omega }_E`$, $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ are indicated by horizontal arrows on the y-axis. Fig. 3 : Plot of $`\mathrm{\Omega }_E`$, $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ against $`p`$ for $`N=100`$ cities, obtained using the optimization programs (exact). The error bars are due to configurational fluctuations. The extrapolated values of $`\mathrm{\Omega }_E`$, $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ are indicated by horizontal arrows on the y-axis. Fig. 4 : Plots of $`\mathrm{\Omega }_E`$ , $`\mathrm{\Omega }_M`$ and of $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ in the $`p0`$ limit, against $`1/N`$. The error bars are due to configurational fluctuations. The extrapolated value of $`\mathrm{\Omega }_E`$ , $`\mathrm{\Omega }_M`$ and $`\mathrm{\Omega }_M/\mathrm{\Omega }_E`$ in this $`p0`$ limit for $`N\mathrm{}`$ are indicated by horizontal arrows on the y-axis.
warning/0002/nucl-ex0002009.html	ar5iv	text	# New Limit for the Half-Life of 2K(2𝜈)-Capture Decay Mode of 78Kr. ## Abstract Features of data accumulated at 1817 hours in the experimental search for 2K(2$`\nu `$)-capture decay mode of <sup>78</sup>Kr are discussed. The new limit for this decay half-life is found to be T$`{}_{1/2}{}^{}2.3\times 10^{20}`$ yr. (90$`\%`$ C.L.). First result for 2K(2$`\nu `$)-capture decay mode of <sup>78</sup>Kr was presented in refs.-. The limit derived from the data collected at 254.2 hours was T$`{}_{1/2}{}^{}0.9\times 10^{20}`$ yr. (90$`\%`$ C.L.). The theoretical predictions for the half-life of <sup>78</sup>Kr(2e,2$`\nu `$)<sup>78</sup>Se capture are $`3.7\times 10^{21}`$ yr. , $`3.7\times 10^{22}`$ yr. and $`6.2\times 10^{23}`$ yr. . Corresponding values of half-life for 2K(2$`\nu `$)-capture decay mode are $`4.7\times 10^{21}`$ yr., $`4.7\times 10^{22}`$ yr. and $`7.9\times 10^{23}`$ yr. if one takes into account that 2K-electron capture part is 78.6$`\%`$ from a total number of 2$`e`$-captures for <sup>78</sup>Kr . The method used in refs.- allows to reach a level of sensitivity for a half-life up to 10<sup>22</sup> yr. and tests some of theoretical models. Results of the next step of measurement are presented. Measurements were performed with use of the multiwire wall-less proportional counter (MWPC) with a krypton sample enriched in <sup>78</sup>Kr. The main features of the counter and measurement conditions have been described in refs.-. The MWPC contents a central main counter (MC) and a surrounding it protection ring counter (RC) in the same body. A common anode wires signal ($`PAC`$) from RC and $`PC1`$ and $`PC2`$ signals from both ends of the MC anode are read out from MWPC. A scheme with signal read out from two sides of the MC anode allows to determine the event relative coordinate $`\beta `$ along the anode ($`\beta =100\times PC1/(PC1+PC2)`$) and to eliminate the events which don’t correspond to a selected working length. A shaping amplifier with 26 $`\mu `$s integration and differentiation shaping times was used for the amplification of the $`PC1`$ and $`PC2`$ pulses to have a good enough energy resolution. A parameter $`f=1000\times P12/(PC1+PC2)`$ was used to obtain an information about a pulse rise time and the pulse front features. Signals $`P12`$ are output pulses of the additional shaping amplifier which amplifies a sum signal ($`PC1+PC2`$) with the 1.5 $`\mu `$s shaping times. A value of parameter $`f`$ depends on the energy space distribution of event in the MC volume. The K-shell double vacancy of daughter <sup>78</sup>Se<sup>∗∗</sup> isotope appears as a result of <sup>78</sup>Kr 2K(2$`\nu `$) <sup>78</sup>Se capture. A total energy released is 2K$`{}_{ab}{}^{}=25.3`$ keV where K<sub>ab</sub> is a couple energy of K-electrons with the Se nuclear. One can obtain that Se<sup>∗∗</sup> summary probability to emit one or two characteristic X-rays is 0.837 in assumption that this double vacancy deexcitation is equivalent to sum of two single vacancies deexcitation. The characteristic X-ray ($`E_{K\alpha }11.2`$ keV, $`E_{K\beta }12.5`$ keV) has a sufficient long path length in a krypton. Two point-like energy releases with a total energy of 2K<sub>ab</sub> (total energy absorption peak) will appear if X-ray will absorb in the MC working volume. One part is the X-ray energy release and the second one is the release of Auger electron cascade energy accompany with the characteristic L-shell X-rays energy. The X-rays may leave the MC volume. One- or two-point event would be detected in this case (escape peak). All single electron background events such as Compton electrons or inner $`\beta `$-decay electrons will have one-point energy releases. A multi-point event pulse $`P12`$ would represent a sequence of short pulses with a different time overlap. A number of pulses in the burst corresponds to a number of the local regions where a total ionization distributed. An amplitude and duration of each pulse in a burst depend on a local track length, orientation and distance from the MC anode. The ADCs used to record the $`PAC`$, $`PC1`$, $`PC2`$ and $`P12`$ signals are triggered with the input pulse amplitude maximum. The $`P12`$ signal triggering will be done for the first amplitude maximum which corresponds to an energy released in the anode nearest local region. Peaks corresponding to one-point amplitudes $`P12`$ for a fixed event total energy appear in the event number distribution as a function of parameter $`f`$ ($`f`$-distribution). Events with energy released in MC only and the one in the MC and RC simultaneously named ”Type 1” and ”Type 2” events, respectively. A krypton enriched up to 94$`\%`$ in <sup>78</sup>Kr was used to search for the <sup>78</sup>Kr(2K,2$`\nu `$)<sup>78</sup>Se capture mode. It content an admixture of the natural $`\beta `$-radioactive <sup>85</sup>Kr (T$`{}_{1/2}{}^{}=10.7`$ yr, $`E_{\beta max}=670`$ keV) with the volume activity of 0.14 Bk/l. Measurements were done in the underground laboratory of the Baksan Neutrino Observatory of the Institute for Nuclear Research RAS (Moscow) at a depth of 4900 m w.e.. The MWPC was placed in the low background shield formed by 15 cm of lead, 8 cm of borated polyethylene, and 11 cm of copper. The own background of the MWPC filled up to 4.8 atm with pure xenon without radioactive contamination was measured preliminary. A background energy spectrum 1 collected at 973.9 h and a conveniently scaled spectrum 2 of a <sup>109</sup>Cd source ($`E_\gamma =88`$ keV) are shown on Fig. 1. The spectra consist of ($`PC1+PC2`$) signals from the type 1 events. There are a peak at $`E_\gamma =88`$ keV and the xenon escape peak at $`E=E_\gamma E_{XeK\alpha }=8828.9=58.2`$ keV on the curve 2. The peak at 88 kev is not symmetrical because of the radiation scattered in the counter wall. The energy resolution of the 88 keV $`\gamma `$-line is 13.7$`\%`$. The background spectrum has some features. The main peaks correspond to the energy values of 16, 35, 50, 68, 82 and 92 keV. In the energy regions $`35÷68`$ keV, the are initial peaks accompanied by the escape peaks. The background counting rate in the energy range $`20÷100`$ keV is 91 h<sup>-1</sup>. Energy spectra 1 and 2 of ($`PC1+PC2`$) signals from type 1 and type 2 events respectively for the <sup>109</sup>Cd calibration source and the krypton filling are shown on Fig.2$`a`$. One can see the 88 keV peak on the spectrum 1. This spectrum was multiplied by coefficient 0.5 for convenient comparison. The energy resolution of this peak is 10.8$`\%`$. The highest energy peak on the curve 2 is the krypton escape peak with the energy of $`E=E\gamma E_{KrK_\alpha }=8812.6=75.4`$ keV. It’s appearing in the type 2 events caused by an absorption in the RC of krypton characteristics radiation from the MC. The escape peak on the spectrum 1 is on the left slope of the total absorption peak. The source radiation scattered in the counter body wall lies in this region too. $`f`$-Distributions correspondent to this spectra are shown on Fig.2$`b`$ with the same scaling and notation. One can see a peak on the curve 1 with a maximum at $`f_1=166`$ which corresponds to two-point events from the total absorption peak when the KrK<sub>α</sub>-ray ionization collected first on the MC anode. If the photoelectron ionization collected first, the events have a peak with a maximum at $`f_2=920`$. Calculated $`f`$-values of these maximums should be equal to $`f_1=1000\times 12.6/88=143`$ and $`f_2=1000\times (8812.6)/88=878`$ at a calibration when the amplitude $`P12`$ equal to ($`PC1+PC2`$) one for the single-point events. Real values calculated from the experimental data differ slightly from the theoretical ones. This could be explained by nonideality of the experimental set up. A value of $`f`$-parameter depends on energy also because of it. The energy spectrum 2 on Fig.2$`a`$ consists of one-point events mainly and it’s $`f`$-distribution 2 has no multi-point peaks. The peaks at $`f=1015`$ and $`f=1239`$ are one-point peaks for the energy region higher than 16 keV and for the krypton 12.6 keV X-ray peak correspondingly. The $`f`$-distribution $`1`$ has a one-point peak with the maximum at $`f=1039`$. The type 1 event background energy spectrum of the MWPC with the krypton is shown on Fig.3$`a`$ (spectrum 1). It collected at 1817 h. A counting rate is 1506 h<sup>-1</sup> for the energy range $`20÷100`$ keV. Corresponding $`\beta `$\- and $`f`$-distributions are shown on Fig.3$`b`$ (curve 1) and Fig.3$`c`$ (curve 1). One can see peaks on the ends of the $`\beta `$-distribution which caused by events from a high energy part of the <sup>85</sup>Kr $`\beta `$-spectrum mainly collected in an ionization mode at the end effect correction anode bulges. To eliminate this background component, it is sufficient perform the $`\beta `$-selection of events in the range $`36\beta 58`$ (Fig.3$`b`$, curve 2). The energy spectrum 2 and $`f`$-distribution 2 correspond to this selection. One can see from the $`f`$-distribution 2 that the background events with $`f710`$ suppressed mainly. The energy spectrum of events with $`36\beta 58`$ and $`f710`$ is shown on Fig.3$`a`$ (spectrum 3). It’s shape repeats main features of the spectrum 1 on the Fig.1 with excluded escape peaks. It means that almost all one-point events from the <sup>85</sup>Kr $`\beta `$-spectrum excluded by the used selection. The curve 5 on the Fig.3$`c`$ shows roughly a shape and a place of a $`f`$-distribution waited for the <sup>78</sup>Kr(2K,2$`\nu `$)<sup>78</sup>Se multi-point events. The events with $`36\beta 58`$ and $`330f710`$ were used for the final analysis because of there is no any peak-like distortions of the residual energy spectrum $`4`$ (Fig.3$`a`$) in the region of interest for the such selection. A low energy part of this spectrum is shown on the Fig.4 (spectrum $`1`$). A sample of a shape and a place of the <sup>78</sup>Kr(2K,2$`\nu `$)<sup>78</sup>Se effect is shown as a spectrum 3 (Fig.4). An energy region $`25.3\pm 3.8`$ keV include 95$`\%`$ of the events. A background was fitted by using of points before and above this region (spectrum 2). The sum fitted background for the $`25.3\pm 3.8`$ keV was found to be 266. The experimental data sum is 262. The difference is $`4\pm 23`$ or $`19\pm 111`$ yr<sup>-1</sup>. Taking into account the efficiency of the events registration (0.22) and the effective counter length (0.6 of working length), we find that the limit of the half-life of <sup>78</sup>Kr with respect to the 2K(2$`\nu `$)-capture mode is T$`{}_{1/2}{}^{}2.3\times 10^{20}`$ yr (at a 90$`\%`$ C.L.). This work was supported by Russian Foundation for Basic Research (project nos 94-02-05954a and 97-02-16052) and in part by the International Science Foundation (grants nos RNT000 and RNT300).
warning/0002/astro-ph0002016.html	ar5iv	text	# Understanding NLR in Seyfert Galaxies: numerical simulation of jet-cloud interaction ## 1 Introduction Extensive HST emission-line imaging of Seyfert galaxies has for the first time resolved details of the structure of their Narrow Line Regions (NLR). In several cases cone-like morphologies have been revealed, similar in shape to - but of much smaller linear extent than - the Extended Narrow Line Regions (ENLR) seen in the lower resolution ground based images (Wilson & Tsvetanov, 1994 and references therein). In the standard model of the NLR, the UV emission of the nucleus is responsible of photoionizing the Interstellar Medium (ISM) of the host galaxy. These conical distributions of the ionized gas have been interpreted as a confirmation of the anisotropy of the nuclear radiation field which, in the framework of the unified scheme for Seyfert galaxies (e.g. Antonucci 1993) is caused by the shadowing of an obscuring circumnuclear torus. However, in galaxies with linear radio structures, the morphology of the emission-line region appears to be directly related to that of the radio emission. In particular, in Seyferts with radio jets (e.g. Mrk 3, Mrk 348, Mrk 6, Mrk 1066, ES0 428-G14), the NLR itself appears jet-like and is spatially coincident with the radio jet, while the emission-line region takes a different form when a radio lobe is present (e.g. Mrk 573, Mrk 78, NGC 3393): arc-like shells of emission, very reminiscent of bow-shocks, surround the leading edge of the lobes (Capetti et al. 1995a, 1995b, 1996; Falcke et al. 1996, 1998). This dichotomy in radio and emission-line morphology is reflected in their different scales: bow-shock structures cover several kiloparsecs, while the jet-like features extend only over a few hundred parsecs. The simplest interpretation of this radio-to-optical correspondence is that the radio emitting outflow creates an expanding and cooling gas halo. The compression induced by the outflow causes the line emission to be highly enhanced in the regions where the jet-cloud interactions occur. A clear confirmation of this scenario came recently from HST spectroscopy of Mrk 3 (Capetti et al. 1999): its NLR has velocity field characteristic of a cylindrical shell expanding at a rate of 1700 km/s. They interpreted this as the consequence of the rapid expansion of a hot gas cocoon surrounding the radio-jet, which compresses and accelerates the ambient gas. HST observations also provided evidence for spatial variations in the NLR ionization structure. In NGC 1068 the material located along the radio jet is in a much higher ionization state than its surroundings. This might suggest the presence of a local source of ionization which dominates over the nuclear radiation field (Capetti, Axon and Macchetto 1997; Axon et al. 1998). In other sources, too distant for such a detailed analysis, the radial variations of the ionization parameter are generally much flatter than expected from pure nuclear photoionization on the basis of the measured density gradients (Capetti et al. 1996, Allen et al. 1999) requiring again a local source of ionizing photons. An appealing possibility of interpreting these data is to invoke the ionizing effects of shocks, originated by jet-cloud interactions: if these shocks are fast enough (velocities $`>`$ a few hundred km s<sup>-1</sup>) the hot, shocked gas could produce a significant flux of ionizing photons (Sutherland, Bicknell and Dopita 1993; Dopita and Sutherland 1995, 1996). Direct evidence for this emission has been found by Axon, Capetti and Macchetto (1999) who showed that the radio-jets in the Seyfert 2 galaxies Mrk 348 and Mrk 3 are associated with an extended linear structure in UV and optical continua. In this picture, the radio-jet would not only determine the morphology of the NLR but is physically involved in its ionization. A radio imaging survey of the CfA sample of Seyfert (Kukula et al. 1995) shows that radio linear structures are present in a large fraction of sources (more than 50%) suggesting that such an interaction is likely to be a quite common phenomenon in this class of objects. The jet interaction with the external medium is clearly a complex physical problem which involves both a hydrodynamical study of the jet propagation as well as a detailed understanding of the microphysics of the induced shocks, which might also be magnetized, and of the radiative processes. In the framework of Seyfert galaxies this issue has been tackled by several authors (Dopita and Sutherland 1995, 1996, Evans et al. 1999, Wilson and Raymond 1999, Allen et al. 1999). Their focus is however mainly on the shocks properties with a very detailed treatment of the emission mechanisms, with simplifying assumptions about the hydrodynamics (e.g. plane parallel geometry, steady-state shock). The comparison with the observations is based on the emitted spectrum and in particular on diagnostic line-ratios, particularly with the aim of distinguishing the different signatures of nuclear versus local photoionization. In this paper we follow a complementary, albeit different, approach by studying in detail the jet hydrodynamics, while adopting a simplified treatment of the radiative processes, as we employ an equilibrium cooling function in an optically thin approximation. This approach allows us to compare the results of simulations with the observed properties of NLR, in particular their morphology, the expansion velocities and the characteristic values of gas density and temperature. More precisely we consider the interaction of the jet with an inhomogeneity in the external medium (cloud) and our aim is that of constraining the jet and cloud physical parameters for which it is possible to reproduce the observed conditions. In this way, in addition of getting a better understanding of the NLR physics, we can also obtain information on the jet properties from the NLR data. Moreover, we can calculate the fraction of the jet power converted in radiation by shocks, resulting from the interaction of the jet with the environment. We then get from the global dynamics a conversion efficiency from kinetic to radiative power and we can determine whether the jet itself, via shocks, can provide an in situ photoionization source for the NLR emitting material, as discussed above. Steffen et al. (1997a) have used a rather similar approach with the main difference that they considered the jet propagating into a uniform medium. It seems that in this situation it is impossible to reach the high densities typical of the NLR with jet-like emission (see discussion below) on which we will focus in the present paper. This is because, at low density, radiation is not efficient enough to give the needed compression factors. The conditions of the emitting material obtained by Steffen et al. seem to be appropriated for the case of the more extended (lobe-like) line emission structure. Steffen et al. (1997b) considered also jet-cloud interactions mainly from an analytical point of view. They found that when a jet interacts with a large number of clouds the most relevant effects on the NLR structure are due to the most massive clouds located along the jet path. This lead us to our choice for the geometry of the simulation in which the jet hits a single dense cloud. The paper is structured as follows: In the next section (Sect. 2), we describe the basic physical problem and the observational constraints, while the equations used and the method of solution are examined in Sect. 3 and 4; the results of simulations are discussed in Sect. 5; conclusions are drawn in Section 6. ## 2 Observational data and astrophysical scenario Observational data provide us with quite detailed information on the physical conditions of the narrow line emitting regions, in particular HST observations can now be used to determine the propertirties of individual NLR clouds: typically, densities are larger than $`10^3`$cm<sup>-3</sup>, temperatures are of the order of $`10^42\times 10^4`$ K, and velocities are $`3001000\mathrm{km}\mathrm{s}^1`$ (Caganoff et al. 1991, Kraemer, Ruiz and Crenshaw 1998, Ferruit et al. 1999, Axon et al. 1998, Capetti et al. 1999). These are the observational constraints that we try to match in our simulations. Results of simulations of a jet impinging on a uniform medium, with properties typical of the ISM, have shown that it is not possible to match, in this situation, the density values reported above (Steffen et al. 1997a, Rossi & Capetti 1998). We will therefore consider throughout the rest of the paper the case of a jet impinging on pre-existing inhomegeneities. We can identify such inhomogeneities with giant molecular clouds (GMCs), that typically populate spiral galaxies. These objects have typically mass $`10^510^6\mathrm{M}_{}`$, radius $`<100`$ pc, and temperature $`10`$ K (Blitz 1993). The resulting particle densities span from a few up to about hundred particles per cm<sup>3</sup>. A supersonic jet, of radius $`10`$ pc, that bore its way through the interstellar medium has a considerably good chance of impinging frontally upon a (much larger) GMC, and this is the case we will consider in our simulations. In any event, this latter case, i.e. the head-on collision with a large cloud, can be considered the most efficient case of interaction, for the compression, acceleration and heating of the NLR material. As discussed below the effects of the jet/cloud interaction last for a time considerably longer than the cloud crossing time. Moreover, the jet crosses the tenuous inter-cloud regions at a much higher speed than while in a cloud. We therefore expect that more than one cloud will be interacting at any given time and they will display simultaneously the different evolutionary stages of the interaction. ## 3 The physical problem We study the evolution of a cylindrical fluid jet impinging upon a cold heavy steady inhomogeneity, namely the cloud, in pressure equilibrium with the external medium. The relevant equations governing the jet evolution, for mass, momentum conservation, and radiative losses, are $$\frac{\rho }{t}+(\rho 𝒗)=0,$$ $`(1\mathrm{a})`$ $$\frac{\rho v_r}{t}+(\rho v_r𝒗)=\frac{p}{r},$$ $`(1\mathrm{b})`$ $$\frac{\rho v_z}{t}+(\rho v_z𝒗)=\frac{p}{z},$$ $`(1\mathrm{c})`$ $$\frac{E}{t}+(E𝒗)=p𝒗,$$ $`(1\mathrm{d})`$ where the fluid variables $`p`$, $`\rho `$, $`𝒗`$ and $`E`$ are, as customary, pressure, density, velocity, and thermal energy ($`p/(\mathrm{\Gamma }1)`$) respectively; $`\mathrm{\Gamma }`$ is the ratio of the specific heats; $``$ represents the radiative energy loss term (energy per unit volume per unit time, Raymond and Smith 1977). The jet occupies initially a cylinder of length $`L`$. The initial flow structure has the following form: $$v_z(r)=\{\begin{array}{cc}\frac{v_z(r=0)}{\mathrm{cosh}[(r/a)^m]}\hfill & zL\hfill \\ & \\ 0\hfill & z>L\hfill \end{array}$$ where $`m`$ is a ‘steepness’ parameter for the shear layer separating the jet from the external medium (see Fig. 1). The choice of separating the jet’s interior from the ambient medium with a smooth transition, instead of a sharp discontinuity, avoids numerical instabilities that can develop at the interface between the jet and the exteriors, especially at high Mach numbers. Regarding the cloud, we fix its initial density $`\rho _{\mathrm{cloud}}`$ and impose pressure equilibrium with respect to external medium; for simplicity we consider a steady cloud, with a thickness equal to the jet diameter. ## 4 The numerical scheme ### 4.1 Integration domain and boundary conditions Integration is performed in cylindrical geometry and the domain of integration ($`0zD`$, $`0rR`$) is covered by a grid of $`1020\times 704`$ grid points. The axis of the beam is taken coincident with the bottom boundary of the domain ($`r=0`$), where symmetric (for $`p`$, $`\rho `$ and $`v_z`$) or antisymmetric (for $`v_r`$) boundary conditions are assumed. At the top boundary ($`r=R`$) and right boundary ($`z=D`$) we choose free outflow conditions, imposing for every variable $`Q`$ null gradient ($`dQ/d(r,z)=0`$). The boundaries are placed as far as possible from the region of the jet where the most interesting evolutionary effects presumably take place by employing a nonuniform grid both in the longitudinal $`(z)`$ and the radial $`(r)`$ directions (Fig. 1, panel a)). In the radial direction the grid is uniform over the first 500 points and then the mesh size is increased assuming $`\mathrm{\Delta }r_{j+1}=1.015\mathrm{\Delta }r_j`$. The jet spans over 200 uniform meshes, while the external boundary is shifted to $`r=10a`$ where $`a`$ is the jet radius. As for the $`z`$–direction, we assume a constant fine grid in the central part of the domain, where the cloud is located, i.e. in a sub-domain of length $`40a`$, between the grid points 180 and 844; conversely, in the remaining part we consider an expanded grid increasing the mesh distance according to the scaling law $`\mathrm{\Delta }z_{j\pm 1}=1.015\mathrm{\Delta }z_j`$, where the minus sign applies in the first 180 grid points and the plus sign above grid point 844. ### 4.2 Integration method The basic equations (1a-d) have been integrated with a two-dimensional version of the Piecewise Parabolic Method (PPM) of Colella & Woodward (1984) (for a discussion of the main characteristics of this code and its merits for this kind of problems see Bodo et al. 1995). Radiative losses are dealt with the operator splitting technique, following which we split a single time step into two parts. In the first part, we advance the dynamical quantities, by using the adiabatic equations. In the second part we update the internal energy, keeping all the other variables constant, by taking into account radiative losses. ### 4.3 Physical parameters and Scaling The physical problem that we are approaching is quite complex, with three different interacting and radiating media, i.e. jet, ambient medium and cloud, each one described by its density, temperature, velocity and size. We note that in the adiabatic simulations of propagating jets, by normalizing to the jet density, sound speed, jet radius and sound crossing time over the jet radius, we are left with only two parameters, namely the density ratio between jet and external medium and the jet Mach number. The presence of radiation complicates the matter (Rossi et al. 1997), in fact temperature is not scale free, since the radiative loss function in Eq. (1c) explicitly depends on its physical value and in addition to the sound crossing time ($`t_{\mathrm{cr}}=a/c_\mathrm{s}`$), we have another typical time scale of the system, i.e. the radiative time scale, defined as $`t_{\mathrm{rad}}=p/[(\mathrm{\Gamma }1)]`$ which depends on the density of the medium. Therefore in this case one has to consider for each medium the value of density and temperature as independent parameters. In addition, as we already noticed we are now considering three media. When the jet passes through the cloud, the evolution of the compressed cloud material is completely different with respect to the case of two media, where the jet continues to push dense material at the head and it does not have any reaccelerations related to the passage from a denser to lighter medium. We would like to stress that the presence of a inhomogeneity is fundamental, in fact only in this case, as we show later, it is possible to reach the proper density for the emitting material. In conclusion, we must assign a large number of physical parameters for defining the initial conditions of our simulations. A thorough investigation of such huge parameter space is unfeasible; however, not all the parameters are equally important and some of them can be well constrained by observational considerations. As a first step we will then fix criteria to minimize the number of free parameters. Concerning the external (uniform) medium, we have to fix $`\rho _{\mathrm{ext}}`$ and $`T_{\mathrm{ext}}`$, having $`v_{\mathrm{ext}}=0`$. With respect to $`\rho _{\mathrm{ext}}`$, we can assume one particle per cubic centimeter, a value which we know to be appropriate to the interstellar medium of our Galaxy (Cox & Reynolds, 1987). In reference to $`T_{\mathrm{ext}}`$, again its choice is not so crucial, since the most important temperature for the emission processes is the shock temperature, depending mainly on jet velocity, in any case observations tell us that the external medium is completely ionized, which means temperatures larger than $`10^4`$ K, and we assume $`T_{\mathrm{ext}}=10^4`$K. The jet is physically described by its density $`\rho _{\mathrm{jet}}`$, temperature $`T_{\mathrm{jet}}`$, initial velocity $`v_{\mathrm{jet}}`$ and radius $`a`$. The radius $`a`$ can be chosen as our length unit in order to scale the other lengths in the system, and following the radio observational suggestions (Pedlar et al. 1993, Kukula et al. 1999) we consider it to be $`10`$ pc. Concerning the jet density we do not have any tightening constraint, so in a first approach, we take it equal to $`\rho _{\mathrm{ext}}`$. Relatively to $`T_{\mathrm{jet}}`$, looking at the loss function (Fig. 2) we can immediately realize that its initial value it is not so crucial, since cooling is fast and soon the jet temperature falls to $`10^4`$ K. Anyway $`T_{\mathrm{jet}}`$ in our simulations is taken to be $`10^6`$K. The jet velocity will be instead an important parameter of our simulations. Finally we consider the cloud, its density $`\rho _{\mathrm{cloud}}`$ is the parameter on which we will focus our investigation. $`T_{\mathrm{cloud}}`$ will be fixed by imposing pressure equilibrium with the external medium. Actually GMC’s are not required to be in pressure equilibrium since they might be autogravitating, however, as discussed for the external temperature, the exact value of $`T_{\mathrm{cloud}}`$ is not crucial for the results of the simulations. For simplicity we consider a steady cloud ($`v_{\mathrm{cloud}}=0`$). The cloud dimensions must lie in the range of GMCs, so we will fix the size longitudinal to the jet to 20 pc (i.e. $`2a`$), with a indefinitely large (with respect to $`a`$) transversal size. In summary, we have three control parameters, namely the initial cloud density, the initial jet velocity and the initial jet density, that we fix equal to one particle per cm<sup>-3</sup>. So we will investigate in details the effects of adopting different values for $`v_{\mathrm{jet}}`$ and $`\rho _{\mathrm{cloud}}`$. ## 5 Results We begin our discussion with a short general description of the complete evolution of the jet-cloud interaction, that can be summarized in three steps (see Fig. 3 for a visualization of the basic features of the three steps for the case $`\rho _{\mathrm{cloud}}=30\mathrm{cm}^3`$ and $`v_{\mathrm{jet}}=6500\mathrm{km}\mathrm{s}^1`$): $``$ The jet hits the cloud, forming a strong shock, the post-shock region becomes hot and blows up, because of its increasing pressure; the jet material is conveyed in a back-flow that squeezes the jet itself. During this process the cloud material is compressed and heated by the shock, at the head the temperature is very high ($`>10^8`$ K), while on the jet sides it is lower $`10^7`$ K, so that it can cool down, to reach the observed line emission conditions. It is in this region, forming a layer around the jet, that the narrow line emission can originate. Our analysis will therefore concentrate on the properties of this region. During this first phase, in which the jet crosses the cloud, the layer is accelerated by the strong inside pressure and cools down, its density thereby increases (see the leftmost panels in Fig. 3). $``$ The second phase begins when the jet is completely out of the cloud, the compressed emitting material reaches a quasi-steady state, during which the emission is almost constant, the inside pressure begins to decay, but the emitting layer is still accelerated. From Fig. 3 (central panels), we can see that the material in the layer has been compressed, its maximum density has increased, while its temperature has decreased. The maximum density is found now at temperatures around $`10^4`$ K, and its velocity has also increased. $``$ In the third phase, the inside pressure has decayed and the emitting layer begins to slow down, the jet flows freely through the cloud and also the emission decreases, eventually disappearing. From Fig. 3 (rightmost panels), we see a decrease in density and velocity, while almost all the layer is found at $`5\times 10^3<T<10^4`$ K. The efficient formation of the line emitting region will therefore depend on the efficiency of radiation during the jet crossing of the cloud. We will then introduce two typical timescales, the cloud crossing time and the radiative time, whose ratio will be a fundamental parameter for determining the evolution of the narrow line emitting layer. Following analytical treatments of the jet-ambient interaction we define the cloud crossing time as $`t_{\mathrm{cc}}=d(1+\sqrt{\rho _{\mathrm{cloud}}/\rho _{\mathrm{jet}}})/v_{\mathrm{jet}}`$, where we have assumed, for the jet head velocity in the cloud, the steady velocity obtained from the 1-D momentum balance in a medium with $`\rho =\rho _{\mathrm{cloud}}`$ (see, e.g., Cioffi and Blondin 1992, Norman et al. 1982). In this way we are actually overestimating the crossing time, since our situation is not steady, however this value is sufficiently accurate for our purposes. Concerning the radiative time, its definition is given in Section 4.3, however, we must notice that for its evaluation we have to assume a value for the temperature, in the following considerations we have taken $`T=10^7`$ K, that is the average of the typical post-shock temperature in the region of our interest, this choice is properly done for all jets with $`v_{\mathrm{jet}}=6500\mathrm{km}\mathrm{s}^1`$ and $`v_{\mathrm{jet}}=32500\mathrm{km}\mathrm{s}^1`$, while it is overestimated for the low velocity cases, that means that $`t_{\mathrm{rad}}`$ for those cases are shorter than the real ones. We have then defined $`\tau t_{\mathrm{cc}}/t_{\mathrm{rad}}`$ as the ratio between crossing and radiative time scales and this, as said before, is an important parameter for the interpretation of the results. As a first step in our analysis, we have performed an exploration of the parameter space. As discussed before, we reduced our parameters to $`\rho _{\mathrm{cloud}}`$ and the initial $`v_{\mathrm{jet}}`$. In Table 1 we report, for each pair of their values, typical values of density, expansion velocity and temperature of the emitting material and the value of $`\tau `$. The density is the median value of density distribution weighted on the emissivity function (that is proportional to $`\rho ^2`$), while velocity and temperature are those corresponding to this density value. All the quantities are evaluated at $`2t_{\mathrm{cc}}`$, this choice is due to the fact that during this period the expansion velocity of the emitting material increases rapidly reaching a maximum and then decreases monotonically, so that, if the expansion velocity does not match the observational constraint within this time, it never will, and the case will not be of interest for our analysis. Radiation must therefore act efficiently during this time, in order to create the needed conditions for radiation, and this poses a lower limit on the value of $`\tau `$. On the basis of the values reported in this table we choose the most promising cases for our investigation. Considering the first column we can immediately realize that jets at low velocity cannot reach conditions comparable to those observed. The values of $`\tau `$ for these simulations are high, meaning that radiation is very efficient. On the other hand, the jet momentum is low and cannot drive the emitting material at high velocities. For the case $`\rho _{\mathrm{cloud}}=30\mathrm{cm}^3`$ we have, in fact, high densities in accord with the high value of $`\tau `$, but very low velocities. For this reason we did not perform simulations for the other two cases of higher density, since jets would produce stronger and cooler compression practically at rest, very far from the observational scenario. Looking at the high velocity case, we see that, in the case of small inhomogeneities, $`\tau `$ has a very low value and, therefore, radiation is inefficient. The jet is very energetic and sweeps the cloud, before radiation becomes effective and so it does not form any condensation (the velocity reported for this case is therefore meaningless). Increasing the cloud density, we increase also the value of $`\tau `$: the maximum density increases, but it is still quite low. Only for the high density cloud ($`\tau =0.5`$), we get values of density and velocity in agreement with observations. Regarding the intermediate velocity, the values of $`\tau `$ are $`>0.3`$: radiation is efficient and thus the emitting layer can reach sufficiently high densities. Only in the lighter cloud case, however, the velocity is comparable to the observed values. From this exploration of the parameter space we can conclude that the observed conditions can be matched only for a narrow range of parameters and that the properties of the emitting layer depend essentially only on one parameter, the ratio between the radiative timescale and the cloud crossing timescale $`\tau `$. For low values of $`\tau `$ ($`\tau <0.3`$), radiation is inefficient and the densities in the layer are too low. For higher values of $`\tau `$ ($`\tau >0.55`$) we find, on the other hand, that the velocity of the emitting layer becomes too small. This is because the cloud density is high and the jet momentum flux is too small to impart to it a large enough velocity. Only for a narrow range of values of $`\tau `$ we can match the observed conditions and, in Table 2, we have translated these limits into limits on velocity range at different cloud densities. ### 5.1 Case of $`\rho _{\mathrm{cloud}}=\mathrm{\hspace{0.17em}120}\mathrm{cm}^1,v_{\mathrm{jet}}=\mathrm{\hspace{0.17em}32500}\mathrm{km}\mathrm{s}^1`$ In this subsection we will discuss in more details the case that best matches the observational scenario. We begin our discussion showing, Fig. 4, a gray-scale image with a snapshot of the density distribution at $`5t_{\mathrm{cc}}`$ and three small panels showing enlargements of the region of interaction between jet and cloud referred to density, temperature and the expansion velocity of emitting gas. The proper physical condition for emission are reached in a thin layer of compressed cloud material, whose width and mass grow in time as the shocked cloud material cools down. The detailed physical properties of this line emitting region are reported in Fig. 5, where we have represented the behavior of density, temperature and velocity along radial cuts through this layer. We note that the proper conditions are matched in a layer of width $`<\mathrm{\hspace{0.17em}2}`$ pc. How the properties of the material contained in this thin layer compare with the physical conditions of gas of the NLR? To answer the question we plot in Fig. 6 the temporal behavior of the mean expansion velocity (panel a) and mass (panel b) of the emitting material shell at two different density limits. We see that from the time when the jet touches the cloud until $`2t_{\mathrm{cc}}`$, when a strong interaction between the jet head and the cloud takes place, the cloud material is accelerated and the quantity of emitting material increases; after this interval the jet flows, essentially freely, across the cloud without any further acceleration of the compressed material shell and the accelerated cloud material slows down monotonically. Notice that the mean expansion velocity, relative to an observer, lies, for the denser material, in the range $`6001200\mathrm{km}\mathrm{s}^1`$ (since one must consider twice the mean expansion velocity), that is in good agreement with the velocity deduced by the line widths detected. Looking more in detail at the emitting mass, we see that its growth begins some time after the jet has initiated to drill its way into the cloud, and this delay corresponds to the cooling time of the shocked material. We also note that, after $`t=2t_{\mathrm{cc}}`$, the jet continues to sweep out material laterally at a pace that is higher for the lighter material, the total mass exceeds $`3\times 10^4/\mathrm{M}_{}`$ at $`t=35,000`$ ys and this would correspond to an $`H_\beta `$ luminosity of $`2\times 10^{40}`$erg s<sup>-1</sup> which, considering also the possibility of having simultaneuously several active clouds, is consistent with the observed values. As the interaction is effective over a timescale much longer than $`t=t_{\mathrm{cc}}`$ the jet will quickly propagate into the low density inter-cloud medium and it will reach any other cloud lying on its path. Thus more than one cloud will be effectively interacting with the jet at any time. Each will display a behaviour typical of its evolutionary stage and the total emitting mass must be considered as the total over all clouds. Furthermore, this will naturally reproduce the jet-like morphology of the NLR. ### 5.2 Energetics As discussed in the Introduction, the source of ionization of the NLR is still matter of debate. While the NLR gas is certainly illuminated by the nuclear source, its interaction with the radio jet also produces regions of high temperature and density which radiates ionizing photons. In this paragraph we derive the conversion efficiency of the jet kinetic power into energy radiated in ionizing photons. To estimate the ionizing energy flux we integrated radiative losses over all regions where $`T>10^5`$ K as above this temperature most of them correspond to the production of photons with energy higher than the hydrogen ionization threshold. In Table 3 we summarize our results reporting the kinetic power referred to the three different velocities ($`P_{\mathrm{kin}}=\rho v_{\mathrm{jet}}^3A`$, where $`\rho `$, $`v_{\mathrm{jet}}`$ and $`A`$ are respectively the density, the velocity and the transverse section of the jet) and the conversion efficiency at peak and after 2 $`t_{cc}`$ for all the cases considered. The peak efficiency reaches in one case a value as high as 10 % but it is usually 0.1 - 2 %. However, over the interaction, the typical value of $`\eta `$ (well represented by its value after 2 $`t_{cc}`$) is much lower $`\eta 10^45.\times 10^3`$. Faster jets have lower efficiency than slower jets and this conspires in producing a very similar amount of energy radiated in ionizing photons, $`10^{40}\mathrm{erg}\mathrm{s}^1`$, in all cases. In Seyfert galaxies $`L_{H\beta }10^{39}10^{42}\mathrm{erg}\mathrm{s}^1`$ (Koski 1978). The minimum ionizing photon luminosity required to produce a given line emission luminosity corresponds to the limiting case in which all ionizing photons are absorbed and all photons have an energy very close to the hydrogen ionization threshold $`\nu _{\mathrm{ion}}`$. In this situation $`L_{\mathrm{ion},\mathrm{min}}=\frac{1}{p_{H\beta }}\frac{\nu _{\mathrm{ion}}}{\nu _{H\beta }}L_{H\beta }50L_{H\beta }`$ where $`p_{H\beta }0.1`$ is the probability that any recombination will result in the emission of an H$`\beta `$ photon. It appears that, even in this most favourable scenario, the radiation produced in shocks can only represent a small fraction of the overall ionization budget of the NLR, particularly as sources with high radio luminosity (in which usually radio-jets are found) also have the highest line luminosity (e.g. Whittle 1985). Nonetheless, in the most promising case examined above ($`\rho _{\mathrm{cloud}}=120\mathrm{cm}^3`$ and $`v_{\mathrm{jet}}=32500\mathrm{km}\mathrm{s}^1`$) at the peak of the conversion efficiency the radiated energy is $`3\times 10^{41}\mathrm{erg}\mathrm{s}^1`$ and it is substained over a crossing time, $`10^4`$ years. Shock ionization may thus produce important ionization effects which, however, can be only both local and transient. ## 6 Conclusions We have studied in detail the dynamics of the interaction of a jet with a large cloud pre-existing in the ISM in order to find the conditions for which it is possible to reproduce the main physical parameters of the NLR emitting material. Following the suggestion by Steffen et al. (1997b) that the most relevant effects of the interaction arise when a jet hits dense massive clouds, we adopted a quite simplified geometry of a single gas condensation with can be astrophysically identified with a giant molecular cloud. As the interaction last for a time considerably longer than the cloud crossing time more than one cloud will be interacting at any given time and they will display simultaneously the different evolutionary stages of the interaction. Furthermore the characteristic jet-line structure of the NLR is thus reproduced. In any event, this case, i.e. the head-on collision with a large cloud, is the most efficient case of interaction, for the compression, acceleration and heating of the NLR material. We concentrated our efforts on the exploration of the parameter plane ($`v_{\mathrm{jet}},\rho _{\mathrm{cloud}})`$, since the other parameters, on which the simulation depends, have little influence on the properties of the optically emitting material. We have found that the condition for obtaining values of density, temperature and velocity in the observed range can be translated in a condition on the parameter $`\tau `$, which is the ratio of the cloud crossing timescale to the radiative timescale ($`0.3<\tau <0.55`$) and which depends on our two fundamental parameters $`v_{\mathrm{jet}}`$ and $`\rho _{\mathrm{cloud}}`$. For small values of $`\tau `$, radiation is inefficient and it is not possible to produce regions dense enough, while, on the other hand, for large values of $`\tau `$, the cloud is too dense and the obtained velocities are too low. We have explored a range of cloud densities which can be considered typical of GMCs and, for this range, the jet velocities span an interval from $`4000`$km s<sup>-1</sup> to $`55,000`$km s<sup>-1</sup>. The jet kinetic power corresponding to these combinations of parameters (for a jet density of 1 cm<sup>-3</sup>) ranges from $`3.2\times 10^{41}`$ to $`8\times 10^{44}`$ erg s<sup>-1</sup>, in general agreement with the estimates of Capetti et al. 1999 for Mrk 3. For jet density much lower than 1 cm <sup>-3</sup>, however, in order to match the observed NLR conditions we would need a correspondingly higher velocity and therefore untenable requirements on the kinetic power which grows with $`v_{\mathrm{jet}}^3`$. We conclude that jets in Seyfert galaxies are unlikely to have densities much lower than 1 cm <sup>-3</sup> and velocities higher than $`50,000`$km s<sup>-1</sup>, and therefore they are very different from their counterparts in radio-galaxies in which densities are much lower and velocities are relativistic. Concerning radio-galaxies we can speculate that with lower jet densities and higher velocities, the gas postshock temperatures and radiative time would be increased with respect to the case of Seyfert galaxies and therefore the conditions for having efficient line emission would be more difficult to meet. In addition the different properties of the jet environment in the elliptical galaxies hosting radio-galaxies render encounters with gas condensations less likely to occur. This probably explain why the association between radio and line emission although often present in radio-galaxies (e.g. Baum and Heckman 1989) is not as strong as in Seyfert galaxies. Finally, the study of the global dynamics allowed us to have estimates of the overall efficiency of the conversion of kinetic to high frequency radiative power in the shocks that form in the interaction between jet and ambient medium. We have found that the efficiency is increased by the presence of the cloud, its peak value is 0.1 - 2 % , its typical value is much lower $`10^45\times 10^3`$ and it decreases with the jet power. These results lead us to the conclusion that radiation emitted in shocks can be only a small fraction of the overall ionization budget of the NLR, although it can have local and transient important effects. * We thank CNAA (Consorzio Nazionale per l’Astronomia e Astrofisica) for supporting the use of supercomputers at CINECA, .
warning/0002/cond-mat0002078.html	ar5iv	text	# Axial Anomaly in Quasi-1D Chiral Superfluids ## Abstract The axial anomaly in a quasi-one-dimensional (quasi-1D) chiral $`p`$-wave superfluid model, which has a $`\epsilon _xp_x+i\epsilon _yp_y`$-wave gap in 2D is studied. The anomaly causes an accumulation of the quasiparticle and a quantized chiral current density in an inhomogeneous magnetic field. These effects are related to the winding number of the gap. By varying the parameters $`\epsilon _x`$ and $`\epsilon _y`$, the model could be applicable to Sr<sub>2</sub>RuO<sub>4</sub> near the second superconducting transition point, some quasi-1D organic superconductors and the fractional quantum Hall state at $`\nu =5/2`$ Landau level filling factor. The chiral superfluidity is realized in the superfluid <sup>3</sup>He-A . Recently, the possibility of the chiral superconductivity is argued. In such superfluids or superconductors, the ground state is the condensate of the Cooper pairs which have orbital angular momentum along a same direction. Therefore, time-reversal symmetry (T) and also parity (P) in two-dimensional space (2D) are violating. We investigate a quasi-1D chiral $`p`$-wave superfluid in 2D. It is revealed that the axial anomaly causes P- and T-violating phenomena related to the quantized number. The axial anomaly has been originally pointed out in the Dirac QED in 3D. It is a phenomenon that a symmetry under the phase transformation $`e^{i\gamma _5\alpha }`$ of the Dirac field in the action at the classical level is broken in the quantum theory. Here, $`\alpha `$ is a constant and $`\gamma _5`$ is a hermitian matrix which anti-commutes with all of the Dirac matrices $`\gamma _\mu `$, where $`\mu `$ is the spacetime index. The Adler-Bardeen’s theorem guarantees the absence of higher order collections to the divergence of the axial current. Therefore, the exact calculation of the two-photon decay rate of neutral $`\pi `$ meson can be done. It has been pointed out that the same results are obtained by using the path-integral formalism and has been clarified the relation between the axial anomaly and topological quantized numbers through the Atiyah-Singer index theorem. It has been pointed out that the axial anomaly also plays important role in the quantum Hall effect (QHE) in the 2D massive Dirac QED. In 2D, the mass term of the Dirac Fermion violates P and T like the magnetic field, and the Hall effect may occur. It was shown that the existence of the Hall current and its quantization are caused by the axial anomaly in 1D. The relation between the axial anomaly and QHE in 2D electron gas in the magnetic field was also discussed, and the quantized Hall conductance is expressed by the winding number of the fermionic propagator in the momentum space. Other applications of the axial anomaly to the condensed matter physics are studied in the field of the superfluid <sup>3</sup>He in 3D and in charge density waves in 1D conductors. The phenomena caused by the axial anomaly are related to the topologically quantized numbers. On the analogy of QHE, it is expected that the axial anomaly also plays important role in other P- and T-violating 2D systems. In this letter, we investigate the axial anomaly in a quasi-1D chiral superfluid model in 2D, which has the spin-triplet $`\epsilon _xp_x+i\epsilon _yp_y`$-wave symmetry. P and T-violation occur whenever both of $`\epsilon _x`$ and $`\epsilon _y`$ are non-zero. We show that the axial anomaly in 1D causes an accumulation of the mass density of the quasiparticle in an inhomogeneous magnetic field. The axial anomaly also causes a chiral current density, which is perpendicular to the gradient of the magnetic field. These effects are related to the winding number of the gap; $`\mathrm{sgn}(\epsilon _x\epsilon _y)`$. Our discussion would be valid for the superconductors by taking into account the Meissner effect. By varying the parameters $`\epsilon _x`$ and $`\epsilon _y`$, the model could be applicable to Sr<sub>2</sub>RuO<sub>4</sub> near the second superconducting transiton point, and some quasi-1D organic superconductors or the fractional quantum Hall (FQH) state at $`\nu =5/2`$ Landau level (LL) filling factor . We use the 2+1-dimensional Euclidian spacetime and the natural unit ($`\mathrm{}=c=1`$) in the present paper. Let us consider a quasi-1D chiral superfluid model. We assume a linearized fermion spectrum and a spin-triplet chiral $`p`$-wave gap near the Fermi surface in the normal state written as, $`ϵ_{\mathrm{R},\mathrm{L}}(𝐩)`$ $`=`$ $`\pm v_\mathrm{F}(p_xp_\mathrm{F}),`$ (1) $`\mathrm{\Delta }(𝐩)`$ $`=`$ $`i\sigma _3\sigma _2{\displaystyle \frac{\mathrm{\Delta }}{\|p_F\|}}(\epsilon _xp_x+i\epsilon _yp_y),`$ (2) where $`v_\mathrm{F}`$ and $`p_\mathrm{F}`$ are the Fermi velocity and the Fermi momentum, respectively. $`ϵ_R(𝐩)`$ ($`ϵ_L(𝐩)`$) is the kinetic energy for the right (left) mover. When $`\epsilon _x<<1`$ and $`\epsilon _y1`$, the model describes the low energy excitations (the quasiparticle excitations around the tiny gap points) of Sr<sub>2</sub>RuO<sub>4</sub> near the second superconducting transition point under the uniaxial pressure in the $`xy`$ plane (the basal plane). For simplicity, we assume a circular Fermi Surface in the normal state (Fig.1(a)). When $`\epsilon _x1`$ and $`\epsilon _y<<1`$, the model describes the excitations near $`p_x=\pm p_\mathrm{F}`$ with the chiral $`p`$-wave gap, whose kinetic energy in the normal state is $`ϵ(𝐩)=2t_x\mathrm{cos}(p_xa)2t_y\mathrm{sin}(p_yb)ϵ_\mathrm{F}`$ $`(t_y<<t_x,ϵ_\mathrm{F}`$; the Fermi energy). The model in this case is applicable to some quasi-1D organic superconductors or the FQH state at $`\nu =5/2`$ LL filling factor (Fig.1(b)). The quasi-1D superconductivity has been observed in organic conductors, such as (TMTSF)<sub>2</sub>X. The NMR knight shift study in Ref. is a evidence supporting a spin-triplet pairing state in (TMTSF)<sub>2</sub>PF<sub>6</sub>. The spin-triplet superconductivity in a quasi-1D system with a nodeless gap is obtained theoretically when an electron-phonon coupling and antiferromagnetic fluctuations are taken into account. Our discussion would be valid for such superconductors if they have the chiral $`p`$-wave pairing symmetry. It has been pointed out that the unidirectional charge density wave state, which has the belt-shaped Fermi sea like Fig. 1(b), seems to be the most plausible compressible state at the half-filled Landau levels in the quantum Hall system. Recently, the FQH effect has been observed at $`\nu =5/2`$ , and the effect could be described by the chiral $`p`$-wave pairing state (the Pfaffian state). Therefore, our model could be a candidate of the $`\nu =5/2`$ FQH state. The Lagrangian of our model is written as $``$ $`=`$ $`\overline{\mathrm{\Psi }}_𝐩[\{i_\tau +\mu {\displaystyle \frac{d\mathrm{B}_z}{dy}}y\sigma _3\}\gamma _\tau `$ (4) $`+{\displaystyle \frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}}\sigma _3(\epsilon _xp_x\gamma _x+\epsilon _yp_y\gamma _y)iv_\mathrm{F}(p_xp_\mathrm{F})]\mathrm{\Psi }_𝐩.`$ Here, we use the Bogoliubov-Nambu representation with an isospin $`\alpha =1,2`$ $$\mathrm{\Psi }(𝐱)=e^{ip_\mathrm{F}x}\left(\begin{array}{c}\psi (p_\mathrm{F},𝐱)\\ i\sigma _2\psi ^{}(p_\mathrm{F},𝐱)\end{array}\right),$$ and $`\mathrm{\Psi }_𝐩`$ is its Fourier transform. $`\psi (p_\mathrm{F},𝐱)`$ and $`\psi (p_\mathrm{F},𝐱)`$ are the slowly varying fields for the right mover and the left mover with a real spin index $`s=1,2`$, respectively. The matrices $`\gamma _\tau =\tau _3,\gamma _x=\tau _2`$ and $`\gamma _y=\tau _1`$ are the $`2\times 2`$ Pauli matrices with isospin indices and $`\sigma _i(i=1,2,3)`$ is the $`2\times 2`$ Pauli matrices with spin indices. The symbol $`\sigma _i\gamma _{\tau ,x,y}`$ shows the direct product. $`\overline{\mathrm{\Psi }}`$ is defined as $`\overline{\mathrm{\Psi }}=i\mathrm{\Psi }^{}\gamma _\tau `$. We assume a magnetic field, which is directed to the $`z`$-axis (the $`c`$-axis in the crystal) and has a constant gradient in the $`y`$-direction, i.e. $`B_z(y)=(dB_z/dy)y`$, and $`(dB_z/dy)=const.`$ The magnetic field couples with the Fermion through the Zeeman term $`\mu B_z\overline{\mathrm{\Psi }}\sigma _3\gamma _\tau \mathrm{\Psi }`$, where $`\mu `$ is the magnetic moment of the Fermion. We note that the Lagrangian is similar to that of 2D Dirac QED in a background scalar potential, except for the last term and the existence of $`\sigma _3`$. The axial anomaly in such a system is discussed in Ref. . Let us calculate the expectation value of the mass density $`\rho _e(𝐱)`$ $`=`$ $`e\overline{\mathrm{\Psi }}(𝐱)\mathrm{\Psi }(𝐱)`$ (5) $`=`$ $`\mathrm{Tr}\left[{\displaystyle \frac{e\sigma _3}{\overline{)}d+\frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x\widehat{1}\gamma _xiv_\mathrm{F}(p_xp_\mathrm{F})\sigma _3\widehat{1}}}\right],`$ (6) where $`e`$ shows a mass of the quasiparticle. It shows an electric charge when we consider a superconductor. A hermitian operator $`\overline{)}d`$ in the $`y`$-direction is defined as $$\overline{)}d=\left\{i_\tau \sigma _3+\mu \frac{dB_z}{dy}y\right\}\gamma _\tau +\frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _yp_y\widehat{1}\gamma _y.$$ (7) We define $`\gamma _5`$ as $`\gamma _5=i\gamma _\tau \gamma _y=\gamma _x`$, and it is anti-commute with $`\gamma _\tau `$ and $`\gamma _y`$, therefore, $`\gamma _5`$ is a hermitian matrix and satisfies $`\{\gamma _5,\overline{)}d\}=0`$. These facts suggest that if an eigenstate $`u_n`$ of $`\overline{)}d`$ with a nonzero eigenvalue $`\xi _n`$ ($`0<n`$) exists (i.e. $`\overline{)}du_n=\xi _nu_n`$), $`\gamma _5u_n`$ should be another eigenstate with an eigenvalue $`\xi _n`$. If zeromodes of $`\overline{)}d`$ exist (i.e. $`\overline{)}du_0=0`$ and $`\overline{)}d\gamma _5u_0=0`$), they are divided into two groups. One of them is $`u_0^{(+)}=(1/2)(1+\gamma _5)u_0`$ with an eigenvalue $`\gamma _5=+1`$ and another is $`u_0^{()}=(1/2)(1\gamma _5)u_0`$ with an eigenvalue $`\gamma _5=1`$, since $`\gamma _5^2=1`$. Let us research eigenmodes of $`\overline{)}d`$. The expectation value of $`\overline{)}d^2`$ is $`(u_n,\overline{)}d^2u_n)`$ $`=`$ $`\|\omega _c\|(n+{\displaystyle \frac{1}{2}})+{\displaystyle \frac{\omega _c}{2}}(u_n,\gamma _5u_n),`$ (8) $`\omega _c`$ $`=`$ $`\mu {\displaystyle \frac{dB_z}{dy}}{\displaystyle \frac{2\mathrm{\Delta }}{\|p_\mathrm{F}\|}}\epsilon _y,`$ (9) where $`u_n=u_n(yy_c(p_\tau ,\sigma _3))`$ is the eigenfunction of the harmonic oscillator with the frequency $`\omega _c`$. The oscillator is centered at $`y_c(p_\tau ,\sigma _3)=(dB_z/dy)^1(p_\tau /\mu )\sigma _3.`$ Eq. (9) indicates that only zeromodes which belong to $`u_0^{}`$ ($`u_0^+`$) exist when $`0<\omega _c`$ ($`\omega _c<0`$). It suggests the nonconservation of the vacuum expectation value of the axial charge which is defined in the second-quantized formalism as $`Q_5`$ $`=`$ $`N_+N_{},`$ (10) $`N_\pm `$ $`=`$ $`{\displaystyle 𝑑p_y\widehat{u}_0^{(\pm )}\widehat{u}_0^{(\pm )}},`$ (11) while the classical 1D theory $`_{1\mathrm{D}}=\overline{\mathrm{\Psi }}\overline{)}d\mathrm{\Psi }`$ has the axial symmetry $`\mathrm{\Psi }e^{i\alpha \gamma _5}\mathrm{\Psi }`$, i.e. the axial anomaly occurs. Here, $`\widehat{u}_0^\pm `$ is a second-quantized fermionic field. The anomaly comes from the spectral asymmetry of zeromodes as same as the discussions in Refs.. In the free system, the energy spectrum of $`u_0^{(\pm )}`$ is $`p_0=\pm \frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _yp_y`$ in Minkowski spacetime, and all of the negative energy states are filled while all of the positive energy states are empty and $`Q_5=0`$. After we turn on the magnetic field adiabatically (for a while, we assume $`0<\omega _c`$), the energy spectrum of $`u_0^{(+)}`$ is lowered and $`N_+`$ decreases (i.e. empty negative energy states arise on the spectrum of $`u_0^{(+)}`$), on the other hand, the energy spectrum of $`u_0^{()}`$ is lifted and $`N_{}`$ increases (i.e. filled positive energy states arise on the spectrum of $`u_0^{()}`$), therefore $`Q_5`$ does not conserve. Finally $`N_+=0`$ and only $`u_0^{()}`$ exists. The nonzero eigenvalues of $`\overline{)}d^2`$ is $`E_n=\omega _c(n+1/2)`$, since the inner product $`(u_n,\gamma _5u_n)`$ vanishes whenever $`\overline{)}du_n0`$ because of the orthogonal relation between the eigenfunctions of the hermitian operator. Next, we consider the eigenvalue problem of a 2D operator $$\overline{)}D=\overline{)}d+\frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x\widehat{1}\gamma _x=\overline{)}d\frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x\widehat{1}\gamma _5.$$ (12) Let $$\phi _n=(\alpha _nu_n+\beta _n\gamma _5u_n)e^{ip_xx}$$ (13) stands for an eigenfunction. We use a representation for the $`n`$-th level such as ($``$) $$\overline{)}d=\left(\begin{array}{cc}\xi _n& 0\\ 0& \xi _n\end{array}\right),u_n=\left(\begin{array}{c}1\\ 0\end{array}\right),\gamma _5=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),$$ (14) where, $$\xi _n=\{\begin{array}{cc}\sqrt{\|\omega _c\|(n+\frac{1}{2})}& (n=1,2,\mathrm{}),\hfill \\ 0& (n=0).\hfill \end{array}$$ Therefore, the eigenvalue equation is written as $$\left(\begin{array}{cc}\xi _n& \frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x\\ \frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x& \xi _n\end{array}\right)\left(\begin{array}{c}\alpha _n\\ \beta _n\end{array}\right)=\zeta _n\left(\begin{array}{c}\alpha _n\\ \beta _n\end{array}\right).$$ (15) There are two eigenstates for an oscillator in the $`n(0)`$-th level written as $`\zeta _n^{(\pm )}(p_x)`$ $`=\pm \sqrt{\xi _n^2+{\displaystyle \frac{\mathrm{\Delta }^2}{p_\mathrm{F}^2}}\epsilon _x^2p_x^2},`$ (16) $`\left(\begin{array}{c}\alpha _n^+\\ \beta _n^+\end{array}\right)`$ $`={\displaystyle \frac{1}{C_+}}\left(\begin{array}{c}\zeta _n^{(+)}+\xi _n\\ \frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x\end{array}\right),`$ (21) $`\left(\begin{array}{c}\alpha _n^{}\\ \beta _n^{}\end{array}\right)`$ $`={\displaystyle \frac{1}{C_{}}}\left(\begin{array}{c}\frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}\epsilon _xp_x\\ \zeta _n^{()}+\xi _n\end{array}\right),`$ (26) where $`C_\pm `$ are normalization constants, but for $`n=0`$, there is only one eigenstate $`\zeta _0(p_x)`$ $`=`$ $`{\displaystyle \frac{\omega _c}{\|\omega _c\|}}{\displaystyle \frac{\mathrm{\Delta }}{\|p_\mathrm{F}\|}}\epsilon _xp_x,`$ (27) $`\left(\begin{array}{c}\alpha _0\\ \beta _0\end{array}\right)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\begin{array}{c}1\\ \omega _c/\|\omega _c\|\end{array}\right),`$ (32) because the solution should satisfy $`\gamma _5\phi _0=(\omega _c/\|\omega _c\|)\phi _0`$. This condition comes from the axial anomaly in the $`y`$-direction. Finally, we show the accumulation of the mass density from Eq. (5), which is derived as $`\rho _e(𝐱)`$ $`=`$ $`\mathrm{Tr}\left[{\displaystyle \frac{e\sigma _3}{\overline{)}Div_\mathrm{F}(p_xp_\mathrm{F})\sigma _3\widehat{1}}}\right]={\displaystyle \underset{n}{}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{dp_\tau }{2\pi }}{\displaystyle \frac{dp_x}{2\pi }\mathrm{tr}\left[\frac{e\sigma _3\|u_n(yy_0(p_\tau ,\sigma _3))\|^2}{\zeta _n(p_x)iv_\mathrm{F}(p_xp_\mathrm{F})\sigma _3}\right]}`$ (33) $`=`$ $`{\displaystyle \frac{e\mu }{2\pi }}{\displaystyle \frac{dB_z}{dy}}{\displaystyle \underset{n0}{}}{\displaystyle \frac{dp_x}{2\pi }\mathrm{tr}\left[\frac{1}{\zeta _n^{(+)}(p_x)iv_\mathrm{F}(p_xp_\mathrm{F})\sigma _3}+\frac{1}{\zeta _n^{()}(p_x)iv_\mathrm{F}(p_xp_\mathrm{F})\sigma _3}\right]}`$ (35) $`{\displaystyle \frac{e\mu }{2\pi }}{\displaystyle \frac{dB_z}{dy}}{\displaystyle \frac{dp_x}{2\pi }\mathrm{tr}\left[\frac{1}{\zeta _0(p_x)iv_\mathrm{F}(p_xp_\mathrm{F})\sigma _3}\right]}=\mathrm{sgn}(\epsilon _x\epsilon _y)e\mu N_{1\mathrm{D}}(0){\displaystyle \frac{dB_z}{dy}},`$ where the symbol $`tr`$ means a trace on the real spin, and we use the normal-orthogonal relation $`𝑑y\|u_n\|^2=1`$. $`\frac{dp_x}{2\pi }=_{p_\mathrm{F}\mathrm{\Lambda }}^{p_\mathrm{F}+\mathrm{\Lambda }}\frac{dp_x}{2\pi }`$ , and $`\mathrm{\Lambda }`$ is a momentum cutoff. We assume a relation $`\|\mathrm{\Delta }\|<<\mathrm{\Lambda }^2/2m<<ϵ_F`$. $`N_{1\mathrm{D}}(0)=(2\pi v_\mathrm{F})^1`$ is the density of state at the Fermi surface in 1D. All of the $`n0`$ parts are canceled out because of the co-existence of the eigenvalues $`\zeta _n^{(+)}`$ and $`\zeta _n^{()}`$. Only the $`n=0`$ part survives because of the axial anomaly in the $`y`$-direction. We can define a chiral transformation in the $`x`$-direction such as $`\psi (\pm p_\mathrm{F},𝐱)e^{\pm i\alpha }\psi (\pm p_\mathrm{F},𝐱)`$, therefore $`\mathrm{\Psi }e^{i\alpha }\mathrm{\Psi },\overline{\mathrm{\Psi }}\overline{\mathrm{\Psi }}e^{i\alpha }`$. The expectation value of the corresponding current density which is perpendicular to $`dB_z/dy`$ is $$j_x^{chi}(𝐱)=ev_\mathrm{F}\overline{\mathrm{\Psi }}(𝐱)\mathrm{\Psi }(𝐱)=\mathrm{sgn}(\epsilon _x\epsilon _y)\frac{e\mu }{2\pi }\frac{dB_z}{dy},$$ (36) and we call it a chiral Hall current density. These two effects are related to the winding number of the gap Eq. (2), $`\mathrm{sgn}(\epsilon _x\epsilon _y)`$ $`=`$ $`{\displaystyle \frac{d^2p}{16\pi }tr[\widehat{𝐠}(\widehat{𝐠}\times \widehat{𝐠})]},`$ (37) $`𝐠(𝐩)`$ $`=`$ $`\left(\begin{array}{c}\mathrm{Re}[\mathrm{\Delta }(𝐩)(i\sigma _2)]\\ \mathrm{Im}[\mathrm{\Delta }(𝐩)(i\sigma _2)]\\ (𝐩^2/2m)ϵ_\mathrm{F}\end{array}\right),`$ (41) where $`=/𝐩`$. It suggests that these effects occur even if $`\epsilon _x`$ and/or $`\epsilon _y`$ are infinitesimally small, and that these effects come from the P- and T-violation of the gap. The accumulated mass density and the chiral Hall current density exist in the bulk region of the superfluid. In the superconductors, the Meissner effect occurs and the magnetic field cannot penetrate into the bulk, therefore the accumulated charge density and the chiral Hall current density would exist near the edge of the superconductors and also around the vortex core. As we mentioned before, our discussion could be applicable to Sr<sub>2</sub>RuO<sub>4</sub> near the second superconducting phase transition point, some quasi-1D organic superconductors and the FQH state at $`\nu =5/2`$ LL filling factor by varying the parameters $`\epsilon _x`$ and $`\epsilon _y`$. Recently, the vortex in chiral superconductors has been discussed , and such a vortex has a fractional charge and a fractional angular momentum. Interesting phenomena related to these fractional quantum numbers and the present effects are expected to occur around the vortex core. The axial anomaly also causes the spin quantum Hall effect (SQHE) in the chiral $`d`$-wave ($`d_{x^2y^2}+id_{xy}`$-wave) superconductors. The low energy quasiparticles in a magnetic field with a constant gradient can be mapped onto the massive Dirac Fermion in a constant electric field, and the spin rotation around the $`z`$-axis for the quasiparticle corresponds to the $`U(1)`$ transformation for the Dirac Fermion. Therefore, according to the discussions in Ref., we can see that the axial anomaly causes the quantized spin Hall current, which is perpendicular to the gradient of the magnetic field. SQHE has been pointed out by Volovik and Yakovenko in superfluid <sup>3</sup>He-A film, which is the chiral $`p`$-wave superfluid. They have described the effect by the Chern-Simons term. It has been clarified the relation between the axial anomaly and the Chern-Simons term in 2D Dirac QED. Therefore, SQHE in <sup>3</sup>He-A could be related to the axial anomaly. According to Ref., SQHE also occurs at the edge or around the vortex core of the superconducting Sr<sub>2</sub>RuO<sub>4</sub> by a magnetic field in the basal plane. The author thanks K. Ishikawa and N. Maeda for useful discussions and encouragement. This work was partially supported by the special Grant-in-Aid for Promotion of Education and Science in Hokkaido University provided by the Ministry of Education, Science, Sports, and Culture, the Grant-in-Aid for Scientific Research on Priority area (Physics of CP violation) (Grant No. 10140201), and the Grant-in-Aid for International Science Research (Joint Research 10044043) from the Ministry of Education, Science, Sports and Culture, Japan.
warning/0002/math0002222.html	ar5iv	text	# Quantization of bending deformations of polygons in 𝔼³, hypergeometric integrals and the Gassner representation ## 1 Introduction In \[KM\] and \[Kly\] certain Hamiltonian flows on the moduli space $`M_r`$ of $`n`$-gon linkages in $`𝔼^3`$ were studied. In \[KM\] these flows were interpreted geometrically and called bending deformations of polygons. In \[Kly\], Klyachko pointed out that the Hamiltonian potentials of the bending deformations gave rise to a Hamiltonian action of $`𝒫_n`$, the Malcev Lie algebra of the pure braid group $`P_n`$ (see §3), on $`M_r`$. It is a remarkable fact, see \[K1, Lemma 1.1.4\], that a representation $`\rho :𝒫_nEnd(V)`$, $`dim(V)<\mathrm{}`$, gives rise to a flat connection $``$ on the vector bundle $`_{}^n\times V`$ over $`_{}^n`$, the space of distinct points in $``$. Accordingly the monodromy representation of $``$ yields a representation $`\widehat{\rho }:P_nAut(V)`$. We see then that if we can find a finite dimensional representation of the Lie algebra $`C^{\mathrm{}}(M_r)`$ generated by the bending Hamiltonians under the Poisson bracket , i.e. if we can “quantize” $``$, then we will obtain a representation of $`P_n`$. Klyachko suggested using a geometric quantization of $`M_r`$ to quantize $``$. This appears to be difficult to carry out because the bending flows do not preserve a polarization. Note however that the problem of quantizing a Poisson subalgebra of $`C^{\mathrm{}}(M_r)`$ can be solved immediately if the functions in the subalgebra have a common critical point $`xM_r`$. For in this case we may simultaneously linearize all the Hamiltonian fields at $`x`$. We are fortunate that simultaneous critical points for the algebra $``$ exist if $`M_r`$ is singular. Indeed, a degenerate $`n`$-gon (i.e. an $`n`$-gon which is contained in a line $`L`$) is a critical point of all bending Hamiltonians. The point of this paper is to compute the representation $`\widehat{\rho }_{ϵ,r}:P_nAut(T_{ϵ,r})`$ associated to a degenerate $`n`$-gon $`P`$. Here $`T_{ϵ,r}=T_P(M_r)`$ and $`ϵ=(ϵ_1,\mathrm{},ϵ_n)`$, $`ϵ_i\{\pm 1\}`$, and $`r=(r_1,\mathrm{},r_n),r_i_+`$, are defined as follows. Fix an orientation on $`L`$. The number $`r_i`$ is the length of the $`i`$-th edge of $`P`$. Define $`ϵ_i`$ to be $`+1`$ if the $`i`$-th edge is positively oriented and $`ϵ_i=1`$ otherwise. We call $`ϵ=(ϵ_1,\mathrm{},ϵ_n)`$ the vector of edge-orientation of $`P`$. Our formula for $`\rho _{ϵ,r}:𝒫_nEnd(T_{ϵ,r})`$ is in terms of certain $`n\times n`$ matrices $`J_{ij}(\lambda )`$ which are called Jordan-Pochhammer matrices. Let $`\lambda =(\lambda _1,\mathrm{},\lambda _n)`$ be an $`n`$-tuple of complex numbers. Define matrices $`J_{ij}(\lambda )`$ for $`1i<jn`$ by $$\begin{array}{cccccccc}& & & i\text{th column}& & j\text{th column}& & \\ & (\mathrm{}& 0\mathrm{}0& 0& 0\mathrm{}0& 0& 0\mathrm{}0& )\mathrm{}\\ i\text{th row}& 0\mathrm{}0& \lambda _j& 0\mathrm{}0& \lambda _j& 0\mathrm{}0\\ & 0\mathrm{}0& 0& 0\mathrm{}0& 0& 0\mathrm{}0\\ & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ j\text{th row}& 0\mathrm{}0& \lambda _i& 0\mathrm{}0& \lambda _i& 0\mathrm{}0\\ & 0\mathrm{}0& 0& 0\mathrm{}0& 0& 0\mathrm{}0\\ & \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}=J_{ij}(\lambda ).$$ Define $`J_{ii}=0`$ and $`J_{ij}(\lambda )=J_{ji}(\lambda )`$ for $`i>j`$. We have (as can be verified easily) ###### Lemma 1.1 The matrices $`\{J_{ij}(\lambda )\}`$ satisfy the infinitesimal braid relations: * $`[J_{ij}(\lambda ),J_{kl}(\lambda )]=0`$ if $`\{i,j\}\{k,l\}=\mathrm{}`$. * $`[J_{ij}(\lambda ),J_{ij}(\lambda )+J_{jk}(\lambda )+J_{ki}(\lambda )]=0`$, $`i,j,k`$ are distinct. Consequently the assignment $`\rho _\lambda (X_{ij})=J_{ij}(\lambda )`$ (see Section 3 for the meaning of $`X_{ij}`$) yields a representation $`\rho _\lambda :𝒫_nM_n()`$ and a flat connection $``$ on $`_{}^n\times ^n`$. Here we realize $`^n`$ as the space of row vectors with $`n`$ components. It is immediate that the subspace $`_0^n^n`$ defined by $$_0^n=\{z^n:\underset{i}{}z_i=0\}$$ is invariant under $`\rho _\lambda `$, in fact $`\rho _\lambda (𝒫_n)(^n)_0^n`$. Now we assume $`_{i=1}^n\lambda _i=0`$. Then $`\lambda _0^n`$ and we see that $`\rho _\lambda (𝒫_n)(\lambda )=0`$. Thus we have a $`𝒫_n`$-invariant filtration $$\lambda _0^n^n.$$ Define $`W_\lambda =_0^n/\lambda `$. Now let $`P`$ be a degenerate $`n`$-gon with side-lengths $`r=(r_1,\mathrm{},r_n)`$ and edge-orientations $`ϵ=(ϵ_1,\mathrm{},ϵ_n)`$. Our first main theorem is Theorem A. There is a $`𝒫_n`$-invariant almost complex structure $`J^ϵ`$ on $`T_{ϵ,r}`$ such that there is an isomorphism of $`𝒫_n`$-modules $`T_{ϵ,r}^{1,0}W_\lambda `$ for $`\lambda :=(\sqrt{1}ϵ_1r_1,\mathrm{},\sqrt{1}ϵ_nr_n)`$. Here $`T_{ϵ,r}^{1,0}=\{wT_{ϵ,r}:J^ϵw=\sqrt{1}w\}`$. We have Corollary. The flat connection on $`_{}^n\times T_{ϵ,r}^{1,0}`$ has the connection form $$\omega =\underset{1i<jn}{}\frac{dz_idz_j}{z_iz_j}J_{ij}(\lambda )$$ with $`\lambda `$ as above. We then adapt the methods of \[K1\] to give formulae for multivalued parallel sections of $``$ in terms of hypergeometric integrals and to compute the monodromy of $``$. Before stating our first formula for the monodromy of $``$ we need more notation. Let $`\gamma _j`$, $`1jn`$, be the free generators of the free group $`𝔽_n`$. Define the character $`\chi :𝔽_n^{}`$ by $`\chi (\gamma _j)=e^{2\pi i\lambda _j}`$, $`1jn`$ (recall that $`\lambda _j=\sqrt{1}ϵ_jr_j`$). Let $`_{\chi ^1}`$ be the 1-dimensional module (over $``$) in which the free group $`𝔽_n`$ acts by $`\chi ^1`$. The pure braid group $`P_n`$ acts by automorphisms on $`𝔽_n`$ so that the character $`\chi `$ is fixed. Thus we have the associated action of $`P_n`$ on $`\text{H}_1(𝔽_n,_{\chi ^1})`$. We let $`\mathrm{\Gamma }_n=\pi _1(^1\{z_1,\mathrm{},z_n\})`$ be the fundamental group of the $`n`$ times punctured sphere. Hence $`\mathrm{\Gamma }_n`$ is the quotient of $`𝔽_n`$ by the normal subgroup generated by $`\gamma _1\mathrm{}\gamma _n`$. Since $`\chi (\gamma _1\mathrm{}\gamma _n)=1`$, the character $`\chi `$ induces a character of $`\mathrm{\Gamma }_n`$. The group $`P_n`$ fixes $`\gamma _1\mathrm{}\gamma _n`$ and consequently acts on $`\mathrm{\Gamma }_n`$ and on $`\text{H}_1(\mathrm{\Gamma }_n,_{\chi ^1})`$. We can now state Theorem B. The monodromy representation of $``$ is equivalent to the representation of $`P_n`$ on $`\text{H}_1(\mathrm{\Gamma }_n,_{\chi ^1})`$. In §10 we define the Gassner representation of the pure braid group, the reduced Gassner representation and their specializations via characters of the free group. Let $``$ is the $``$-algebra of Laurent polynomials on $`t_1,..,t_n`$. Theorem C. The monodromy representation of $``$ is dual to the quotient of the reduced Gassner representation $`Z^1(\mathrm{\Gamma }_n,)`$ specialized at $`t_j=e^{2\pi ϵ_jr_j}`$, where we quotient by the 1-dimensional subspace $`B^1(\mathrm{\Gamma }_n,_\chi )`$ fixed by $`P_n`$. Our results appear to be related to those of \[DM\] and \[Lo\] but there are significant differences. In \[Lo\], D. D. Long linearizes the action of $`P_n`$ on the moduli space of $`n`$-gon linkages in $`S^3`$ obtained from the action of $`P_n`$ on $$\text{Hom}(\pi _1(S^2\{z_1,\mathrm{},z_n\}),SU(2))/SU(2)$$ by precomposition. The corresponding action of $`P_n`$ on $`M_r`$ is trivial in our case, see \[KM, Remark 5.1\]. In \[DM\], Deligne and Mostow arrive at the Gassner representation by considering a variation of Hodge structure over $`_{}^n/PGL_2()M_r`$. They obtain the quotient (by the 1-coboundaries) of the reduced Gassner representation specialized at $`(e^{2\pi ir_1},\mathrm{},e^{2\pi ir_n})`$; we obtain the dual of the quotient of the reduced Gassner representation specialized at $`(e^{2\pi ϵ_1r_1},\mathrm{},e^{2\pi ϵ_nr_n})`$. Here we must assume $`_{i=1}^nr_i=2`$ to be consistent with \[DM\]. Our representation lies in $`GL(n2,)`$; their representation is in $`U(n3,1)`$. Acknowledgements. It is our pleasure to thank Ragnar Buchweitz and Richard Hain for helpful conversations. ## 2 The moduli space of $`n`$-gon linkages in $`𝔼^3`$. Let $`Pol_n(𝔼^3)`$ be the space of (closed) $`n`$-gons with distinguished vertices in the Euclidean space $`𝔼^3`$. An $`n`$-gon $`P`$ is defined to be an ordered $`n`$-tuple of points $`(v_1,\mathrm{},v_n)(𝔼^3)^n`$. The point $`v_i`$ is called the $`i`$-th vertex of $`P`$. The vertices are joined in cyclic order by edges $`e_1,\mathrm{},e_n`$ where $`e_i`$ is the oriented segment from $`v_i`$ to $`v_{i+1}`$. We think of $`e_i`$ as a vector in $`^3`$. Two polygons $`P=(v_1,\mathrm{},v_n)`$ and $`Q=(w_1,\mathrm{},w_n)`$ are identified if and only if there exists an orientation-preserving isometry $`g`$ of $`𝔼^3`$ such that $`g(v_i)=w_i`$, $`1in`$. Let $`r=(r_1,\mathrm{},r_n)`$ be an $`n`$-tuple of positive real numbers. Then $`M_r`$ is defined to be the moduli space of $`n`$-gons with the side-lengths $`r_1,\mathrm{},r_n`$ modulo isometries as above. An element of $`M_r`$ will be called a closed $`n`$-gon linkage. We will also need the moduli space space $`N_r`$ of “open” $`n`$-gon linkages. To obtain $`N_r`$ we repeat the above construction of $`M_r`$ except we do not assume the end vertex $`v_{n+1}`$ of the edge $`e_n`$ is equal to $`v_1`$. The starting point of \[KM\] was the observation that $$M_r=\{e=(e_1,\mathrm{},e_n)\underset{i=1}{\overset{n}{}}S^2(r_i):e_1+\mathrm{}+e_n=0\}/SO(3).$$ This equality exhibits $`M_r`$ as the symplectic quotient of $`_{i=1}^nS^2(r_i)`$ and has many consequences. First $`M_r`$ is a complex analytic space with isolated (quadratic) singularities. The smooth part of $`M_r`$ is a Kähler manifold. The singular points of $`M_r`$ are the equivalence classes of degenerate $`n`$-gons. Thus $`M_r`$ is singular if and only if $`r`$ is the set of side-lengths of a degenerate $`n`$-gon. In \[KM\] we introduced bending deformations of closed polygonal linkages in $`𝔼^3`$, see also \[Kly\]. Suppose $`P=e=(e_1,\mathrm{},e_n)`$. Let $`I\{1,\mathrm{},n\}`$ be a subset and define $`f_IC^{\mathrm{}}(M_r)`$ by $$f_I(e)=\underset{iI}{}e_i^2.$$ Then $`f_I`$ is the Hamiltonian potential of a Hamiltonian vector field $`B_I`$. The vector $`e_I=_{iI}e_i`$ is constant along an integral curve of $`B_I`$. By \[KM, Lemma 3.5\], $`B_I(e)=(\delta _1,\mathrm{},\delta _n)`$, where $`\delta _i=e_I\times e_i`$, $`iI,`$ and $`\delta _i=0`$ for $`iI`$. The integral curves of $`B_I`$ are obtained as follows. Define an element $`ad(e_I)so(3)`$ by $$ad(e_I)(v)=e_I\times v$$ and a one-parameter group $`R_I(t)SO(3)`$ by $$R_I(t)=\mathrm{exp}(tad(e_I)).$$ Then the integral curve $`e(t)`$ of $`B_I`$ passing through $`e`$ is given by $$e_i(t)=R_I(t)e_i,iI$$ $$e_j(t)=e_j,jI.$$ This motion of a polygon $`P`$ has a simple geometric interpretation if the elements of $`I`$ are consecutive. In this case $`e_I`$ is a diagonal and it divides the polygon into two parts. Keep one part fixed and bend the polygon by rotating the other part around the diagonal with the angular speed $`e_I`$. For this reason we call the above motion a bending deformation of the polygon. We will be specifically interested in the case $`I=\{i,j\}`$, $`i<j`$. We abbreviate $`f_{\{i,j\}}`$ to $`f_{ij}`$ and $`B_{\{i,j\}}`$ to $`B_{ij}`$ . We have: $$f_{ij}(e)=e_i+e_j^2.$$ ###### Lemma 2.1 Let $`eM_r`$ be a degenerate polygon. Then $`B_{ij}(e)=0`$ for all $`i,j`$. Proof: The bending field $`B_{ij}`$ is given by $$B_{ij}(e)=(0,\mathrm{},(e_i+e_j)\times e_i,0,\mathrm{},(e_i+e_j)\times e_j,0,\mathrm{})=(0,\mathrm{},e_j\times e_i,0,\mathrm{},e_i\times e_j,0,\mathrm{}).$$ If $`e`$ is degenerate then $`e_i`$ and $`e_j`$ are linearly dependent, so $`e_i\times e_j=0`$. $`\mathrm{}`$ ###### Remark 2.2 In fact $`B_I(e)=0`$ for all $`I`$ if $`e`$ is degenerate. Define $`\stackrel{~}{N}_r:=_{i=1}^nS^2(r_i)`$ where $`S^2(r_i)`$ is the round 2-sphere of the radius $`r_i`$. We also define $`\stackrel{~}{M}_r\stackrel{~}{N}_r`$ by $$\stackrel{~}{M}_r=\{e\stackrel{~}{N}_r:\underset{i=1}{\overset{n}{}}e_i=0\}.$$ Hence $`N_r`$ is the quotient of $`\stackrel{~}{N}_r`$ by $`SO(3)`$ and $`M_r`$ is the quotient of $`\stackrel{~}{M}_r`$ by $`SO(3)`$. ## 3 The Malcev Lie algebra of the pure braid group. Let $`P_n`$ be the pure braid group on $`n`$ strands in $``$ (see \[C, §1\]). Let $`_{}^n`$ denote the subset of $`^n`$ consisting of distinct $`n`$-tuples. Then $`P_n`$ is isomorphic to the fundamental group of $`_{}^n`$. Let $`𝒫_n`$ be the Malcev Lie algebra of $`P_n`$, see \[ABC\]. Kohno found the following presentation for $`𝒫_n`$ in \[K2\] (see also \[I, Proposition 3.2.1\]). ###### Lemma 3.1 The Lie algebra $`𝒫_n`$ is the quotient of the free Lie algebra over $``$ generated by $`X_{ij},1i,jn`$, subject to the relations: 1. $`X_{ii}=0`$, $`1in`$. 2. $`X_{ij}=X_{ji}`$, $`1i,jn`$ 3. $`[X_{ij},X_{kl}]=0`$ if $`\{i,j\}\{k,l\}=\mathrm{}`$. 4. $`[X_{ij},X_{ij}+X_{jk}+X_{ki}]=0`$, $`i,j,k`$ are distinct. We will now see that any finite dimensional representation of $`𝒫_n`$ induces a finite dimensional representation of $`P_n`$ on the same vector space. This remarkable fact is an immediate consequence of the following lemma of Kohno \[K1, Lemma 1.1.4\]. ###### Lemma 3.2 Suppose $`V`$ is a finite dimensional vector space and $`A_{ij},1i,jn`$, are elements of $`End(V)`$ such that $`A_{ii}=0`$ and $`A_{ij}=A_{ji}`$. Let $``$ be the connection on the trivial $`V`$ bundle over $`_{}^n`$ with connection form $$\omega =\underset{1i<jn}{}\frac{dz_idz_j}{z_iz_j}A_{ij}.$$ Then $``$ is flat if and only if the relations (3) and (4) for $`𝒫_n`$ are satisfied by the $`A_{ij}`$’s. Thus there is a 1-1 correspondence between Lie algebra homomorphisms $`\rho :𝒫_nEnd(V)`$ and flat connections $``$ on $`_{}^n\times V`$ of the above form. Suppose we are given $`\rho `$ as above. Since $`\pi _1(_{}^n,z)P_n`$ ($`z`$ is a base-point), the monodromy representation of $``$ gives an induced representation of $`P_n`$ to $`Aut(V)`$. Let $`F:_{}^nV`$ be a smooth map. Then $`F`$ induces a parallel section of $``$ if and only if $`F`$ satisfies the equation (of the $`V`$-valued 1-forms on $`_{}^n`$) $$dF=\underset{1i<jn}{}\frac{dz_idz_j}{z_iz_j}A_{ij}(F).$$ ## 4 A Hamiltonian action of $`𝒫_n`$ on $`M_r`$. We define the function $`f_{ij}`$ on $`\stackrel{~}{N}_r`$ by $$f_{ij}(e)=e_i+e_j^2.$$ The next proposition was proved in \[Kly\]. Since it is central to our paper we give a proof here. ###### Proposition 4.1 1. $`f_{ij}=f_{ji}`$. 2. $`\{f_{ij},f_{kl}\}=0`$, if $`\{i,j\}\{k,l\}=\mathrm{}`$. 3. $`\{f_{ij},f_{ij}+f_{jk}+f_{ki}\}=0`$, if $`i,j,k`$ are distinct. Proof: The assertions (1) and (2) are obvious. The third assertion will be a consequence of the following discussion. Since $`\stackrel{~}{N}_r`$ is a symplectic leaf of the Lie algebra $`(^3,\times )`$ equipped with the Lie Poisson structure it suffices to prove (3) for the functions $`f_{ij}`$ extended to $`(^3)^n`$ using the same formula. Let $`g_{ij}:(^3)^n`$ be given by $`g_{ij}(e)=e_ie_j`$ and $`h_{ijk}:(^3)^n`$ be given by $`h_{ijk}(e)=e_i(e_j\times e_k)`$. ###### Lemma 4.2 $`\{g_{ij},g_{jk}\}=h_{ijk}`$. Proof: It suffices to prove the lemma for $`i=1,j=2,k=3`$. We use coordinates $`(x_i,y_i,z_i),1in`$, on $`(^3)^n`$. Then $$\{x_i,y_i\}=z_i,\{y_i,z_i\}=x_i,\{z_i,x_i\}=y_i,1in.$$ We have $$\{g_{12},g_{23}\}=\{x_1x_2+y_1y_2+z_1z_2,x_2x_3+y_2y_3+z_2z_3\}=$$ $$=\{x_1x_2,y_2y_3\}+\{x_1x_2,z_2z_3\}+\{y_1y_2,x_2x_3\}+\{y_1y_2,z_2z_3\}+\{z_1z_2,x_2x_3\}+\{z_1z_2,y_2y_3\}=$$ $$=x_1y_3z_2x_1y_2z_3x_3y_1z_2+x_2y_1z_3+x_3y_2z_1x_2y_3z_1=e_1(e_2\times e_3).\text{ }\text{ }\mathrm{}$$ ###### Corollary 4.3 $`\{f_{ij},f_{jk}\}=4e_i(e_j\times e_k)`$. Proof: $`f_{ij}=f_{ii}+f_{jj}+2g_{ij}`$. But $`f_{ii}`$ and $`f_{jj}`$ are Casimirs. $`\mathrm{}`$ We now prove the 3-rd assertion. $$\{f_{ij},f_{ij}+f_{jk}+f_{ki}\}=\{f_{ij},f_{jk}\}+\{f_{ij},f_{ki}\}=\{f_{ij},f_{jk}\}+\{f_{ji},f_{ik}\}=$$ $$\{f_{ij},f_{jk}\}+\{f_{ji},f_{ik}\}=4e_ie_j\times e_k4e_je_i\times e_k=4e_ie_j\times e_k+4e_i(e_j\times e_k)=0.$$ $`\mathrm{}`$ Since the function $`f_{ij}`$ is $`SO(3)`$-invariant it induces a function (which is again denoted by $`f_{ij}`$) on $`M_r`$. The Poisson bracket of these functions remain the same and we obtain ###### Theorem 4.4 There exists a Hamiltonian action of the Lie algebra $`𝒫_n`$ on the symplectic manifold $`\stackrel{~}{N}_r`$. This action induces an action on $`M_r`$. ¿From Lemma 3.1 and Proposition 4.1 we see that if we can find a finite-dimensional representation of the Lie subalgebra of $`C^{\mathrm{}}(M_r)`$ generated by $`\{f_{ij},1i<jn\}`$ then we will get a representation of $`P_n`$. As explained in the introduction we obtain such a representation on $`T_e(M_r)`$ for a degenerate $`n`$-gon $`e`$. ## 5 Linearization of the bending fields at degenerate polygons. This section is the heart of the paper. We compute $`A_{ij}End(T_e(M_r))`$, the linearization of the bending field $`B_{ij}`$ at a degenerate polygon $`eM_r`$. Now assume that $`e`$ is degenerate, so we may write $$e=(r_1ϵ_1u,\mathrm{},r_nϵ_nu)$$ for some vector $`uS^2`$ and $`ϵ_i=\pm 1`$. Let $`M`$ be a manifold, $`mM`$. We recall the definition of the linearization $`A_XEnd(T_m(M))`$ of a vector field $`X`$ at a point $`m`$ where $`X(m)=0`$. Choose a connection $``$ on $`T(M)`$. Let $`uT_m(M)`$, then $$A_X(u):=(_uX)(m)$$ Since $`X(m)=0`$, $`A_X`$ is independent of the choice of connection. For the case in hand the above definition must be modified since $`M_r`$ is singular at $`e`$. There is a commutative algebra version of the above construction that goes as follows. Assume $`M`$ is a real affine variety, $`mM`$ and $`X`$ is a vector field on $`M`$ satisfying $`X(m)=0`$. Let $`𝔪`$ be the maximal ideal of $`m`$. Then (since $`X(m)=0`$) we have $`X𝔪𝔪`$ whence $`X𝔪^2𝔪^2`$ and $`X`$ induces an element of $`End(𝔪/𝔪^2)=End(T_m^{}(M))`$. By duality we obtain $`A_XEnd(T_m(M))`$. The reader will verify that if $`m`$ is a smooth point of $`M`$ then the two definitions coincide. We now compute the linearization of $`B_{ij}`$ at $`e`$ in $`M_r`$. Recall that we have a diagram $$\begin{array}{ccc}\stackrel{~}{M}_r& & \stackrel{~}{N}_r\\ & & \\ M_r& & N_r\end{array}$$ where $`\stackrel{~}{N}_r=S^2(r_1)\times \mathrm{}\times S^2(r_n)`$ and $`\stackrel{~}{M}_r=\{e\stackrel{~}{N}_r:_{i=1}^ne_i=0\}`$. Define $`g_{ij}:\stackrel{~}{N}_r`$ by $`g_{ij}(e)=e_i+e_j^2`$. Hence $`g_{ij}\|\stackrel{~}{M}_r`$ is $`SO(3)`$-invariant and descends to the function $`f_{ij}`$ on $`M_r`$. Let $`\stackrel{~}{B}_{ij}`$ be the Hamiltonian vector field of $`g_{ij}`$. Then $$\stackrel{~}{B}_{ij}(e)=(0,\mathrm{},e_j\times e_i,0,\mathrm{},e_i\times e_j,0,\mathrm{})$$ and hence $`\stackrel{~}{B}_{ij}`$ vanishes at $`e`$ and is tangent to $`\stackrel{~}{M}_r`$. The induced field on $`\stackrel{~}{M}_r`$ will be denoted $`B_{ij}^{}`$. Then $`B_{ij}^{}`$ projects to $`B_{ij}`$ on $`M_r`$. We note $`dimT_e(\stackrel{~}{N}_r)=2n`$, $`dimT_e(\stackrel{~}{M}_r)=2n2`$ and $`dimT_e(M_r)=2n4`$. ###### Remark 5.1 Since $`e`$ is a singular point of $`M_r`$ we have $$dimT_e(M_r)=2n4>dimM_r=2n6.$$ We will first compute the linearization of $`\stackrel{~}{B}_{ij}`$ at $`e`$ in $`\stackrel{~}{N}_r`$ ($`e`$ is a smooth point on $`\stackrel{~}{N}_r`$ so we use the first procedure) to obtain $`\stackrel{~}{A}_{ij}End(T_e(\stackrel{~}{N}_r))`$. Then $`\stackrel{~}{A}_{ij}`$ will preserve the subspace $`T_e(\stackrel{~}{M}_r)T_e(\stackrel{~}{N}_r)`$ whence we obtain an induced element $`A_{ij}^{}End(T_e(\stackrel{~}{M}_r))`$. But there is an exact sequence $$V_eT_e(\stackrel{~}{M}_r)T_e(M_r)$$ where $`V_e=\{\delta :v^3\text{ such that }\delta _i=e_i\times v,1in\}`$, we note that $`dimV_e=2`$. We will verify that $`A_{ij}^{}(V_e)V_e`$ (in fact $`A_{ij}(V_e)=0`$). Hence $`A_{ij}^{}`$ will descend to $`T_e(M_r)`$. The resulting element of $`End(T_e(M_r))`$ will be $`A_{ij}`$, the linearization of $`B_{ij}`$ at $`e`$. Accordingly we begin by computing the linearization $`\stackrel{~}{A}_{ij}`$ of $`\stackrel{~}{B}_{ij}`$ on $`T_e(\stackrel{~}{N}_r)`$. Thus $`\stackrel{~}{A}_{ij}`$ will be $`2n\times 2n`$ matrix (instead of a $`2n4\times 2n4`$ matrix). Another advantage in passing to $`\stackrel{~}{N}_r`$ is that $`T_e(\stackrel{~}{N}_r)`$ is now a direct sum of the tangent bundles of the factors $$T_e(\stackrel{~}{N}_r)=_{i=1}^nT_{e_i}(S^2(r_i)).$$ The Riemannian connection on $`\stackrel{~}{N}_r`$ is a direct sum of the Riemannian connections on the summands. Thus we may write (for $`\delta T_e(\stackrel{~}{N}_r)`$) $$\stackrel{~}{A}_{ij}(\delta )=(0,\mathrm{},_\delta (e_j\times e_i),0,\mathrm{},_\delta (e_i\times e_j),0,\mathrm{}).$$ We will suppress the zeroes in the above row vectors henceforth. ###### Lemma 5.2 $$\stackrel{~}{A}_{ij}(\delta )=(u\times \delta _i,u\times \delta _j)\left[\begin{array}{cc}ϵ_jr_j& ϵ_jr_j\\ ϵ_ir_i& ϵ_ir_i\end{array}\right]$$ Proof: In the above formula for $`\stackrel{~}{A}_{ij}(\delta )`$ we use the Riemannian connection $``$ on $`S^2`$. We will compute using the flat connection $`\overline{}`$ on $`T(^3)\|S^2`$ and then project back into $`T(S^2)`$ to get $``$. We have $$\overline{}_\delta (e_j\times e_i)=\delta _j\times e_i+e_j\times \delta _i$$ $$\overline{}_\delta (e_i\times e_j)=\delta _i\times e_j+e_i\times \delta _j.$$ Evaluating at $`e`$ we obtain $$\overline{}_\delta (e_j\times e_i)\|_e=ϵ_ir_i\delta _j\times u+ϵ_jr_ju\times \delta _i=ϵ_jr_ju\times \delta _iϵ_ir_iu\times \delta _j.$$ Since the right-hand side is in $`T_e(S^2)`$ we have also $$_\delta (e_j\times e_i)\|_e=ϵ_jr_ju\times \delta _iϵ_ir_iu\times \delta _j.$$ Finally $`_\delta (e_i\times e_j)\|_e=_\delta (e_j\times e_i)\|_e`$ and the lemma follows. $`\mathrm{}`$ We now relate the action of $`𝒫_n`$ on $`T_e(\stackrel{~}{N}_r)`$ we have just computed to the action on $`T_e(M_r)`$. We recall that $`\stackrel{~}{M}_r=\{e\stackrel{~}{N}_r:_{i=1}^ne_i=0\}`$ whence $`T_e(\stackrel{~}{M}_r)=\{\delta T_e(\stackrel{~}{N}_r):_{i=1}^n\delta _i=0\}`$. We have the 2-dimensional subspace $`V_e`$ of tangents to the $`SO(3)`$-orbit through $`e`$ described above. Thus we have a filtration $`F_{}`$ given by $$V_eT_e(\stackrel{~}{M}_r)T_e(\stackrel{~}{N}_r)$$ and a canonical isomorphism $$T_e(\stackrel{~}{M}_r)/V_eT_e(M_r).$$ We now show that $`𝒫_n`$ preserves the above filtration. ###### Lemma 5.3 1. $`𝒫_nT_e(\stackrel{~}{N}_r)T_e(\stackrel{~}{M}_r)`$. 2. $`𝒫_nV_e=0`$. Proof: (1) is immediate. We prove (2). Suppose $`\delta V_e`$. We claim $$\delta _j\times e_i+e_j\times \delta _i=0,1i<jn.$$ Indeed, $$\delta _j\times e_i+e_j\times \delta _i=(e_j\times v)\times e_i+e_j\times (e_i\times v)=(e_j\times v)\times e_i+(e_j\times e_i)\times v+e_i\times (e_j\times v).$$ But $`e`$ is degenerate, so $`e_i\times e_j=0`$. $`\mathrm{}`$ We collect our results in ###### Theorem 5.4 1. There is a $`𝒫_n`$-stable filtration $$V_eT_e(\stackrel{~}{M}_r)T_e(\stackrel{~}{N}_r).$$ 2. $`T_e(M_r)T_e(\stackrel{~}{M}_r)/V_e`$. 3. There is an isomorphism $$\varphi :T_e(\stackrel{~}{N}_r)T_u(S^2)^n$$ such that $`\varphi \rho \varphi ^1(X_{ij})=aduJ_{ij}(ϵ_ir_i,ϵ_jr_j)`$. 4. $`\varphi (T_e(\stackrel{~}{M}_r))=T_u(S^2)_0^n`$ and $`\varphi (V_e)=T_u(S^2)v(ϵ,r)`$. Here $`_0^n=\{(x_1,\mathrm{},x_n):_{i=1}^nx_i=0\}`$ and $`v(ϵ,r)=(ϵ_1r_1,\mathrm{},ϵ_nr_n)`$. Here $`^n`$ is realized as the space of row vectors with $`n`$ components. ## 6 The action on the holomorphic tangent space. The point of this section is that $`T_e(\stackrel{~}{N}_r)`$ has a $`𝒫_n`$-invariant almost complex structure that descends to $`T_e(M_r)`$. We will compute the corresponding action of $`𝒫_n`$ on the holomorphic tangent space. Define an almost complex structure $`JEnd(T_e(\stackrel{~}{N}_r))`$ by $$J(\delta )=\eta \text{ such that }\eta _i=u\times \delta _i,1in.$$ ###### Lemma 6.1 1. $`J`$ is $`𝒫_n`$-invariant. 2. The filtration $`F_{}`$ is invariant under $`J`$. Proof: The first assertion is immediate. It is also clear that $`T_t(\stackrel{~}{M}_r)`$ is invariant under $`J`$. It remains to check that $`V_e`$ is invariant under $`J`$. Suppose $`\delta V_e`$. Hence there exists $`v^3`$ such that $`\delta _i=ϵ_ir_iu\times v`$, $`1in`$. Then $`J\delta _i=u\times (ϵ_ir_iu\times v)=ϵ_ir_iu\times (u\times v)`$. Hence if we put $`w=u\times v`$ then $$J\delta _i=ϵ_ir_ru\times w,1in.$$ Therefore $`J\delta V_e`$. $`\mathrm{}`$ ###### Remark 6.2 The almost complex structure $`J`$ is not the one induced by the complex structure on $`\stackrel{~}{N}_r=_{i=1}^nS^2(r_i)`$. We have changed the complex structure on $`S^2(r_i)`$ to its conjugate for each $`i`$ such that $`e_i`$ is a back-track (i.e. $`ϵ_i=1`$). We can decompose $`T_e(\stackrel{~}{N}_r)`$ into the $`+i`$-eigenspace of $`J`$ denoted by $`T_e^ϵ(\stackrel{~}{N}_r)`$ and the $`i`$-eigenspace denoted by $`T_e^ϵ(\stackrel{~}{N}_r)`$. Accordingly we have $$T_e^ϵ(\stackrel{~}{N}_r)=\{\delta T_e(\stackrel{~}{N}_r):u\times \delta _j=\sqrt{1}\delta _j\}$$ Similarly we denote the $`+i`$-eigenspaces of $`J`$ acting on $`T_e(\stackrel{~}{M}_r)`$ and $`V_e`$ by $`T_e^ϵ(\stackrel{~}{M}_r)`$ and $`V_e^ϵ`$ respectively. We denote the quotient $`T_e^ϵ(\stackrel{~}{M}_r)/V_e^ϵ`$ by $`T_e^ϵ(M_r)`$. Clearly the latter space is the $`+i`$-eigenspace of $`J`$ acting on $`T_e(M_r)`$. Now we recall that we have an isomorphism $$\varphi :T_e(\stackrel{~}{N}_r)T_u(S^2)^n$$ complexifying we obtain $$\varphi :T_e(\stackrel{~}{N}_r)T_u(S^2)_{}^n.$$ We see that $`\varphi `$ conjugates $`J`$ to $`adu1`$ and we have an induced isomorphism (again denoted by $`\varphi `$) $$\varphi :T_e^ϵ(\stackrel{~}{N}_r)T_u^{1,0}(S^2)_{}^n.$$ Under $`\varphi `$ the action of $`X_{ij}`$ transforms to $`\sqrt{1}IJ_{ij}(ϵ_ir_i,ϵ_jr_j)`$. We note that $`dim_{}T_u^{1,0}(S^2)=1`$ and we obtain a canonical isomorphism $$\psi :T_e^ϵ(\stackrel{~}{N}_r)^n.$$ This isomorphism has the property: $$\psi (T_e^ϵ(\stackrel{~}{M}_r))=_0^n,\psi (V_e^ϵ)=v(ϵ,r).$$ We have completed our computation of the action of $`𝒫_n`$. ###### Theorem 6.3 1. There is a canonical isomorphism $`\psi :T_e^ϵ(\stackrel{~}{N}_r)^n`$ 2. $`\psi `$ induces the action of $`X_{ij}𝒫_n`$ on $`^n`$ by $`\sqrt{1}J_{ij}(ϵ_ir_i,ϵ_jr_j)`$. 3. $`^n`$ admits a $`𝒫_n`$-invariant filtration by $`\psi (T_e^ϵ(\stackrel{~}{M}_r))=_0^n`$, $`\psi (V_e^ϵ)=v(ϵ,r)`$. 4. There is an $`𝒫_n`$-invariant complex structure $`J`$ on $`T_e(M_r)`$. The induced action of $`𝒫_n`$ on the $`+i`$-eigenspace of $`J`$ in $`T_e(M_r)`$ corresponds to the action of $`𝒫_n`$ on the quotient $`_0^n/v(ϵ,r)`$. Here $`^n`$ is realized as the space of row vectors with $`n`$ components. ## 7 The associated hypergeometric equation. As discussed in the introduction we use the linear operators $`A_{ij}End(^n)`$ to obtain a flat holomorphic connection $``$ on the trivial $`T_e^ϵ(\stackrel{~}{N}_r)`$-bundle $``$ over $`=_{}^n`$. The connection form $`\omega `$ of $``$ is $$\omega =\underset{1i<jn}{}\frac{dz_idz_j}{z_iz_j}A_{ij}.$$ A (multivalued) holomorphic section of $``$ corresponds to a row vector $`F=(F_1,\mathrm{},F_n)`$ of (multivalued) holomorphic functions. The hypergeometric equation comes from the condition that $`F`$ be parallel for the connection $``$: $$dF=F\omega $$ or equivalently $$dF_i=\underset{j,ji}{}(\lambda _jF_i\lambda _iF_j)\frac{dz_idz_j}{z_iz_j}$$ (1) with $`\lambda _j=\sqrt{1}ϵ_jr_j`$. We will refer to (1) as the hypergeometric equation. We observe that the operators $`A_{ij}`$ leave invariant the subspace $`_0^n`$ and annihilate the line $`V_\lambda =(\lambda _1,\mathrm{},\lambda _n)`$. We obtain a diagram of flat bundles over $`_{}^n`$: $$\begin{array}{ccc}_{}^n\times _0^n& & _{}^n\times ^n\\ & & \\ _{}^n\times _0^n/V_\lambda & & \end{array}$$ The monodromies of these bundles will be the representations of $`P_n`$ corresponding to the actions of $`𝒫_n`$ on $`T_e(\stackrel{~}{N}_r)`$, $`T_e(\stackrel{~}{M}_r)`$, $`T_e(M_r)`$. ## 8 Solving the hypergeometric equation by hypergeometric integrals. Let $`\lambda _1,\mathrm{},\lambda _n`$ be complex numbers with $`\lambda _j,1jn`$. Let $`(\xi ,z_1,\mathrm{},z_n)(^{n+1})_{}`$ and $`\mathrm{\Phi }(\xi ,z_1,\mathrm{},z_n)`$ be the hypergeometric integrand $$\mathrm{\Phi }(\xi ,z_1,\mathrm{},z_n):=(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}.$$ Let $`\chi :=\chi _\lambda :𝔽_n^{}`$ be the character defined by $`\chi (\gamma _j)=\mathrm{exp}(2\pi \sqrt{1}\lambda _j)`$, $`1jn`$. Recall that $`\{\gamma _1,\mathrm{},\gamma _n\}`$ is a generating set for $`𝔽_n`$, the free group of rank $`n`$. Here we identify $`𝔽_n`$ with the fundamental group $`\pi _1(M,b)`$, where $`M=\{z_1,\mathrm{},z_n\}`$, so that the conjugacy class of $`\gamma _j`$ is represented by a sufficiently small loop which goes once around $`z_j`$ in the counterclockwise direction. Note that $`\chi (\gamma _j)1`$, $`1jn`$. For any character $`\chi :𝔽_n^{}`$ we let $`L_\chi `$ be the local system over $`M`$ given by $$L_\chi =\stackrel{~}{M}\times /((x,z)(\gamma x,\chi (\gamma )z)).$$ We define a multivalued parallel section $`\sigma `$ of $`L_\chi `$ by $`\sigma (x)=[x,1]`$ (where $`[x,z]`$ denotes the equivalence class of $`(x,z)`$). Note that the lift of $`\sigma `$ to the universal cover satisfies $$\sigma (\gamma x)=[\gamma x,1]=[x,\chi (\gamma )^1]=\chi (\gamma )^1\sigma (x).$$ The following lemma is obvious: ###### Lemma 8.1 The $`L_\chi `$-valued 1-forms $`\zeta _j`$, $`1jn`$, defined by $$\zeta _j(\xi )=(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\frac{d\xi }{\xi z_j}\sigma $$ are single-valued on $`M`$. Hence $`\zeta _j`$ gives rise to a class $`[\zeta _j]`$ in the de Rham cohomology group $`\text{H}_{dR}^1(M,L_\chi )`$. Let $`\gamma \text{H}_1(M,L_{\chi ^1})`$. Let $`G_j`$ be the Kronecker pairing $`\zeta _j,\gamma `$ considered as a function of $`z_1,\mathrm{},z_n`$. This Kronecker pairing is traditionally represented as an integral. To make this precise let $`\gamma =_{i=1}^ka_i\tau _i`$, where each $`a_i`$, $`1ik`$, is a 1-simplex and $`\tau _i`$ is a parallel section of $`^1\|a_i`$. Then $`\zeta _j,\gamma `$ is given by $$G_j(z_1,..,z_n)=\underset{i=1}{\overset{k}{}}_{a_i}(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\sigma ,\tau _i\frac{d\xi }{\xi z_j}.$$ We will use the following more economical notation: $$G_j(z_1,..,z_n)=_\gamma (\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\frac{d\xi }{\xi z_j}\sigma .$$ Now we let $`z=(z_1,\mathrm{},z_n)`$ vary. Let $`\pi :_{}^{n+1}_{}^n`$ be the map that forgets the first component. Then $`\pi ^1(z)`$ is isomorphic to $`\{z_1,\mathrm{},z_n\}`$. By \[DM, 3.13\], the flat line bundle $`L_\chi `$ on $`\pi ^1(z)`$ is the restriction of a flat line bundle $`\stackrel{~}{L}_\chi `$ on $`_{}^{n+1}`$. As $`z`$ varies, the forms $`\zeta _1,\mathrm{},\zeta _n`$ give rise to relative holomorphic 1-forms on $`_{}^{n+1}`$ with coefficients in $`\stackrel{~}{L}_\chi `$. We recall that a relative holomorphic form on the total space $`E`$ of a holomorphic fiber bundle $`p:EB`$ is an element of the quotient differential graded algebra $$\mathrm{\Omega }^{}(E)/(p^{}\mathrm{\Omega }^{}(B)^+).$$ Here $`\mathrm{\Omega }^q`$ denotes the holomorphic $`q`$-forms and $`(p^{}\mathrm{\Omega }^{}(B)^+)`$ denotes the differential ideal in $`\mathrm{\Omega }^{}(E)`$ generated by the pull-backs to $`E`$ of holomorphic forms on $`B`$ of positive degree. A relative holomorphic $`q`$-form $`\eta `$ is relatively closed if $`d\eta `$ is in the above ideal. The forms $`\zeta _1,\mathrm{},\zeta _n`$ are relatively closed, hence they induce holomorphic sections $`[\zeta _1],\mathrm{},[\zeta _n]`$ of the vector bundle $`^1`$ over $`_{}^n`$ with fiber over $`z`$ given by $$\text{H}^1(\pi ^1(z),\stackrel{~}{L}_\chi \|\pi ^1(z)).$$ Precisely, $`[\zeta _i](z)`$ is the class of the 1-form $`\zeta _i(z)`$ on $`\pi ^1(z)`$ in the above cohomology group. The bundle $`^1`$ has a flat connection, the Gauss-Manin connection, whose definition we now recall. Note first that a local trivialization of $`\pi `$ induces a local trivialization of $`^1`$. Then a smooth section of $`^1`$ is parallel for the Gauss-Manin connection if it is constant when expressed in terms of all such induced local trivializations. The bundle $`_1`$ of the first homology groups with coefficients in $`\stackrel{~}{L}_{\chi ^1}`$ admits an analogous flat connection. Now let $`p:\stackrel{~}{}_{}^n_{}^n`$ denote the universal cover of $`_{}^n`$. We obtain a pull-back fiber bundle $`\stackrel{~}{\pi }:E\stackrel{~}{}_{}^n`$ of $`n`$-punctured complex lines over $`\stackrel{~}{}_{}^n`$ and pull-back flat vector bundles $`\stackrel{~}{}^1`$ and $`\stackrel{~}{}_1`$. Choose a base-point $`z^0=(z_1^0,\mathrm{},z_n^0)`$ in $`_{}^n`$. We use $`M`$ to denote $`\{z_1^0,\mathrm{},z_n^0\}`$ henceforth. Choose a base-point $`\stackrel{~}{z}^0`$ in $`\stackrel{~}{}_{}^n`$ lying over $`z^0`$. We may identify the fiber of $`\stackrel{~}{}_1`$ over $`\stackrel{~}{z}^0`$ with $`\text{H}_1(M,L_{\chi ^1})`$. Hence given $`\gamma \text{H}^1(M,L_{\chi ^1})`$ there is a unique parallel section $`\stackrel{~}{\gamma }`$ of $`\stackrel{~}{}_1`$ such that $`\stackrel{~}{\gamma }(\stackrel{~}{z}^0)=\gamma `$. We can now define a global holomorphic function $`G_j(z)`$ on $`\stackrel{~}{}_{}^n`$ by $$G_j(z)=_{\stackrel{~}{\gamma }}(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\frac{d\xi }{\xi z_j}\sigma .$$ Here we have used the same notation for corresponding (under pull-back) objects on $`_{}^n`$ and $`\stackrel{~}{}_{}^n`$. We may also write $$G_j(z)=[\zeta _j(z)],\stackrel{~}{\gamma }$$ where $`,`$ is the fiberwise pairing between $`\stackrel{~}{}^1`$ and $`\stackrel{~}{}_1`$. We have ###### Lemma 8.2 $$dG_i(z)=\underset{j=1}{\overset{n}{}}(_\gamma \frac{}{z_j}(\frac{\mathrm{\Phi }}{\xi z_i})𝑑\xi \sigma )dz_j.$$ Proof: We have $$dG_i(z)=[\zeta _i(z)],\stackrel{~}{\gamma }$$ where $``$ is the Gauss-Manin connection. We will need another formula for the Gauss-Manin connection, see \[KO\] or Remark 8.3 below. Before stating the formula we need more notation. Let $`F^q\mathrm{\Omega }^q(E)`$ denote the subspace of holomorphic $`q`$-forms on $`E`$ that are multiples of pull-backs of $`q`$-forms from the base $`\stackrel{~}{}_{}^n`$ by elements of $`𝒪(E)`$. Then we have a canonical isomorphism (because the fibers of $`\stackrel{~}{\pi }`$ have complex dimension $`1`$) $$\frac{\mathrm{\Omega }^2(E)}{dF^1\mathrm{\Omega }^1(E)+F^2\mathrm{\Omega }^2(E)}\mathrm{\Omega }^1(\stackrel{~}{}_{}^n,\stackrel{~}{}^1).$$ Now the formula for $``$ is $$[\zeta _i]=[d\zeta _i].$$ Here $`d\zeta _i`$ denotes the exterior differential of $`\zeta _i`$ where $`\zeta _i`$ is considered as a 1-form on $`E`$ (modulo $`F^1\mathrm{\Omega }^1(E)`$) with values in the line bundle $`p^{}\stackrel{~}{L}_\chi `$. The symbol $`[d\zeta _i]`$ denotes the class of $`d\zeta _i`$ modulo $`dF^1\mathrm{\Omega }^1(E)+F^2\mathrm{\Omega }^2(E)`$. The lemma follows from the formula $$d\zeta _i\underset{j=1}{\overset{n}{}}\frac{}{z_j}(\frac{\mathrm{\Phi }}{\xi z_i})dz_jd\xi \sigma $$ together with the observation that integration over $`\stackrel{~}{\gamma }`$ factors through $`[]`$. $`\mathrm{}`$ ###### Remark 8.3 The above formula for $``$ can be proved as follows. First note that the formula does indeed define a connection, to be denoted $`^{}`$ on $`^1`$. To show that $``$ and $`^{}`$ agree it suffices to show they agree locally. Since they are both invariantly defined it suffices to prove that they agree on trivial bundles. But it is clear that in this case a section of $`^1`$ is parallel for $`^{}`$ if and only if it is constant. The proof of the next lemma is a modification of \[K1, Proposition 2.2.2\]. ###### Lemma 8.4 The functions $`G=(G_1,\mathrm{},G_n)`$ satisfy $$dG_i=\underset{j,ji}{}(\lambda _jG_i\lambda _jG_j)\frac{dz_idz_j}{z_iz_j}\sigma ordG^T=\omega G^T.$$ Proof: We will drop the $`\sigma `$ for the course of the proof: $$G_i(z)=_\gamma \mathrm{\Phi }\frac{d\xi }{\xi z_i}.$$ Whence by Lemma 8.2 $$dG_i=\underset{j=1}{\overset{n}{}}[_\gamma \lambda _j\mathrm{\Phi }(\xi z_j)^1(\xi z_i)^1𝑑\xi ]dz_j+[_\gamma \mathrm{\Phi }(\xi z_i)^2𝑑\xi ]dz_i$$ $$=\underset{ji}{}[_\gamma \lambda _j\mathrm{\Phi }(\xi z_j)^1(\xi z_i)^1𝑑\xi ]dz_j[_\gamma (\lambda _i1)\mathrm{\Phi }(\xi z_i)^2𝑑\xi ]dz_i.$$ We simplify the first term using $$\frac{1}{\xi z_i}\frac{1}{\xi z_j}=\frac{1}{z_iz_j}(\frac{1}{\xi z_i}\frac{1}{\xi z_j})$$ to obtain $$=\underset{ji}{}\frac{\lambda _j}{z_iz_j}[_\gamma \mathrm{\Phi }\frac{d\xi }{\xi z_i}_\gamma \mathrm{\Phi }\frac{d\xi }{\xi z_j}]dz_j[_\gamma (\lambda _i1)\mathrm{\Phi }(\xi z_i)^2𝑑\xi ]dz_i=$$ $$=\underset{ji}{}\frac{\lambda _jG_i}{z_iz_j}dz_j+\underset{ji}{}\frac{\lambda _jG_j}{z_iz_j}dz_j[_\gamma (\lambda _i1)\mathrm{\Phi }(\xi z_i)^2𝑑\xi ]dz_i.$$ Now we have $$d(\mathrm{\Phi }(\xi z_i)^1)=(\lambda _i1)\mathrm{\Phi }(\xi z_i)^2d\xi +\underset{ji}{}\lambda _j\mathrm{\Phi }(\xi z_i)^1(\xi z_j)^1d\xi .$$ Thus by Stokes’ Theorem $$_\gamma (\lambda _i1)\mathrm{\Phi }(\xi z_i)^2𝑑\xi =_\gamma \underset{ji}{}\lambda _j\mathrm{\Phi }(\xi z_i)^1(\xi z_j)^1d\xi =$$ $$_\gamma \underset{ji}{}\lambda _j\mathrm{\Phi }\frac{1}{z_iz_j}(\frac{1}{\xi z_i}\frac{1}{\xi z_j})d\xi =$$ $$\underset{ji}{}\frac{\lambda _j}{z_iz_j}G_i\underset{ji}{}\frac{\lambda _j}{z_iz_j}G_j$$ hence $$[_\gamma (\lambda _i1)\mathrm{\Phi }(\xi z_i)^2𝑑\xi ]dz_i=\underset{ji}{}\frac{dz_i}{z_iz_j}(\lambda _jG_i\lambda _jG_j).$$ We obtain $$dG_i=\underset{ji}{}\frac{dz_idz_j}{z_iz_j}(\lambda _jG_i\lambda _jG_j).\text{ }\text{ }\mathrm{}$$ ###### Remark 8.5 The simplification using Stokes’ Theorem above is equivalent to observing that $$\mathrm{\Phi }(\xi z_i)^1dz_i\sigma F^1\mathrm{\Omega }^1(E),1in,$$ and we work modulo $`dF^1\mathrm{\Omega }^1(E)`$ in computing $``$. We now define $`F_i:=\lambda _iG_i`$, $`1in`$. ###### Lemma 8.6 $`F=(F_1,\mathrm{},F_n)`$ is a solution of the hypergeometric equation (1). Proof: $$dF_i=\lambda _idG_i=\underset{ji}{}\frac{dz_idz_j}{z_iz_j}(\lambda _i\lambda _jG_i\lambda _i\lambda _jG_j)=$$ $$=\underset{ji}{}\lambda _j(\lambda _iG_i)\lambda _i(\lambda _jG_j)\frac{dz_idz_j}{z_iz_j}=$$ $$=\underset{ji}{}(\lambda _jF_i\lambda _iF_j)\frac{dz_idz_j}{z_iz_j}.\text{ }\text{ }\mathrm{}$$ We have proved ###### Theorem 8.7 Let $`\gamma `$ be an element of $`\text{H}_1(M,L_{\chi ^1})`$ and $`\sigma `$ a flat multivalued section of $`L_\chi `$. For $`\lambda =(\lambda _1,\mathrm{},\lambda _n)^n`$ define a holomorphic function on $`\stackrel{~}{}_{}^n`$ by $$F_i:=\lambda _i_{\stackrel{~}{\gamma }}(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\frac{d\xi }{\xi z_i}\sigma .$$ Then $`F=(F_1,\mathrm{},F_n)`$ is a solution of the hypergeometric equation. ## 9 The monodromy representation of the hypergeometric equation and the action on homology. We have seen that for $`\gamma \text{H}_1(M,L_{\chi ^1})`$ we obtain a solution $`S=(F_1,\mathrm{},F_n)`$ of the hypergeometric equation by the formula $$F_i:=\lambda _i_{\stackrel{~}{\gamma }}(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\frac{d\xi }{\xi z_i}\sigma .$$ It is important to recall that $`_{j=1}^n\lambda _j=0`$. The differential forms $$\eta _j=\lambda _j(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\frac{d\xi }{\xi z_i}\sigma $$ are de Rham representatives of the cohomology classes $`[\eta _j],1jn,`$ in $`\text{H}^1(M,L_{\chi ^1})`$. Note that $$d((\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\sigma )=\eta _1+\mathrm{}\eta _n$$ hence we have the relation $$[\eta _1]+\mathrm{}+[\eta _n]=0$$ (2) ###### Lemma 9.1 The span of the cohomology classes $`[\eta _j],1jn`$, has dimension $`n1`$. Proof: First since $`_{j=1}^n\lambda _j=0`$ we have $`\chi (\gamma _1\gamma _2\mathrm{}\gamma _n)=1`$. Thus $`L_\chi `$ extends to a flat line bundle over $`^1\{z_1,\mathrm{},z_n\}`$. Also, $`\eta _j`$ extends meromorphically over infinity with a simple pole at infinity. Next we extend the flat line bundle $`L_\chi `$ to a holomorphic line bundle $`^{hol}`$ on $`^1`$ so that $`(\xi z_j)^{\lambda _j}\sigma `$ is a local basis around $`z_j`$. Then $`(\xi z_1)^{\lambda _1}\mathrm{}(\xi z_n)^{\lambda _n}\sigma `$ is a holomorphic section of $`^{hol}`$ which has no zeroes or poles. We can now prove the lemma. We have a flat line bundle $`L_\chi `$ over $`M`$ (with trivial monodromy around $`\mathrm{}`$). The argument of \[DM, §2.7\] proves that we can compute the group $`\text{H}^1(M,L_\chi )`$ as the 1-st cohomology group of the complex $`(\mathrm{\Omega }^{}(^1,D,L_\chi ),d)`$ of holomorphic $`L_\chi `$-valued forms on $`M`$ which have at worst poles at $`z_1,\mathrm{},z_n,\mathrm{}`$. Here the (additive) divisor $`D`$ is defined by $`D=z_1+\mathrm{}+z_n+\mathrm{}`$. Now $`\eta _j\mathrm{\Omega }^1(^1,D,L_\chi )`$ and $$\mathrm{\Omega }^0(^1,D,L_\chi )=\{f\mathrm{\Phi }\sigma :\text{ so that }f\text{ has at worst poles at }D\}.$$ First note that $`Span(\eta _1,\mathrm{},\eta _n)\mathrm{\Omega }^1(^1,D,L_\chi )`$ has dimension $`n`$ since the forms $`\eta _j`$ have singularities at distinct points of $``$. Suppose that there exists $`f\mathrm{\Phi }\sigma \mathrm{\Omega }^0(^1,D,L_\chi )`$ and $`c_1,\mathrm{},c_n`$ such that $$d(f\mathrm{\Phi }\sigma )=c_1\eta _1+\mathrm{}+c_n\eta _n.$$ We claim that $`f`$ cannot have any poles. Indeed, assume $`f`$ has a pole of order $`k1`$ at $`z_i`$. Then $$f(\xi )=\frac{c}{(\xi z_i)^k}+\mathrm{}$$ We are assuming $$df\mathrm{\Phi }+fd\mathrm{\Phi }=\underset{i=1}{\overset{n}{}}c_i\eta _i$$ or $$df\mathrm{\Phi }+(f\underset{i=1}{\overset{n}{}}\frac{\lambda _i}{\xi z_i}d\xi )\mathrm{\Phi }=\underset{i=1}{\overset{n}{}}c_i\eta _i.$$ (3) Equating the coefficients of $`(\xi z_i)^{k1}`$ in the equation (3) from each side we obtain $`kc+\lambda _ic=0`$, or $`\lambda _i=k`$. This contradicts the assumption that each $`\lambda _i`$ is pure imaginary. It remains to check that $`f`$ is not a polynomial. Assume $`f`$ has a pole of order $`k1`$ at $`\mathrm{}`$, whence $`f(\xi )=a_0+a_1\xi +\mathrm{}+a_k\xi ^k`$. We equate the coefficients at $`\xi ^{k1}d\xi `$ on each side of (3) to obtain $`ka_k+(_{i=1}^n\lambda _i)a_k=0`$ or $`ka_k=0`$. This contradiction proves the claim. Hence $`fc`$ and hence $$df=c\underset{i=1}{\overset{n}{}}\eta _i$$ which means that the dimension of the subspace of coboundaries in $`Span(\eta _1,\mathrm{},\eta _n)`$ is $`1`$. $`\mathrm{}`$ In the group cohomology computations that follow $`\gamma _1,\mathrm{},\gamma _n`$ will be a generating set of $`𝔽_n`$ and $`b_1,\mathrm{},b_n`$ will be its image under abelianization in $`^n`$. Here the loop representing $`\gamma _i`$ is obtained by connecting the small circle $`a_i`$ going around $`z_i`$ to the base-point $`\{z_1,\mathrm{},z_n\}`$. We recall that $`P_n`$ acts on $`𝔽_n`$ preserving the conjugacy classes of the generators $`\gamma _j`$. Hence the induced action on $`^n`$ is trivial and $`P_n`$ fixes any character $`\chi :𝔽_n^{}`$. Hence $`P_n`$ acts on $`\text{H}^1(𝔽_n,_\chi )`$. Here we let $`_\chi `$ denote the 1-dimensional space on which $`𝔽_n`$ acts via $`\chi `$. We next need ###### Lemma 9.2 Suppose that $`\chi :𝔽_n^{}`$ satisfies $`\chi (\gamma _i)1`$ for all $`i`$. Then $`dim_{}\text{H}^1(𝔽_n,_\chi )=n1`$. Proof: The Euler characteristic $`E(𝔽_n,_1)=1n`$. Hence $`E(𝔽_n,_\chi )=1n`$. On the other hand, $`\text{H}^0(𝔽_n,_\chi )=0`$. $`\mathrm{}`$ ###### Corollary 9.3 $`dim_{}\text{H}_1(M,L_{\chi ^1})=n1`$ and the classes $`[\eta _1],\mathrm{},[\eta _{n1}]`$ form a basis for $`\text{H}^1(M,L_\chi )`$ We can construct an explicit basis $`w_1,\mathrm{},w_{n1}`$ for $`\text{H}_1(M,^1)`$ following \[DM, §2\] as follows. We write $`w_i=\gamma _i\sigma _i+\gamma _{i+1}\sigma _{i+1}`$, where $`\sigma _i,\sigma _{i+1}`$ are multivalued flat sections along $`\gamma _i,\gamma _{i+1}`$ respectively and the jump experienced by $`\sigma _i`$ (at the base-point) after parallel translating along $`\gamma _i`$ cancels that of $`\sigma _{i+1}`$ along $`\gamma _{i+1}`$. Define flat sections $`S_i`$, $`1in1`$, of $`\stackrel{~}{}_{}^n\times _0^n`$ by $$S_i:=(S_{i1},\mathrm{},S_{in}),\text{ where }S_{ij}=\lambda _j_{\stackrel{~}{w}_i}\eta _j.$$ We see then that $`S_1,\mathrm{},S_{n1}`$ are multivalued parallel sections of $`_{}^n\times _0^n`$. The desired representation $`\rho :P_nAut(_0^n)`$ is obtained by parallel translation of $`S_1,\mathrm{},S_{n1}`$ along loops in $`_{}^n`$. The resulting automorphisms leave invariant the line $`\lambda `$ where $`\lambda =(\lambda _1,\mathrm{},\lambda _n)`$. Before stating the main result of this section we need to define a special class $`w_{\mathrm{}}`$ in $`\text{H}_1(M,L_\chi ^1)`$. Let $`a_{\mathrm{}}`$ be a circle whose interior contains all the punctures $`z_1,\mathrm{},z_n`$. Since $`\lambda _1+\mathrm{}\lambda _n=0`$, the monodromy of $`L_\chi ^1`$ around $`a_{\mathrm{}}`$ is trivial. Hence there is a nonzero parallel section $`\sigma ^{}`$ of $`L_\chi ^1\|a_{\mathrm{}}`$. We let $`w_{\mathrm{}}`$ be the homology class represented by $`a_{\mathrm{}}\sigma ^{}`$. Let $`\tau :P_nAut\text{H}_1(M,L_\chi ^1)`$ be the homomorphism induced by the inclusion $`P_nAut(𝔽_n)`$ (recall that $`P_n`$ acts trivially on the sheaf of parallel sections of $`L_\chi ^1`$). ###### Lemma 9.4 (1) $`_w_{\mathrm{}}\eta _i=\lambda _i`$, in particular $`w_{\mathrm{}}0`$. (2) The class $`w_{\mathrm{}}`$ is fixed by $`P_n`$. Proof: To prove (1) we apply the residue theorem and note that $$\mathrm{\Phi }(\xi ,z)\|_{\xi =\mathrm{}}=1$$ and the residue of $`(\xi z_i)^1d\xi `$ at $`\xi =\mathrm{}`$ is $`1`$. To verify (2) we identify $`P_n`$ with a subgroup of the mapping class group of $`M`$. Then we choose representatives for the elements of $`P_n`$ so that they act by the identity on the closure of the exterior of the circle $`a_{\mathrm{}}`$. $`\mathrm{}`$ We now have ###### Theorem 9.5 (i) The monodromy representation of the flat bundle $`_{}^n\times _0^n`$ is equivalent to $`\tau `$. (ii) Under the above equivalence the invariant line $`V_\lambda _0^n`$ corresponds to the line $`w_{\mathrm{}}\text{H}_1(M,_\chi ^1)`$. (iii) We obtain an induced equivalence of the monodromy representation of $`_{}^n\times _0^n/V_\lambda `$ and the induced action of $`P_n`$ on $`\text{H}_1(^1\{z_1,\mathrm{},z_n\},L_\chi ^1)`$. Proof: We have an isomorphism $`\mathrm{\Psi }`$ from $`\text{H}_1(M,L_\chi ^1)`$ onto the space of parallel sections on $`\stackrel{~}{}_{}^n\times _0^n`$ given by $`\mathrm{\Psi }(w)=S_w`$ where $$S_w=(_{\stackrel{~}{w}}\eta _1,\mathrm{},_{\stackrel{~}{w}}\eta _n)=([\eta _1],\stackrel{~}{w},\mathrm{},[\eta _n],\stackrel{~}{w}).$$ We claim that $`\mathrm{\Psi }`$ intertwines the representations $`\tau `$ and $`\rho `$ (see above) of $`P_n`$. The monodromy representation $`\rho :P_nAut(_0^n)`$ is defined by $$S_w(g^1z)=S_w(z)\rho (g).$$ In order to go further we will need to lift the $`P_n`$ action on $`\stackrel{~}{}_{}^n`$ to the total space of $`\stackrel{~}{\pi }:E\stackrel{~}{}_{}^n`$. We note that from the fiber bundle $`\pi :_{}^{n+1}_{}^n`$ we get an exact sequence $`𝔽_nP_{n+1}P_n`$. We may split this sequence by mapping $`P_n`$ to the subgroup of $`P_{n+1}`$ which consists of those elements that do not involve the first string of a braid – recall that $`\pi `$ forgets the first point. Let $`\stackrel{~}{}_{}^{n+1}`$ be the universal cover of $`_{}^{n+1}`$. Then $`P_{n+1}`$ acts on $`\stackrel{~}{}_{}^{n+1}`$. But $`E=\stackrel{~}{}_{}^{n+1}/𝔽_n`$, whence $`P_n=P_{n+1}/𝔽_n`$ acts on $`E`$ as the group of deck transformations of the cover $`E_{}^{n+1}`$, and we obtain the required lift $`\stackrel{~}{g}`$ of elements $`gP_n`$ to $`Aut(E)`$. We now can give a formula for the monodromy representation $`\tau `$, namely $$\stackrel{~}{w}(gz)=\stackrel{~}{g}_{}\tau (g)^1\stackrel{~}{w}(z)$$ or $$\stackrel{~}{w}(g^1z)=\stackrel{~}{g}_{}^1\tau (g)\stackrel{~}{w}(z).$$ We can now prove the claim. Observe that since $`\eta _i`$ is an invariantly defined 1-form with values in $`L_\chi `$ on $`_{}^{n+1}`$ we have $$\eta _i(gz)=(\stackrel{~}{g}^1)^{}\eta _i(z)$$ or $$\eta _i(g^1z)=(\stackrel{~}{g})^{}\eta _i(z).$$ Hence $$S_w(z)\rho (g)=S_w(g^1z)=(_{\stackrel{~}{w}(g^1z)}\eta _1(g^1z),\mathrm{},_{\stackrel{~}{w}(g^1z)}\eta _n(g^1z))=$$ $$(_{g_{}^1\tau (g)\stackrel{~}{w}(z)}\stackrel{~}{g}^{}\eta _1(z),\mathrm{},_{g_{}^1\tau (g)\stackrel{~}{w}(z)}\stackrel{~}{g}^{}\eta _n(z))=$$ $$(_{\tau (g)\stackrel{~}{w}(z)}\eta _1(z),\mathrm{},_{\tau (g)\stackrel{~}{w}(z)}\eta _n(z))$$ and the claim is proved. Hence (i) follows. To verify (ii) it suffices to observe that $`S_w_{\mathrm{}}=(\lambda _1,\mathrm{},\lambda _n)`$, which follows from Lemma 9.4. From (i) and (ii) we deduce that the monodromy representation of $``$ on $`^n/V_\lambda `$ is equivalent to the action of $`P_n`$ on $`\text{H}_1(M,_\chi ^1)/w_{\mathrm{}}`$. But it is clear from the exact sequence of the pair $`(M,^1\{z_1,\mathrm{},z_n\})`$ that we have a natural isomorphism $`\text{H}_1(M,_\chi ^1)/w_{\mathrm{}}\text{H}_1(^1\{z_1,\mathrm{},z_n\},_\chi ^1)`$. $`\mathrm{}`$ ###### Remark 9.6 Since we have seen that $`T_e(\stackrel{~}{M}_r)`$ contains an invariant line, the corresponding representation of $`P_n`$ must be on $`\text{H}_1(M,^1)`$, not on $`\text{H}^1(M,)`$ (the latter has an invariant hyperplane). ## 10 The Gassner Representation. We will follow \[Bi\] and \[Mo\] for our treatment of the Gassner representation. We begin with a quick review of the Fox calculus. Let $`G`$ be a finitely generated group and $`M`$ a $`G`$-module. Let $`[G]`$ be the group ring. ###### Definition 10.1 A derivation $`D:[G]M`$ is a $``$-linear map satisfying $$D(fh)=(D(f))ϵ(h)+fD(h)$$ where $`ϵ:[G]`$ is the augmentation. We let $`Der(G,M)`$ denote the space of derivations. ###### Remark 10.2 The restriction of each derivation $`D`$ to $`G`$ is a 1-cocycle $`\delta Z^1(G,M)`$. Conversely, given a 1-cocycle $`\delta Z^1(G,M)`$ we define a derivation $`D`$ by $$D(\underset{i=1}{\overset{n}{}}c_ig_i)=\underset{i=1}{\overset{n}{}}c_i\delta (g_i).$$ Thus $`Der(G,M)`$ and $`Z^1(G,M)`$ are canonically isomorphic. We will identify them henceforth. In the case $`G`$ is the free group $`𝔽_n`$ on the generators $`\{x_1,..,x_n\}`$ there is a unique derivation $`\frac{}{x_i}Der(𝔽_n,[F_n])`$ given by $$\frac{}{x_i}(x_j)=\delta _{ij},1i,jn.$$ Then $`Der(𝔽_n,[F_n])`$ is free over $`[F_n]`$ with the basis $`\frac{}{x_1},\mathrm{},\frac{}{x_n}`$. Note that the projection $`p:𝔽_n\text{H}_1(𝔽_n)^n`$ induces a ring-homomorphism $`p:[𝔽_n][\text{H}_1(𝔽_n)]`$ and a push-forward map on derivations $$p_{}:Der(𝔽_n,[𝔽_n])Der(𝔽_n,[\text{H}_1(𝔽_n)]).$$ We may identify $`[\text{H}_1(𝔽_n)]`$ with the $``$-algebra $``$ of Laurent polynomials in $`t_1,\mathrm{},t_n`$. The space $`Der(𝔽_n,)`$ is free over $``$ with the basis $`p_{}\frac{}{x_1},\mathrm{},p_{}\frac{}{x_n}`$. We will drop $`p_{}`$ henceforth. The main point in the construction of the Gassner representation is that there is a homomorphism $`\sigma :P_nAut(𝔽_n)`$. This homomorphism is described in terms of formulas in \[Bi, Corollary 1.8.3\]. There is an elementary description of $`\sigma `$ in terms of “pushing a loop along the braid”, see \[Mo, Page 87\]. In both cases the action of $`P_n`$ on $`𝔽_n`$ is a right action, i.e. there is $`\overline{\sigma }`$ such that $`\overline{\sigma }(p_1p_2)=\overline{\sigma }(p_2)\overline{\sigma }(p_1)`$. Therefore, the homomorphism $`\sigma `$ is actually given by $`\sigma (p):=\overline{\sigma }(p^1)`$. Next we note that we have an action of $`P_n`$ on $`Der(𝔽_n,)`$: $$gD(x)=D(\sigma (g)^1x).$$ Since $`P_n`$ acts trivially on $``$, $`gD`$ is still a derivation and the operator $`g`$ is $``$-linear. ###### Remark 10.3 In \[Bi\] and \[Mo\] the action of $`P_n`$ on $`Der(𝔽_n,)`$ is defined by $`gD(x)=D(\overline{\sigma }(g)x)`$. But $`\overline{\sigma }(g)=\sigma (g)^1`$ and hence $`gD=gD`$. The composition of two right actions is a homomorphism! We can now define the Gassner representation. ###### Definition 10.4 The Gassner representation $`\rho :P_nAut_{}(Der(𝔽_n,))`$ assigns to each $`gP_n`$ the operator $`g`$ on $`Der(𝔽_n,)`$, where $`Der(𝔽_n,)`$ is considered as a free $``$-module of rank $`n`$. It is traditional to represent $`\rho (g)`$ as an element $`(a_{ij})`$ of $`GL_n()`$ using the basis $`\frac{}{x_1},\mathrm{},\frac{}{x_n}`$, see \[Bi, Page 119\], \[Mo, Page 194\]: $$a_{ij}=\frac{}{x_j}\overline{\sigma }(g)x_i\|_{x_i=t_i}.$$ The Gassner representation is reducible. We will see shortly that $`Der(𝔽_n,)`$ contains the $`P_n`$-fixed line $`B^1(𝔽_n,)`$ and the $`P_n`$-invariant hyperplane $`Der(\mathrm{\Gamma }_n,)`$. The line does not intersect the hyperplane, nor it is complementary to it ($``$ is not a field). We begin by describing the line. We have seen that $`Der(𝔽_n,)Z^1(𝔽_n,)`$. Consequently, $`Der(𝔽_n,)`$ contains $`B^1(𝔽_n,)`$, the Eilenberg-MacLane 1-coboundaries. Since $`C^0(𝔽_n,)`$ and $`P_n`$ acts trivially on $``$, $`P_n`$ will also act trivially on $`B^1(𝔽_n,)`$. ###### Lemma 10.5 $`B^1(𝔽_n,)`$ is a free rank 1 submodule of $`Z^1(𝔽_n,)`$ with the basis $`_{i=1}^n(1t_i)\frac{}{x_i}`$. Proof: Recall that the coboundary $`\delta :C^0(𝔽_n,)C^1(𝔽_n,)`$ is given by $$\delta \mathrm{}(x_i)=\mathrm{}x_i\mathrm{}=\mathrm{}t_i\mathrm{}=(1t_i)\mathrm{}$$ But $`(1t_i)\mathrm{}=\mathrm{}\delta 1(x_i)`$, thus $`\delta `$ is $``$-linear and $`B^1(𝔽_n,)=(\delta 1)`$. We conclude by observing that $$\delta 1=\underset{i=1}{\overset{n}{}}(1t_i)\frac{}{x_i}\text{ }\text{ }\mathrm{}$$ We now describe the hyperplane. The element $`x_{\mathrm{}}=x_1\mathrm{}x_n𝔽_n`$ is fixed by $`P_n`$. We define $$Der(𝔽_n,)^{\mathrm{}}:=\{DDer(𝔽_n,):Dx_{\mathrm{}}=0\}$$ ###### Lemma 10.6 (i) $`Der(𝔽_n,)^{\mathrm{}}`$ is a free summand of $`Der(𝔽_n,)`$ of rank $`n1`$. (ii) The quotient map $`𝔽_n\mathrm{\Gamma }_n`$ induces an isomorphism $`Der(\mathrm{\Gamma }_n,)Der(𝔽_n,)^{\mathrm{}}`$ of $`P_n`$-modules. Proof: Let $`\{y_1,\mathrm{},y_n\}`$ be the basis for $`𝔽_n`$ given by $`y_i=x_1\mathrm{}x_i`$, $`1in`$. Then $`Der(𝔽_n,)`$ is free on $`\frac{}{y_1},\mathrm{},\frac{}{y_n}`$ and $`Der(𝔽_n,)^{\mathrm{}}`$ is free on $`\frac{}{y_1},\mathrm{},\frac{}{y_{n1}}`$. The statement (ii) is clear. $`\mathrm{}`$ ###### Definition 10.7 The reduced Gassner representation is the restriction of the action of $`P_n`$ from $`Der(𝔽_n,)`$ to $`Der(\mathrm{\Gamma }_n,)`$: $$\rho :P_nAut_{}(Der(\mathrm{\Gamma }_n,)).$$ We may represent $`\rho (g)`$, $`gP_n`$ as elements of $`GL_{n1}()`$ relative to the basis $`\frac{}{y_1},\mathrm{},\frac{}{y_{n1}}`$. Observe that $`B^1(𝔽_n,)`$ does not intersect $`Der(\mathrm{\Gamma }_n,)`$, indeed $$\mathrm{}\delta 1(x_{\mathrm{}})=\mathrm{}(1t_1\mathrm{}t_n)0.$$ ###### Remark 10.8 We will see below that there exist homomorphism images of $`Der(𝔽_n,)`$ such that the image of $`B^1(𝔽_n,)`$ is contained in the image of $`Der(\mathrm{\Gamma }_n,)`$. Hence $`B^1(𝔽_n,)`$ is not a complement to $`Der(\mathrm{\Gamma }_n,)`$. Note also that there is a representation of $`P_n`$ on $`\text{H}^1(𝔽_n,)=Z^1(𝔽_n,)/B^1(𝔽_n,)`$. We do not know whether or not $`\text{H}^1(𝔽_n,)`$ is a free $``$-module. We now have ###### Definition 10.9 Let $`\alpha =(\alpha _1,\mathrm{},\alpha _n)`$ with $`\alpha _j^{}`$, $`1jn`$ and $``$ be an $``$-module. Then the specialization $`_\alpha `$ of $``$ at $`\alpha `$ is defined by $`_\alpha =_{}_\alpha `$. Here $`_\alpha `$ is the complex line equipped with the $``$-module structure $`t_iz=\alpha _iz`$, $`z`$. More concretely, $`_\alpha `$ is the quotient of $``$ by the submodule of elements $`\{(t_j\alpha _j)m,1jn,m\}`$. Suppose that $`TEnd_{}()`$. Then $`T`$ induces an element $`T_\alpha =T1`$ of $`End(_\alpha )`$. Now assume that $``$ is free on $`m_1,\mathrm{},m_n`$. Then $`m_11,\mathrm{},m_n1`$ is a vector space basis for $`_\alpha `$. The matrix of $`T_\alpha `$ relative to this basis is obtained from a matrix of $`T`$ relative to $`m_1,\mathrm{},m_n`$ by substituting $`\alpha _j`$ for $`t_j`$, $`1jn`$. Now we return to the case in hand. We have $`\lambda _1,\mathrm{},\lambda _n`$ with $`\lambda _1+\mathrm{}+\lambda _n=0`$. Define $`\alpha _j:=e^{2\pi i\lambda _j}`$, $`1jn`$; whence $`\alpha _1\mathrm{}\alpha _n=1`$. ###### Lemma 10.10 Suppose that $`\alpha =(\alpha _1,\mathrm{},\alpha _n)`$ satisfies $`\alpha _1\mathrm{}\alpha _n=1`$. Then in the specialization $`Der(𝔽_n,)_\alpha `$ the image of the fixed line $`B^1(𝔽_n,)`$ is contained in the image of the invariant hyperplane $`Der(\mathrm{\Gamma }_n,)`$. Proof: $`\delta 1(x_{\mathrm{}})=1\alpha _1\mathrm{}\alpha _n=0`$. $`\mathrm{}`$ ###### Corollary 10.11 The specialization $`Der(\mathrm{\Gamma }_n,)_\alpha `$ contains a $`P_n`$-fixed line $`B^1(𝔽_n,)_\alpha `$. Now we observe that $`Z^1(𝔽_n,)_\alpha =Z^1(𝔽_n,_\chi )`$, the group of 1-cocycles with values in the 1-dimensional module defined by $`\chi (x_j)=\alpha _j,1jn`$. Moreover $$Z^1(\mathrm{\Gamma }_n,)_\alpha =Z^1(\mathrm{\Gamma }_n,_\chi ),$$ the group of $`_\chi `$-valued 1-cocycles that annihilate $`x_{\mathrm{}}`$ and $$B^1(𝔽_n,)_\alpha =B^1(\mathrm{\Gamma }_n,_\chi ).$$ We obtain ###### Proposition 10.12 Suppose $`\alpha =(\alpha _1,\mathrm{},\alpha _n)`$ satisfies $`\alpha _1\mathrm{}\alpha _n=1`$. Then the specialization of the reduced Gassner representation at $`\alpha `$ contains a $`P_n`$-invariant line. The quotient of the representation of $`P_n`$ by this line is $`\text{H}^1(\mathrm{\Gamma }_n,_\chi )`$. Theorem C follows. Michael Kapovich: Department of Mathematics, University of Utah, Salt Lake City, UT 84112, USA ; kapovich$`\mathrm{@}`$math.utah.edu John J. Millson: Department of Mathematics, University of Maryland, College Park, MD 20742, USA ; jjm$`\mathrm{@}`$math.umd.edu
warning/0002/astro-ph0002110.html	ar5iv	text	# The SIMBAD astronomical database ## 1 Introduction ### 1.1 The CDS The Centre de Données astronomiques de Strasbourg (CDS) defines, develops, and maintains services to help the astronomers find the information they need from the very rapidly increasing wealth of astronomical information, and particularly of on-line information. CDS is operated at the Strasbourg astronomical Observatory, under an agreement between French Institut National des Sciences de l’Univers (INSU) and Université Louis Pasteur, Strasbourg (ULP). CDS personnel created and implemented the Simbad data bank and maintain its data and software system. A detailed description of the CDS on-line services can be found, e.g., in Egret et al. (cds-amp2 (1995)) and in Genova et al. (cds-hub (1996), cds (1998), cds2000 (2000)), or at the CDS web site<sup>1</sup><sup>1</sup>1Internet address: http://cdsweb.u-strasbg.fr/. Questions or comments about the CDS services can be sent to the hot line question@simbad.u-strasbg.fr. ### 1.2 SIMBAD The Simbad database contains identifications, ‘basic data’, bibliographical references, and selected observational measurements for more than 2.7 million astronomical objects (November 1999). Data and information published in Simbad come from selected catalogues and tables and from the whole astronomical literature. The specificity of the Simbad database is to organize the information per astronomical object, thus offering a unique perspective on astronomical data. This is done through a careful cross-identification of objects from catalogues, lists, and journal articles. The ability to gather together any sort of published observational data related to stars or galaxies has made Simbad a key tool used worldwide for all kinds of astronomical studies. Simbad is the acronym for Set of Identifications, Measurements and Bibliography for Astronomical Data. The main access point to Simbad is the WWW home page<sup>2</sup><sup>2</sup>2http://simbad.u-strasbg.fr/Simbad; there is a mirror copy at SAO, Harvard<sup>3</sup><sup>3</sup>3http://simbad.harvard.edu/Simbad. ### 1.3 Historical background Building a reference database for stars – and, later, for extragalactic objects and all astronomical objects outside the Solar System – has been the first goal of the CDS: Simbad is the result of an on-going effort which started soon after the creation of CDS in 1972. Simbad was created by merging the Catalog of Stellar Identifications (CSI, Ochsenbein et al. csi (1981)) and the Bibliographic Star Index (Ochsenbein bsi (1982)) as they existed until 1979. The resulting data base (at that time, about 400,000 objects, mainly stars) was then expanded by the addition of source data from the catalogs and tables, and by new literature references. The database was extended to galaxies and other non-stellar objects in 1983 (Dubois et al. dubois83 (1983)). For details about the early developments of Simbad see Egret (story (1983)). The first on-line interactive version of Simbad was released in 1981, and operated at the Strasbourg Cronenbourg computer center until December 1984, when it was moved to Université Paris-Sud at Orsay, and operated there until June 30, 1990. The database is now hosted on a Unix server, at the Strasbourg Observatory. The original command line interface has been complemented by an interactive X-Window interface (XSimbad) in 1994, and by a World-Wide Web interface in 1996. There is also a client/server mode, providing quick responses to simple queries, essentially for the name resolution in archives and information systems (see Section 5). For descriptions of earlier stages of the database, see Heck & Egret (messenger (1986)), and Egret et al. (ampersand (1991)). ## 2 SIMBAD main features Simbad is, in the first place, a database of identifications, aliases and names of astronomical objects: in principle any name found in the literature – provided it is given as a syntactically correct character string – can be submitted to Simbad in order to retrieve basic information known for this object, as well as pointers to complementary data and bibliography. This implies a continuous careful cross-identification of objects from catalogues, lists, and journal articles. This ability to gather together any sort of published observational data related to stars or galaxies is the first key feature of Simbad. Another unique feature is the complete and up-to-date bibliographic survey of the astronomical literature: objects are associated with the references of all papers in which they are mentioned, independently of the aliases used to name the object. In addition, the Dictionary of Nomenclature (Section 8) is an essential tool for managing the very complex nomenclature of objects found in the literature, and for matching naming variations with those adopted or simply accepted by Simbad. It also includes hints for helping to solve ambiguities, according to the type of object, or to the format. This is complemented by the sesame module within Simbad, for the management of possible variations in the naming of astronomical objects. The database management system of Simbad (Section 4) has been developed in-house at CDS, using the concepts of object-oriented programming. Simbad is kept up-to-date (Section 7) on a daily basis, as the result of the collaboration of CDS with bibliographers in Institut d’Astrophysique de Paris and the Paris and Bordeaux observatories. The statistical contents of the database (Section 3) can be summarized in a few figures as follows (the figures quoted are statistics of November 1999): * entries for about 2.7 million astronomical objects (stars, galaxies and all astronomical objects outside the solar system); * a cross-index of 7.5 million identifiers related to 4500 astronomical catalogues and tables, lists, and observation logs of space missions; * observational data from some 25 different types of data catalogues and compilations; * a bibliographic survey covering the astronomical literature since 1950 for stars, and since 1983 for extragalactic objects: more than 3 million citations from 110,000 different papers. ## 3 SIMBAD astronomical contents ### 3.1 The objects The Simbad data base presently contains information for about: * 1,500,000 stars; * 450,000 galaxies; * 100,000 other non-stellar objects (planetary nebulae, clusters, H ii regions, etc.); * and some 650,000 additional objects observed at various wavelengths (Radio, IR, X), and for which classification is not yet assigned. The only astronomical objects specifically excluded from Simbad are the Sun and Solar System bodies. The Simbad database is primarily organized per astronomical object. The aim is to provide, as much as possible, the user with all published information (identifications, observational data, bibliographical references, and pointers towards external archives) concerning any given object, or list of objects. The two main channels for feeding the database are the following: * the daily scanning of papers published in the astronomical literature provides new references and new identifiers for existing objects, as well as opportunities to create new objects, using the basic data possibly given in the article; * the complete (or partial) folding of selected catalogues into the database serves as a basis for improving the completeness and multi-wavelength coverage of the database. Catalogues are selected for integration with priority given to those which can help provide an optimal support for the large projects conducted within the astronomical community. A large effort was, for instance, devoted in recent years to stellar catalogues (PPM, HIC, CCDM), in the context of the Hipparcos project, and to multi-wavelength identifications (IRAS PSC, Einstein 1E and 2E catalogues, older X–ray catalogues, the IUE Merged Log, etc.). The Hipparcos and Tycho catalogues (ESA tyc (1997)) have been recently included, and inclusion of the ROSAT All Sky Survey is planned in the near future. In parallel, the systematic scanning of the bibliography reflects the diversity and general trends of research in astronomy, and takes into account shorter lists. The published lists from the microlensing surveys, or e.g. the EUVE catalogues, were folded into the database as a result of this scanning. When an object is found in the literature or in a catalogue, its possible cross-identification with objects already in Simbad is systematically studied, before entering the reference and the new object name in the database. About 4500 different names of catalogues or object lists from published papers can currently be found in Simbad, covering all the wavelength domains from high energy astrophysics to radio. When no proper name is suggested by the authors, or when the acronym generates an ambiguity with already existing ones, the current practice, shared with the NED database, is to create an acronym within brackets using the initials of the last names of the first three authors, and the year of publication. For example, \[HFE83\] 366 is the 366th entry in the main table of a paper by Helmer, Fabricius, Einicke and colleagues published in 1983. From the year 2000 on, the year will be noted with four digits (e.g., \[ABC2000\]). Many objects have more than one name, since the database contains more than 7.5 million object names for 2.7 million objects. Examples of objects with more than 50 identifiers, are the galaxy M 87 in Virgo, the bright stars Procyon and Capella, the quasar 3C 273, the Crab Nebula. To help the users with the complex nomenclature of astronomical objects, the CDS now maintains and distributes on–line the Dictionary of Nomenclature of Celestial Objects (first developed by Lortet et al. dic2 (1994); see Section 8). ### 3.2 The Data In the following, the word object will be used to designate a star, non stellar object, or collection of objects such as a cluster, which corresponds to an individual entry in Simbad. For each object, the following data are included when available: * Basic data: : object type, coordinates, proper motion, $`B`$ and $`V`$ magnitudes, spectral type, parallax, radial velocity; : object type, coordinates, blue and visual integrated magnitudes, morphological type, dimension, radial velocity or redshift. : object type, position, $`B`$ and $`V`$ magnitudes. * Cross-identifications from some 4500 catalogues and tables, either completely or partially included in the data base. * Observational data (also called measurements), for some 25 data types. A list of these types is given in the Table 2. * General bibliography, including references to all published papers since 1983 citing the object under any of its designations. For stars, the bibliography starts as early as 1950, but with a smaller coverage of the literature. Simbad also includes a few hundred references before 1950, but without any systematic trend. The bibliography gives access to abstracts and electronic articles when available (either directly from publishers, or through ADS). Currently about 100 journals covering the complete astronomical literature are regularly scanned. A complete list is available on line<sup>4</sup><sup>4</sup>4http://simbad.u-strasbg.fr/guide/chH.htx. In the following, a more detailed description of some of these elements is given. #### 3.2.1 Object type The object type refers to a hierarchical classification of the objects in Simbad, derived by the CDS team on the basis of the catalogue identifiers (as proposed by Ochsenbein and Dubois type (1992)). From Star to Maser source, or Cluster of Galaxies, some 70 different categories, general, or very specific, are proposed (see examples in Table 1). A complete list is available on line<sup>5</sup><sup>5</sup>5http://simbad.u-strasbg.fr/guide/chF.htx. This classification is intended to help the user select objects out of the database (e.g. through the filter procedure, see Section 4.5). It is also a powerful tool for data cross–checking and quality control. It has been designed to be practical and useful, and complements other features also available in Simbad (morphological type or spectral type information, catalogues, and measurements). It can follow the evolution of astronomy, with the introduction of new categories recently appeared in the literature (e.g., in the last years, Low-Mass or High-Mass X-Ray binary, Microlensing event, or Void). Each class has normally a standard designation, a condensed one (used in tables) and an extended explanation. The classification uses a hierarchy with four levels, reflecting our knowledge of the characteristics of the astronomical object. For instance, an object can be classified as a “Star” (this is level 1). If photometric observations have shown variability of the object, it can be classified as a “Variable star” (this is level 2). Examples of level 3 and 4 are “Pulsating variable”, and “Cepheid”. This hierarchy of object types (and their possible synonyms) is managed in the database in such a way that selecting variable stars (V\) is understood as selecting objects classified as V\, and all subdivisions (e.g. PulsV\, Mira, Cepheid, etc.). If the user is only interested in RR Lyrae type stars, he/she will use the RRLyr type, leaving aside all other variable stars for which the variability mode is different, or not known. The classification emphasizes the physical nature of the object rather than a peculiar emission in some region of the electromagnetic spectrum or the location in peculiar clusters or external galaxies. Therefore objects are classified as peculiar emitters in a given wavelength (such as UV or IR source) only if nothing more about the nature of the object is known, i.e. it cannot be decided on the sole basis of the basic data whether the object is a star, a multiple system, a nebula, or a galaxy. For instance, if an object appears only in the IRAS catalogue, it is automatically classified as IR object: it is left to the user to decide to go further and to derive, e.g. on the basis of the IRAS colors, the probability for the source to be stellar or extragalactic. Because there is at most one object type per object, this classification should be used with caution when extracting samples out of the database. This is typically the case for the wavelength types: using IR or X as a criterion cannot generate a sample of all IRAS sources, or all X-ray emitting objects, since a number of them are in fact classified as stars, galaxies, etc. #### 3.2.2 Coordinates, Proper motion, Parallax, and Radial Velocity or Redshift The coordinates were originally stored in the database in the FK4 system for equinox and epoch 1950.0. A major change was undergone in 1999, when they were moved to the International Celestial Reference System (ICRS, see Feissel & Mignard ICRS (1998)) at epoch 2000.0, after the publication of the Hipparcos and Tycho catalogues. The position data frame has become more complex, grouping together all data needed for computing the coordinates into any reference frame, at any epoch and equinox: the coordinates themselves, the proper motion, the parallax and the radial velocity or redshift. All these data contain the same subfields: the original data, displayed with a number of digits consistent with the announced precision of the data; a quality code from ’A’ (reference data) to ’E’ (unreliable origin); an error box (either a standard error, or an ellipse), and the bibliographic reference of the data. In earlier versions of Simbad, the determination of the position for another equinox used to take only precession into account. In the current version, a change of equinox takes into account not only the precession but also the proper motion, the reference frame (FK4, FK5, ICRS), and, when they are known, the parallax and radial velocity. When no epoch is specified, the year of the equinox is used by default. Data come from various sources. When astrometric data are available, the most accurate one has been selected for the ’basic data’. Other values may be available as measurements (in the pos type). The Hipparcos and Tycho catalogues (ESA tyc (1997)) constitute the major source of positions for stars. The coordinates precision may vary from $`1\mathrm{°}`$ to $`1/10`$ mas. The default display format provides equatorial coordinates in the ICRS system at epoch 2000.0, and in the FK5 system at equinoxes 2000 and 1950, as well as galactic coordinates. Coordinates in the FK4 system, and ecliptic or super-galactic coordinates can be computed on request. The proper motions ($`\mu _\alpha \mathrm{cos}\delta ,\mu _\delta `$) are given in mas/year, together with their standard errors (in mas/year). The primary source of proper motions is the Hipparcos and Tycho catalogues (ESA tyc (1997)). The errors for positions or proper motions are expressed as error ellipses, made of three numbers, within brackets: the major axis, the minor axis, and the position angle of the major axis (measured from North to East). Major and minor axes are expressed in mas for the position, and mas/yr for the proper motion; the position angle is expressed in degrees, in the range $`[0\mathrm{°},180\mathrm{°}[`$. When available, the stellar parallax is given in mas, together with the associated error within brackets. The primary source is the Hipparcos and Tycho catalogues (ESA tyc (1997)). Radial velocity (in km/s), or redshift (for extragalactic objects) are currently available for some 160,000 objects. They are stored in their original type (either redshift, or radial velocity in km/sec), associated with the standard error. Display can be done in the original type or forced to be one of the two types, using the corresponding translation formula. Stellar radial velocity data have been compiled with the collaboration of Observatoire de Marseille. For extragalactic objects, up-to-date redshift information has recently been imported from the NASA/IPAC Extragalactic Database (NED, Helou et al. ned (2000)) as a result of the ongoing exchange agreement: the Simbad team is providing NED with bibliographic coverage of extragalactic objects for all astronomical journals, and is being given access, in return, to extragalactic data collected by NED. Tables from individual articles constitute the other major source of information. #### 3.2.3 Magnitudes $`B`$ and $`V`$ magnitudes are given, when possible, in the Johnson’s $`UBV`$ system. Both magnitudes may be followed by a semicolon meaning they cannot be made homogeneous to the $`UBV`$ system. In addition the following flags may appear: a ‘D’ flags a joint magnitude in a double or multiple system; * a ‘V’ indicates a variable magnitude and is followed by a coded index giving a rough estimate of the amplitude: \| code \| definition \| \| --- \| --- \| \| 1 \| 1/100 mag. \| \| 2 \| 1/10 mag. \| \| 3 \| 1 mag. \| \| 4 \| more than 1 mag. \| \| ? \| suspected variable \| When possible the magnitudes have been taken from the Tycho Reference Catalogue (Høg et al. TRC (1998)) where $`B`$ and $`V`$ magnitudes are derived from the original $`B_T`$ and $`V_T`$. Another major source is the $`UBV`$ compilation of Mermilliod (UBV (1987)). Otherwise the data would come from one of the published papers associated to the object. #### 3.2.4 Stellar Spectral type The spectral types of stars have been selected preferably in the Michigan Catalogues of Two-Dimensional Spectral Types for the HD stars (Houk mss (1975), and seq.), or in the bibliographical surveys of MK classifications (Jaschek MK-MJ (1978)). In the absence of a full MK classification, the HD spectral type is recorded. Most of the spectral types need less than 5 characters, but this field can be as long as 36 characters. #### 3.2.5 Morphological type and Dimension of galaxy The morphological types of galaxies have been selected primarily from the Uppsala General Catalogue of Galaxies (UGC, Nilson ugc (1973)), the Morphological Catalogue of Galaxies (MCG, Vorontsov-Velyaminov, mcg (1962)), and other catalogues (see Dubois et al. dubois83 (1983)). In complement, the following data, primarily from UGC, are given, when available, for galaxies: \| $`\mathrm{log}D_{25}`$ \| logarithm of the major axis $`a`$ expressed in tenths of arc minutes; \| \| --- \| --- \| \| $`\mathrm{log}R_{25}`$ \| logarithm of the ratio $`a/b`$ where $`a`$ and $`b`$ are the major and minor axis; \| \| orientation \| orientation angle (in degrees) \| \| (inclination) \| inclination (in units of $`15\mathrm{°}`$ from 0 to 7) \| ### 3.3 Cross–identifications #### 3.3.1 Aliases Cross–identifications of stars and galaxies have been searched for Simbad entries from (currently) about 4500 source catalogues and tables, included, either completely or partially, in the data base. The index of 7.5 million aliases, thus constituted, is one of the unique features of the Simbad database. Aliases may serve as entry points for related catalogues and tables (e.g. in VizieR). Cross-fertilization of a given research with previous studies of the same object published in the astronomical literature is made directly possible from the alias list. The index of names and aliases constitutes the basis for the Simbad name resolver which provides, in response to any object name, the set of coordinates corresponding to the object position on the celestial sphere, or the list of papers citing the object. The name resolving power of Simbad is used by many archives and information systems (such as the archives of Hubble Space Telescope or European Southern Observatory, the High Energy Astrophysics Science Archive Center, the Astrophysics Data System, servers of the Digitized Sky Surveys, etc.). There is no Simbad preferred name for objects<sup>6</sup><sup>6</sup>6In the early times of the Catalog of Stellar Identifications (Ochsenbein et al. csi (1981)), the Durchmusterung number had been used as a preferred name for stars.: all aliases can be equally used. A short list of major catalogues is used internally to put at the top of the list the most common name according to the object type (e.g., Messier or NGC identifier for galaxies and nebulae). All other identifiers are presented in alphabetical order. A command of the Simbad native node (‘selectid’), and an option in the sampling form of the WWW interface, allow the user to impose a list of identifiers to be used when displaying object lists. #### 3.3.2 Multiple systems It is to be noted that for a double system in which the components can be observed separately, Simbad frequently includes three entries: A and B components, and an additional entry for the joint system (AB), the latter entry carrying the observational data and references related to the system as a whole. This has to be taken into account in statistical studies such as stellar counts. ### 3.4 Observational data Observational data are presently given for the measurement types listed in Table 2. For each data type, one can retrieve individual data with their bibliographical references, and, when available, weighted means computed from existing observed values by specialists in the related field. When measurements are listed as a result of a Simbad query, they are normally preceded by a header providing a very short title to each listed parameter. The important rôle now played by the VizieR database of catalogues (Ochsenbein et al. vizier (2000)), coming with easier interoperability of services, is changing the strategy for inclusion of observational measurements into Simbad. Let us take the example of the Hipparcos and Tycho catalogues (ESA tyc (1997)): once the HIP or TYC identifier is available from Simbad it appears convenient enough to provide the user with a WWW link to the corresponding data in VizieR rather than overloading the Simbad database with the full Hipparcos and Tycho catalogues. This functionality is currently being implemented for important catalogues which have already been cross-identified. As a complement, the WWW interface includes pointers to external archives, currently: the INES database of the IUE project (Rodriguez-Pascual et al. ines (1999)); the high-energy observational archives at heasarc (HEASARC team heasarc (1995)). ### 3.5 Bibliographical references One of the key features of the Simbad astronomical database is the unique coverage of bibliographical references to objects. The bibliographic index contains references to stars from 1950 onwards, and to galaxies and all other objects outside the solar system from 1983 onwards. Presently (November 1999) there are about 3 million references taken from 110,000 papers published in the 100 most important astronomical periodical publications. #### 3.5.1 Bibliographical data Articles are scanned in their entirety, and references to all objects mentioned in the title, in the abstract, in the text, in the figures, or in the tables are included in the bibliography. Tables larger than 1000 objects are usually considered as catalogues and processed separately. No assessment is made of the relevance of the citation in terms of astronomical contents: the paper can be entirely devoted to the object, or simply give a side mention of it – in both cases this gives a reference in Simbad. Note that, for instance, the NED team (Helou et al. ned (2000)) applies a different strategy when covering bibliography of extragalactic objects, and tends to select only those papers that appear most relevant. Clearly, Simbad approach favours exhaustivity, at the cost of increased information noise. A code (nicknamed bibcode) is assigned to each considered paper: this 19-digit bibcode contains in principle enough information to locate the article (including year of publication, journal, volume, page, etc.). When one retrieves the bibliography of a Simbad object, a list of codes is usually given, and – according to the options used – these codes are automatically matched against a bibliographic file which provides the full reference, title and list of authors for each citation, together with an anchor pointing to the electronic version of the article. Currently, in Simbad, about 50% of the objects have no bibliographic reference. Among the most cited objects (more than 2000 references) are the Large Magellanic Cloud, M 31, 3C 273, and the supernova SN 1987A. #### 3.5.2 Bibliographic reference coding convention The structure of the 19-digit bibcode has been defined in close collaboration with the NED group at NASA/IPAC so that both databases share the same coding system (Schmitz et al. bibcode (1995)). It is also used, with some adjustments, by the Abstract Service of the Astrophysics Data System (ADS, Kurtz et al. ADS (2000)), and by the electronic journals (see e.g., Boyce & Dalterio epub (1996)). Reference codes have the following general structure: YYYYJJJJJVVVVMPPPPA Year of the publication. Standard abbreviation for the periodical. Volume number (for a journal) or, in the second character of this field, one of the following abbreviations for another publication: B (book), C (catalogue), P (preprint), R (report), S (symposium), T (thesis), U (unpublished). Specific qualifier for a paper: L letter p pink page (in MNRAS) a-z issue number within a volume A-K issue designation used by publisher Q-Z to distinguish articles on the same page. Page number. First letter of the first author’s last name (or ‘:’ if the first author cannot be identified). Example: 1991A&A...246L..24M for Astron. Astrophys. 246, L24, 1991, a Letter to the Editor of Astronomy & Astrophysics, by Motch et al. For a complete description see Schmitz et al. (bibcode (1995)), or the WWW server<sup>7</sup><sup>7</sup>7http://cdsweb.u-strasbg.fr/simbad/refcode.html. #### 3.5.3 Comments in the references Several types of comments are associated with the references in Simbad and normally displayed after the reference: * General comments: they are often comments added by the bibliographers, about the problems encountered while cross-identifying the objects mentioned in the paper, typos in object names, etc. * Notes about the existence of associated electronic tables, or abstracts in the CDS server. Papers including no object are also flagged. * Information on how the quoted objects are named in Simbad (comments related to the Dictionary of Nomenclature of Celestial Objects). ### 3.6 Statistical aspects of the Data Contents The astronomical content of Simbad results from the complex process of folding into the database a selection of important catalogues, and of surveying the complete astronomical literature. This can be illustrated by the histogram in $`V`$ magnitudes of Figure 1. The coverage is reasonably complete up to beyond magnitude 10 for stars, after the inclusion of the Tycho catalogue. Many objects in the range 12 to 26th mag. come from extensive studies of objects in selected sky areas: deep fields, external galaxies, etc. Some well-known very large catalogues are not part of Simbad: for instance the Hubble Telescope Guide Star Catalogue (GSC, Lasker et al. gsc (1990)) is not systematically included (even if GSC identifiers have been added for all Tycho stars present in Simbad). This results from a compromise aiming to save database load as well as manpower for cross-identification and quality control. Note that VizieR and Aladin give access to the full GSC catalogue (and to even larger catalogues and databases such as USNO-A, DENIS, 2MASS). Fig. 2 illustrates the increase of the data contents of the database in the years 1990 to 1999. ## 4 SIMBAD structure and query management Simbad query mechanism can be summarized by the following key features: * Database queries can be made mainly through: + identifiers (names of astronomical objects) and lists of identifiers, + sets of coordinates (retrieving one object by its position on the sky, or extracting all objects lying in a given direction), and + sampling criteria (or filters). * Data output is driven by formats. The user may define his/her own formats or modify existing ones. Output files can be saved and mailed to the user. * The user interface is adaptable to user preferences. The database management system of Simbad has been developed by the CDS, using the concepts of object-oriented programming. ### 4.1 Object-oriented concepts The command language is using the concepts of objects (or agents). Typical object classes are: astronomical object, object list, database, session, reference list, filter, format. Examples of methods are: display, describe, bye (quit). This structure is only visible for the user of the command line interface. The WWW interface is rendered quite independent of the database structure. ### 4.2 Indexing Simbad is organized for optimized access by identifier (through an index table of object names) and by position, through an index of small regions. : A B-tree file contains all identifiers allowing a fast access to any of them. For each identifier, a record contains a pointer to the astronomical object itself in the main database. : Indexing by coordinates is done in two steps: the coordinates are mapped into a set of boxes. Simbad uses the spherical-cubic projection – a technique also used, e.g., for the Cosmic Background Explorer (COBE) data: the celestial sphere is projected onto the six faces of a cube, giving six boxes at the first level. By dividing each face into four parts, one obtains a partition at level two. Further levels are obtained by further divisions of each box into four sub-boxes. Through this mechanism one obtains 6144 boxes at the level 5 with an average size of 6 square degrees and an average number of objects of 500. Box #6145 contains all objects without recorded position. In order to optimize access to objects in a coordinate box, all objects belonging to a box should be physically grouped in a common place in the database. This is done through a clustering mechanism placing objects from the same box in data blocks linked together in the database files. When a set of criteria includes some limitations in coordinates, this generates the definition of a list of boxes including the requested area: all entries from these boxes are read and checked against the whole set of criteria. When a set of criteria includes no limits in sky position, the complete database must be scanned – a long and somewhat expensive operation, which takes typically 15 minutes in the current hardware configuration. ### 4.3 Query by identifier In principle any name found in the literature – provided it is given as a syntactically correct character string – can be submitted to the database in order to retrieve information known for this object. The general syntax of an identifier is the abbreviated catalogue name (or acronym: generally one to four characters), followed by a number or a name (character string) within the catalogue. Object names such as Vega and Altair, but also Barnard’s star, Crab Nebula, Sgr A, HDFN, or HDFS are stored in the database in a specific catalog called ‘name’, while star names in constellations, such as $`\alpha `$ Lyrae, are stored in the catalogue ‘’, and variable stars (such as RR Lyrae) in the catalogue ‘var’ (also called ‘V\’). The user can generally type Vega, Altair, alf Lyrae (or alf Lyr): the sesame name resolving module (Section 8.2) is used for guessing the catalogue and making the internal conversion. There are however some difficult cases in which the name keyword remains necessary, such as in name sgr 1900+14 where sgr stands for Soft Gamma Repeater. In addition the following hints can help the user understand the best way to submit an identifier to Simbad: Simbad is not case-sensitive at this level: ALF AQL or alf Aql are, for instance, both valid. There are some exceptions to the rule, such as the cases of the star cluster RMC 136a, or the star in a multiple system VdBH 25a A, for which case-sensitivity may be necessary for solving format ambiguities. should be abbreviated as three letters: alf, bet, for $`\alpha `$ and $`\beta `$, but also mu. nu. and pi. (with a dot), for $`\mu `$, $`\nu `$ and $`\pi `$. constellation names should be abbreviated with the usual three letters: alf Boo, del Sct, FG Sge, NOVA Her 1991. The full list is available on-line<sup>8</sup><sup>8</sup>8http://simbad.u-strasbg.fr/guide/chB.htx. Identifiers of a multiple system may generate a list of the objects of the system. For instance, ADS 5423 calls for the four components, A to D, of the stellar system around Sirius. This is true only for some specific identifiers. Clusters which have no NGC or IC number are named under the generic appellation Cl followed by the cluster name and number: e.g., Cl Blanco 1 is the 1st stellar cluster named by Blanco. Stars in clusters may belong to a ‘main’ designation list, or to subsequent lists. NGC 5272 692 is star 692 in the list by Von Zeipel, considered as the main list. Subsequent lists have designations starting with Cl\. Examples: Cl\ NGC 5272 AC 968 (list by Auriere & Cordoni); Cl\* Melotte 25 VA 13 (13th star in the list by Van Altena for Melotte 25 – the Hyades cluster); Cl\* Collinder 110 DI 1101 (list by Dawson & Ianna – there is no ‘main’ list for this cluster). More details are available in the on-line description<sup>9</sup><sup>9</sup>9http://simbad.u-strasbg.fr/guide/chC.htx. If the object name seems unknown to Simbad, the user is advised to enter the coordinates of the object: the object may actually exist in the database under a different designation. Submitting the identifier, or the name of the first author of the catalogue, to the Dictionary of Nomenclature may also give useful clues. Figure 3 illustrates the response received from the database after submitting the identifier ‘M 81’. In the identifier list, the meaning of acronyms, such as \[VDD93\], is explained through a link to the on-line Dictionary of Nomenclature. The user interface provides an option for querying around objects, with a radius set by default at 10′. This is equivalent to a query by position using the object coordinates. It is also possible to generate the list of 10 or 25 next objects following a given identifier, or to submit a list of object names, stored in a file with one identifier per line. ### 4.4 Query by coordinates Query by coordinates can be used to retrieve all objects in a circular field defined by the coordinates of the center and a radius. The coordinates can be replaced by the name of an object lying at the center of the field, in which case the coordinates are found through an internal query to Simbad. The radius can have any size (default value is 10′). Queries with a radius smaller than 1–$`2\mathrm{°}`$ are answered quite instantaneously. ### 4.5 Sampling The sampling mode (also named filter) allows users to define criteria for selecting objects in Simbad. The user may extract objects which satisfy one set of coordinate criteria, several physical criteria (using a simple syntax), objects which have specified identifiers or measurements, and, finally, objects having citations within a range of years. The WWW interface provides an interactive form which presents all possible sampling options. The resulting list may be ordered according to sort criteria and, furthermore, it is possible, through the command line mode, to define precisely the output format. Note that reading the whole database for extracting a sample spread on the whole celestial sphere is possible, but quite time-consuming (as mentioned above). The user is thus encouraged to test the filter on a limited region of the sky, before applying it to the whole database. ### 4.6 Charts and sky maps After a sampling by position the user can ask for the corresponding sky plot. This feature is only available through the WWW interface and is generally optimized for a radius range of 10–60 minutes. The maps display the objects with different symbols according to object type; symbol size for stars varies with object magnitude (see Figure 4). The maps are clickable and return the object in Simbad corresponding to cursor position. The WWW interface provides also direct access to the Aladin interactive digitized atlas (Bonnarel et al. aladin (2000)) as illustrated in Figure 5. ### 4.7 Batch mode Simbad can be queried in batch mode, by submitting a mail to the special address smbmail@simbad.u-strasbg.fr. This is especially useful in case of poor interactive connectivity, or for submitting time-consuming queries or lists. A WWW form<sup>10</sup><sup>10</sup>10http://cdsweb.u-strasbg.fr/simbad/batch.html helps to prepare the submission. ### 4.8 Resolving a bibliographical reference code It is possible to obtain a complete bibliographical reference, by entering the corresponding reference code (bibcode). A reference code can be supplied without indicating all the fields: the first reference corresponding to the truncated code will be displayed. An ampersand (&) should be added at the end of the truncated bibcode. ### 4.9 Additional tools Additional tools include special commands for querying auxiliary databases, on-line help, log files, etc. More details can be found in the Simbad User’s Guide or on the Web pages. ## 5 User Interfaces to SIMBAD There are several user interfaces to Simbad. New users are advised to go directly to the WWW interface, unless they have very specific needs. * The Web interface is currently the easiest access mode to Simbad. This interface takes benefit of the WWW features to provide the user with additional links to internal documentation, associated services (Aladin, VizieR), and external archives (currently: the INES database of the IUE project; the high-energy observational archives at heasarc). Some features, such as the finding chart, have been specifically designed for the Web and are not available through the other modes. * The command line interface is the basic underlying interface to the database, which serves as a basis for the other more user-friendly interfaces. During many years it was the sole access mode to the database, and many users who are accustomed to the commands may find it quicker and more versatile. It implies a remote login (through telnet) on the simbad machine in Strasbourg, and the user needs to have a user name and password (see Section 6.1). * A graphical interactive user interface to Simbad, XSimbad, taking benefit of the X Window environment has been developed in 1993-94 for distribution to users working in a Unix environment. It is now obsolete, because all the functionalities, and additional ones, are more easily available through the Web. * A client/server package is distributed on request to data managers of archives and information systems, when they need to organize the most efficient access to Simbad for the resolution of object names into the corresponding position, or the retrieval of other information provided by the database, such as the reference list for a given object. Written currently in C language, it can easily be plugged into any application able to access C routines. Distribution is subject to CDS approval. ## 6 SIMBAD usage ### 6.1 Charging policy Simbad is a charged service. The telnet access is protected by userid/password, and the WWW access is protected either by IP address or by password. Users have to register, and get a userid/password from the CDS staff (or from the U.S. agent for American users). In the U.S. there is no invoicing for the end-user because the charges are covered globally by NASA for all U.S. users. In Europe, the same situation is also true for users from ESO and ESA member states, thanks to an agreement signed with European Southern Observatory and ESA Space Telescope European Coordinating Facility (since January 1995). Special arrangements also exist or are currently being negotiated with other countries, including Canada, Australia, and Japan. ### 6.2 Usage statistics There are currently (November 1999) some 7000 user accounts from 64 different countries. The development of the WWW access makes difficult to keep precise track of the individual usage statistics, but the global statistics show that the world wide interest for accessing the Simbad database continues to increase regularly over the years. The number of Simbad queries evolved from about 30,000 per month in 1997 to about 100,000 per month in 1999. About 50% of the queries come through the client/server mode. A mirror copy of Simbad has been established at CfA (Harvard) for convenience of US users, and about one third of the queries are currently processed on the mirror site (including name resolving activities for the ADS and major US NASA archives). The mirror copy is managed by CDS and updated every night. ## 7 SIMBAD updates and quality control ### 7.1 Updating SIMBAD Simbad is kept up-to-date on a daily basis, as the result of the collaboration of the CDS team, in Strasbourg, with bibliographers in Observatoire de Paris (DASGAL), Institut d’Astrophysique de Paris, and Observatoire de Bordeaux (Laloë et al. maj93 (1993); Laloë maj95 (1995)) who systematically scan the articles published in some 100 astronomy journals. The references are updated very soon after reception of the journal issues, and in some cases directly from the journal table of contents, through agreements with the Editors. New data concerning the objects (identifiers, basic data), and new acronyms for catalogs or tables are being entered when appropriate. The inclusion of a large catalogue in the database is often a long-term task which may span over several months or years; the collaboration of specialists in the different fields is systematically sought. The improvement of the Simbad astronomical contents relies on a network of collaborations: a list of the main current contributors is given in Table 3. More generally, help of other contributing institutes and authors, too numerous to be cited here, is gratefully acknowledged. ### 7.2 Quality control The data contained in Simbad are also permanently updated, as a result of errata, remarks from the bibliographers (during the scanning of the literature), integration of lists and catalogues, quality controls, or special efforts initiated by the CDS team to better cover some specific domains (e.g., multi-wavelength emitters and complex objects). Requests for corrections, errata, or suggestions are regularly received from Simbad users through a dedicated hot line, at e-mail address question@simbad.u-strasbg.fr. A few dozens of messages are usually received every week, and processed on a daily basis by the member of the team who is on duty for that week, or transmitted to the key person in case of specialized questions. Remarks received from the users by this way are especially welcome, as they help the CDS team to improve the database contents through the scrutiny of specialists’ eyes. Developing new tools for quality control of the database is a major challenge for the future, and CDS is exploring possible solutions. Multivariate analysis applied to bibliographic information retrieval has been proven a possible tool for developing quality control in a database such as Simbad (Lesteven lesteven (1995)). ### 7.3 Towards automation of updating procedures The advent of electronic publishing brings new perspectives for improvement and automation of the updating procedures. In a first place, tables of contents of the major journals are now received electronically through the network, thanks to journal Editors and Publishers, thus reducing the risk of errors. Regularly, a number of electronic lists of objects are also folded into Simbad through semi-automatic procedures. The next step will be the automatic flagging of object names in the text of the articles: this has now become a very interesting medium-term goal. Two ways of achieving this flagging are currently being considered: * the first one is to ask the authors, with the help of the Editors of electronic journals, to flag astronomical object names in their text; this can be done, for instance, by the use of a `\object{ }` command within the or source, which will be eventually used to build an anchor pointing towards Simbad, or another database, in the on-line version made available on the network. This approach has been adopted by the Editors of Astronomy & Astrophysics. * another approach is the use of intelligent search tools for identifying object names within the electronic version of the paper, using a set of syntactic and semantic rules, and the Dictionary of Nomenclature as a reference database for already known objects. The first approach seems safer, provided the authors understand what exactly they are being required, and accept this (minor) additional work load. The latter implies a lot of fine tuning from the system developers. The current experience with the handling of publications (Lesteven et al. lesteven2 (1998)) suggests that both approaches may be needed, and that a careful quality control, including final check by an expert, will probably remain necessary to avoid errors or misinterpretations, and to ensure appropriate completeness. ## 8 Nomenclature ### 8.1 The Dictionary of Nomenclature Designations of astronomical objects are often confusing. A complete list of astronomical designations has been collected and published by Lortet et al. (dic2 (1994)) in the Dictionary of Nomenclature of Celestial Objects outside the Solar System. This information is available on-line through the info command, or on the WWW<sup>11</sup><sup>11</sup>11http://vizier.u-strasbg.fr/cgi-bin/Dic-Simbad. This service is the electronic look-up version of the Dictionary which is now under the responsibility of CDS. It is kept up-to-date on a weekly basis; about 15 new acronyms are incorporated every week. The Dictionary currently provides full references and usages about some 5000 different acronyms. It is used by the International Astronomical Union as a reference for its recommendations related to nomenclature. ### 8.2 The sesame module The sesame module is used inside Simbad for the management of possible variations in the naming of astronomical objects. It is based on a list of rules, written as regular expressions, allowing translation of the submitted name into its Simbad canonical form; it is only made visible to the user when a message mentions the submitted syntax and its translation. There are cases where ambiguities cannot be solved. This is actually specific to the broad context of Simbad. Let us give an example: in the context of extragalactic objects ‘N’ is a possible abbreviation for ‘NGC’ (accepted by NED); but people studying Novae would frequently use ‘N’ as an abbreviation for Nova, people studying H ii regions would use it for naming nebulae in the Magellanic Clouds (LHA 120-N or LHA 115-N), and ‘N’ has also been found in the literature for cluster stars studied by Nordlund in NGC 2099 (Cl\* NGC 2099 N), for stars studied by Neckel (\[N78\]), or even for ‘New’ parts of the galaxy NGC 1275 (\[NJS93\] in Simbad). When a name like ‘N 1992’ is submitted to Simbad the ambiguity cannot be solved without requesting additional information from the user. ## 9 Integration of SIMBAD into the CDS Hub While the CDS databases have followed different development paths, the need to build a transparent access to the whole set of CDS services has become more and more obvious with the easy navigation permitted by hypertext tools (Genova et al. cds2000 (2000)). Aladin has become the prototype of such a development, by giving comprehensive simultaneous access to Simbad, the VizieR Catalogue service, and to external databases such as NED, using a client/server approach and, when possible, standardized query syntax and formats. In order to be able to go further, the CDS has built a general data exchange model, taking into account all types of information available at the Data Center, known under the acronym of GLU for Générateur de Liens Uniformes – Uniform Link Generator (Fernique et al. glu (1998)). In the current stage of development, the WWW interface to Simbad provides access to Aladin previewer (reference image around one object), and to the Aladin interactive Java program (see Bonnarel et al. aladin (2000)). There are also links between Simbad and the bibliographic services developed or mirrored at CDS, and more generally to the ADS and the electronic journals. While this article is written stronger links between Simbad and VizieR are just being created allowing even easier transfers of data and information between both services. This will also make easier to build new links pointing to distributed data archives, beyond those already existing (currently: IUE/INES and HEASARC). ## 10 Future developments In the near future, the CDS team expects to go on enriching the database contents and system functionality. The users play an important role in that respect, by giving feedback on the desired features, on the user-friendliness of the interfaces, etc. In the context of interoperability of distributed services, as currently discussed within the ISAIA project (Interoperable Systems for Archival Information Access; Hanisch isaia (2000)), Simbad is prepared to deliver resource profiles and to format the query outputs in a standard way, for instance XML (Ochsenbein et al. adass9 (2000)). As larger and larger astronomical datasets are being produced, the CDS is studying the concepts of a new generation database of several billion objects, instead of the current several million objects. We expect Simbad to remain an essential tool for astronomical research in the years to come. ###### Acknowledgements. CDS acknowledges the support of INSU-CNRS, the Centre National d’Etudes Spatiales (CNES), and Université Louis Pasteur (ULP, Strasbourg). Many of the current developments of Simbad have been made possible by long-term support from NASA, ESA, and ESO. We thank more specifically J. Mead and G. Riegler (NASA), P. Benvenuti (ESA/ST-ECF), and P. Quinn (ESO) for their help in setting up the current agreements. Developing and maintaining the database is a collective undertaking to which many contributors – too numerous to be listed here – are associated. A special mention shall be made of M.-J. Wagner, F. Woelfel, J. Marcout (Strasbourg), A. Beyneix, G. Chassagnard (IAP, Paris), N. Ralite, S. Pasquier (Bordeaux), E. Davoust (Toulouse), and B. Skiff (Lowell Observatory), who are watching with great care over the Simbad contents. We want finally to thank Jean Delhaye, Jean Jung, Carlos Jaschek and Michel Crézé for their leadership and their vision in the consecutive early phases of the Simbad project.
warning/0002/cond-mat0002352.html	ar5iv	text	# d-wave Superconductivity in the Hubbard Model ## Introduction The discovery of high-$`T_c`$ superconductors has stimulated strong experimental and theoretical interest in the field of strongly correlated electron systems. After a decade of intensive studies we are still far from a complete understanding of the rich physics observed in high-$`T_c`$ cuprates . Angle resolved photoemission experiments on doped materials show a $`d`$-wave anisotropy of the pseudogap in the superconducting state . In underdoped materials even in the normal state this pseudogap persists , which is believed to cause the unusual non-Fermi-liquid behavior in the normal state. This emphasizes the importance of achieving a better understanding of the superconducting phase, i.e. the physical origin of the pairing mechanism, the nature of the pairing state and the character of low energy excitations. On a phenomenological basis the $`d`$-wave normal state pseudogap as well as the transition to a superconducting state with a $`d`$-wave order parameter has been described within theories where short-ranged antiferromagnetic spin fluctuations mediate pairing in the cuprates . On a microscopic level it is believed that the Hubbard model or closely related models like the t-J model should capture the essential physics of the high-$`T_c`$ cuprates . However, despite years of intensive studies, these models remain unsolved except in one or infinite dimensions. Finite size QMC calculations for the doped 2D Hubbard model in the intermediate coupling regime with Coulomb repulsion $`U`$ less than or equal to the bandwidth $`W`$, support the idea of a spin fluctuation driven interaction mediating $`d`$-wave superconductivity . But the fermion sign problem limits these calculations to temperatures too high to observe a possible Kosterlitz-Thouless transition for the 2D system . Another problem encountered in QMC calculations is their finite size character, which makes statements for the thermodynamic limit dependent on a scaling ansatz. These limitations do not apply to approximate many particle methods like the Fluctuation Exchange Approximation (FLEX) . Results of FLEX calculations for the Hubbard model are in agreement with QMC results, i.e. they show evidence for a superconducting state with $`d`$-wave order parameter at moderate doping for sufficiently low temperatures. But the FLEX method as an approximation based on a perturbative expansion in $`U`$ breaks down in the strong coupling regime $`U>W`$, where $`W`$ is the bare bandwidth. On the other hand it is believed that a proper description of the high-$`T_c`$ cuprates in terms of the one-band Hubbard model requires $`U>W`$. Calculations within the Dynamical Mean Field Approximation (DMFA) can be performed in the strong coupling regime and take place in the thermodynamic limit. But the lack of non-local correlations inhibits a transition to a state with a non-local ($`d`$-wave) order parameter. The recently developed Dynamical Cluster Approximation (DCA) is a fully causal approach which systematically incorporates non-local corrections to the DMFA by mapping the lattice problem onto an embedded periodic cluster of size $`N_c`$. For $`N_c=1`$ the DCA is equivalent to the DMFA and by increasing the cluster size $`N_c`$ the dynamic correlation length can be gradually increased while the DCA solution remains in the thermodynamic limit. Using a Nambu-Gorkov representation of the DCA we observe a transition to a superconducting phase in doped systems at sufficiently low temperatures. This occurs in the intermediate to strong coupling regime $`U>W`$ and the corresponding order parameter has $`d`$-wave symmetry. ## Method A detailed discussion of the DCA formalism was given in previous publications where it was shown to systematically restore momentum conservation at internal diagrammatic vertices which is relinquished by the DMFA. However, the DCA also has a simple physical interpretation based on the observation that the self energy is only weakly momentum dependent for systems where the dynamical intersite correlations have only short spatial range. The corresponding self-energy is a functional of the interaction $`U`$ and the Green function propagators. The latter may be calculated on a coarse grid of $`N_c=L^D`$ selected $`𝐊`$-points only, where $`L`$ is the linear dimension of the cluster of $`𝐊`$-points. According to Nyquist’s sampling theorem , this sampling of the reciprocal space at intervals of $`\mathrm{\Delta }𝐤=2\pi /L`$ implies that the DCA incorporates nonlocal dynamical correlations with a spatial range $`L/2`$ and cuts off longer ranged dynamical correlations. Knowledge of the momentum dependence on a finer grid may be discarded to reduce the complexity of the problem. To this end the first Brillouin zone is divided into $`N_c`$ cells of size $`(2\pi /L)^D`$ around the cluster momenta $`𝐊`$ (see Fig.1). The Green functions used to form the self-energy $`\mathrm{\Sigma }(𝐊,\omega )`$ are coarse grained, or averaged over the momenta $`𝐊+\stackrel{~}{𝐤}`$ surrounding the cluster momentum points $`𝐊`$ (cf. Fig.1). Thus, the coarse grained Green function is $$\widehat{\overline{G}}(𝐊,\omega )=\frac{N_c}{N}\underset{\stackrel{~}{𝐤}}{}\widehat{G}(𝐊+\stackrel{~}{𝐤},\omega )\text{,}$$ (1) where the sum runs over all vectors $`𝐤=𝐊+\stackrel{~}{𝐤}`$ within a cell around the cluster momentum $`𝐊`$. Note that the choice of the coarse grained Green function has two well defined limits: For $`N_c=1`$ the sum over $`\stackrel{~}{𝐤}`$ runs over the entire Brillouin zone, $`\widehat{\overline{G}}`$ is the local Green function, thus the DMFA algorithm is recovered. For $`N_c=\mathrm{}`$ the $`\stackrel{~}{𝐤}`$-summation vanishes and the DCA becomes equivalent to the exact solution. The dressed lattice Green function takes the form $$\widehat{G}(𝐤,\omega )=\left(\omega \mathrm{𝟙}ϵ_𝐤\tau _3\widehat{\mathrm{\Sigma }}(𝐊,\omega )\right)^1\text{,}$$ (2) with the self-energy $`\widehat{\mathrm{\Sigma }}(𝐤,\omega )`$ approximated by the cluster self-energy $`\widehat{\mathrm{\Sigma }}(𝐊,\omega )`$. To allow for a possible transition to the superconducting state we utilized the Nambu-Gorkov matrix representation in (2) where the self-energy matrix $`\widehat{\mathrm{\Sigma }}`$ is most generally written as an expansion $`\widehat{\mathrm{\Sigma }}=_i\mathrm{\Sigma }_i\tau _i`$ in terms of the Pauli matrices $`\tau _i`$. The diagonal components of $`\widehat{\mathrm{\Sigma }}`$ represent quasiparticle renormalizations, whereas the offdiagonal parts are nonzero in the superconducting state only. Since the self-energy $`\widehat{\mathrm{\Sigma }}(𝐊,\omega )`$ does not depend on the integration variable $`\stackrel{~}{𝐤}`$, we can write $$\widehat{\overline{G}}(𝐊,\omega )=(\omega \mathrm{𝟙}\overline{ϵ}_𝐊\tau _3\widehat{\mathrm{\Sigma }}(𝐊,\omega )\widehat{\mathrm{\Gamma }}(𝐊,\omega ))^1\text{,}$$ (3) where $`\overline{ϵ}_𝐊=N_c/N_{\stackrel{~}{𝐤}}ϵ_{𝐊+\stackrel{~}{𝐤}}`$. This has the form of the Green function of a cluster model with periodic boundary conditions coupled to a dynamic host described by $`\widehat{\mathrm{\Gamma }}(𝐊,\omega )`$. Here we employ the NCA to calculate the cluster Green function and self-energy respectively. A detailed discussion of the NCA-algorithm applied to the cluster model for the paramagnetic state was given in . The NCA for the superconducting state has to be extended in order to account for the hybridization to the anomalous host, which couples cluster states with different particle numbers. The self-consistent iteration is initialized by calculating the coarse grained average $`\widehat{\overline{G}}(𝐊)`$ (Eq. 1) and with Eq. 3 the host function $`\widehat{\mathrm{\Gamma }}(𝐊)`$, which is used as input for the NCA. The NCA result for the cluster self-energy $`\widehat{\mathrm{\Sigma }}(𝐊)`$ is then used to calculate a new estimate for the coarse grained average $`\widehat{\overline{G}}(𝐊)`$ (Eq. 1). The procedure continues until the self-energy converges to the desired accuracy. ## Results We investigate the single particle properties of the doped 2D Hubbard Model $$H=\underset{ij,\sigma }{}t_{ij}c_{i\sigma }^{}c_{j\sigma }+U\underset{i}{}n_in_i\text{,}$$ (4) where $`c_i^{}`$ ($`c_i`$) creates (destroys) an electron at site $`i`$ with spin $`\sigma `$ and $`U`$ is the on-site Coulomb repulsion. For the Fourier transform of the hopping integral $`t_{ij}`$ we use $$ϵ_𝐤=ϵ_o\mu 2t(\mathrm{cos}k_x+\mathrm{cos}k_y)4t^{}\mathrm{cos}k_x\mathrm{cos}k_y\text{,}$$ (5) accounting for both, nearest neighbor hopping $`t`$ and next nearest neighbor hopping $`t^{}`$. We set $`t=0.25\mathrm{eV}`$ and $`U=3\mathrm{e}\mathrm{V}`$, well above the bandwidth $`W=8t=2\mathrm{e}\mathrm{V}`$. For this choice of parameters the system is a Mott-Hubbard insulator at half filling as required for a proper description of the high-$`T_c`$ cuprates. To allow for symmetry breaking we start the iteration procedure with finite offdiagonal parts of the self-energy matrix $`\widehat{\mathrm{\Sigma }}`$. As we mentioned above one expects the order parameter of a possible superconducting phase to have $`d`$-wave symmetry. Therefore we work with a 2x2-cluster ($`N_c=4`$), the smallest cluster size incorporating nearest neighbor correlations. For the set of cluster points we choose $`𝐊_{\alpha l}=l\pi `$, where $`l=0,1`$ and $`\alpha =x`$ or $`y`$. Fig.1 illustrates this choice of $`𝐊`$-points along with a sketch of the $`d`$-wave order parameter and the coarse graining cells. Obviously, for symmetry reasons, in the case of $`d`$-wave superconductivity, we expect the coarse grained anomalous Green function to vanish at the zone center and the point $`(\pi ,\pi )`$. Whereas the anomalous parts at the points $`(0,\pi )`$ and $`(\pi ,0)`$ should be finite with opposite signs. Fig.2 shows a typical result for the local density of states (DOS) in the superconducting state along with the anomalous coarse grained Green function $`\overline{G}_{12}(𝐊,\omega )=N_c/N_{\stackrel{~}{𝐤}}c_{𝐊+\stackrel{~}{𝐤}};c_{(𝐊+\stackrel{~}{𝐤})}_\omega `$ at the cluster $`𝐊`$-points for $`t^{}=0`$, temperature $`T=137\mathrm{K}`$ and doping $`\delta =0.19`$. The anomalous coarse grained Green function vanishes at the cluster points $`(0,0)`$ and $`(\pi ,\pi )`$ but is finite at the points $`(\pi ,0)`$ and $`(0,\pi )`$, consistent with a $`d`$-wave order parameter. Note that this result is independent of the initialization of the self-energy, i.e. an additional initial $`s`$-wave contribution vanishes in the course of the iteration. Thus a possible $`s`$-wave contribution to the order parameter can be ruled out. The finite pair amplitude is also reflected in the local density of states (DOS) depicted in Fig.2a, where we show the lower sub-band of the full spectrum near the Fermi energy. It displays a pseudogap at zero frequency as expected for a $`d`$-wave order parameter. Fig.3 shows the DOS near the Fermi energy for the same parameters as in Fig.2, fixed temperature $`T=137\mathrm{K}`$, but for various dopings. Obviously, the pseudogap size, measured as the peak to peak distance, as well as the density of states at the Fermi energy do not depend strongly upon doping. However the drop in the density of states from the gap edge to the $`\omega =0`$ value first increases, reaches a maximum at about 19% doping, then decreases again. This behavior originates in the doping dependence of the anomalous Green function. In the inset we plot the coarse grained anomalous equal time Green function $`\overline{G}_{12}(𝐊,\tau =0)=N_c/N_{\stackrel{~}{𝐤}}c_{𝐊+\stackrel{~}{𝐤}}c_{(𝐊+\stackrel{~}{𝐤})}`$ for $`𝐊=(\pi ,0)`$. This number as a measure of the superconducting gap shows exactly the same behavior as the pseudogap in the density of states. The anomalous components $`\overline{G}_{12}(𝐊,\omega )`$ and hence the pseudogap in the DOS become smaller with increasing temperature and eventually vanish at a critical temperature $`T_c`$ depending on the set of parameters. The phase diagram is shown in Fig.4. As a function of doping, $`T_c(\delta )`$ has a maximum $`T_c^{max}150\mathrm{K}`$ at $`\delta 19\%`$ and strongly decreases with decreasing or increasing $`\delta `$. The qualitative behavior of $`T_c(\delta )`$ in the calculated $`T\delta `$ region agrees well with the generic phase diagram of the high-$`T_c`$ cuprates. Unfortunately, due to the break-down of the NCA at very low temperatures we are not able to extend the phase diagram beyond the region shown in Fig.4. This means in particular that we cannot predict reliable values for $`\delta _c(T=0)`$, beyond which superconductivity vanishes. The inset of Fig.4 shows the transition temperature dependence $`T_c(t^{},\delta =\mathrm{const}.)`$ on the next nearest neighbor hopping amplitude $`t^{}`$ for fixed doping $`\delta =0.18`$. As compared to $`t^{}=0`$ $`T_c`$ strongly decreases with growing negative $`t^{}`$ but increases for $`t^{}>0`$. The shape of the phase diagram as well as the $`t^{}`$-dependence of $`T_c`$ can be qualitatively understood in terms of the phenomenological picture, where spin fluctuations mediate the electron-electron interaction, which then is strong at the antiferromagnetic wave vector $`𝐐=(\pi ,\pi )`$. In Fig.5 we display the coarse grained spectra $`\frac{1}{\pi }\mathrm{}m\overline{G}_{11}(𝐊,\omega )`$ at $`𝐊=(\pi ,0)`$ in the normal state ($`T=290\mathrm{K}`$) for the next nearest neighbor hopping $`t^{}=0`$ and $`t^{}=0.05eV`$ (left and right hand side) for different dopings $`\delta `$. At the bottom we show the Fermi surfaces of the corresponding noninteracting systems ($`U=0`$) in the first quadrant of the BZ. The diagonal thick solid line indicates the set of $`𝐤`$-points which fulfill the nesting condition, i.e. which can be connected by $`𝐐`$ to equivalent $`𝐤`$-points in the opposite quadrant of the BZ. The doping dependence of the $`t^{}=0`$ and $`t^{}=0.05`$ spectra is qualitatively different. In the $`t^{}=0.05`$ case the parts of the Fermi surface near $`𝐊=(\pi ,0)`$ and $`(0,\pi )`$ fulfill the nesting condition roughly for the whole doping range, the quasiparticles couple strongly to the spin fluctuations and hence the corresponding spectra display a pseudogap at zero frequency over the entire doping range. The $`t^{}=0`$ spectra in contrast exhibit the pseudogap in the underdoped regime only ($`\delta =0.05`$), where the spin fluctuations are strong, but show a quasiparticle peak at optimal doping $`\delta =0.19`$, where the points on the Fermi surface near $`𝐊=(\pi ,0)`$ and $`(0,\pi )`$ are far from being nested. The suppression of the density of states at the Fermi energy results in a suppression of superconductivity and hence the transition temperatures drop with decreasing doping as well as decreasing $`t^{}<0`$. For positive $`t^{}`$ we obtain similar spectra and Fermi surfaces as for $`t^{}=0`$, but with a slightly enhanced density of states at the Fermi energy, resulting in higher transition temperatures. ## Summary We have used the recently developed DCA to study the long open question of whether the 2D Hubbard model shows instabilities towards a superconducting state in the intermediate to strong coupling regime. We find conclusive evidence that at moderate doping a transition to a state with offdiagonal long range order occurs and that the corresponding order parameter has pure $`d`$-wave symmetry. The corresponding temperature-doping phase diagram agrees qualitatively with the generic high-$`T_c`$ phase diagram. ## Acknowledgements It is a pleasure to acknowledge useful discussions with P.G.J. van Dongen, M. Hettler and H.R. Krishnamurthy. This work was supported by NSF grants DMR-9704021, DMR-9357199 and the Graduiertenkolleg “Komplexität in Festkörpern”. Computer support was provided by the Ohio Supercomputer Center and the Leibnitz-Rechenzentrum, Munich.
warning/0002/hep-ph0002170.html	ar5iv	text	# Bound on the tau neutrino magnetic moment from the TRISTAN experiments ## Abstract We set limits on the magnetic moment and charge radius of the tau neutrino by examining an extra contribution to the electroweak process $`e^+e^{}\nu \overline{\nu }\gamma `$ using VENUS, TOPAZ and AMY results. We find that $`\kappa (\nu _\tau )<9.1\times 10^6`$ (i.e. $`\mu (\nu _\tau )<9.1\times 10^6\mu _B,\mu _B=e/2m_e`$) and $`r^2<3.1\times 10^{31}`$cm<sup>2</sup> with Poisson statistics by combining their results. Whereas, similar to this method, with the Unified Approach we find that $`\kappa (\nu _\tau )<8.0\times 10^6`$ and $`r^2<2.7\times 10^{31}`$cm<sup>2</sup>. , and PACS number: 14.60.St The electromagnetic properties of neutrinos have been vigorously examined in recent years, since they are related to the neutrino mass and to solar neutrino problems. Neutrino oscillation between $`\nu _\mu `$ and $`\nu _\tau `$, which means a finite neutrino mass, has become realistic since evidential results by Super-Kamiokande . Gninenko has shown the bound on the tau neutrino magnetic moment from the S-K atmospheric neutrino data to be $`1.3\times 10^7\mu _B`$. Chua and Hwang recently remarked that the third-generation neutrino magnetic moment induced by leptoquarks might be of the order of $`10^{10}10^{13}\mu _B`$. If neutrinos have mass, the tau neutrino would be the most massive among $`\nu _e`$, $`\nu _\mu `$ and $`\nu _\tau `$. Therefore, the tau neutrino magnetic moment might show a relatively large value because the neutrino magnetic moment is estimated to be proportional to its mass according to the standard model extended to have the right-handed neutrino singlet ($`\nu _R`$), $`\mu _\nu =3eG_Fm_\nu /(8\pi ^2\sqrt{2})=(3.20\times 10^{19})m_\nu \mu _B`$ . Assuming $`m_{\nu _\tau }<18.2`$ MeV , the upper limit on the magnetic moment is less than $`5.8\times 10^{12}\mu _B`$, where $`e,G_F`$, $`m_\nu `$ and $`\mu _B`$ are the electron charge magnitude, the Fermi coupling constant, the neutrino mass in eV and the Bohr magneton ($`=e/2m_e`$), respectively. Thereupon, we estimated the tau neutrino magnetic moment and charge radius limits with the $`e^+e^{}\nu \overline{\nu }\gamma `$ results from three TRISTAN experimental groups. We classify the methods to evaluate the tau neutrino magnetic moment as follows: (A)cosmological estimation, (B)fixed target experiments ($`\nu _\tau e^{}\nu _\tau e^{}`$) and (C)$`e^+e^{}`$ colliding beam experiments. Method(A) gives a very strong upper limit of the magnetic moment such as $`\mu _{\nu _\tau }<6.2\times 10^{11}\mu _B`$ , $`2\times 10^{12}\mu _B`$ and $`6\times 10^{14}\mu _B`$ . However, these values are based on many cosmological assumptions. Incidentally, Grifols and Massó have argued that primordial nucleosynthesis also constrains the neutrinos charge radii to satisfy $`r^2<7\times 10^{33}\mathrm{cm}^2`$. Their argument, however, also has an implicit dependence on the neutrino mass which may allow them to be evaded. Method(B) gives $`\mu _{\nu _\tau }<5.4\times 10^7\mu _B`$ . It assumes the form factor ratio of $`f_{D_s}/f_\pi =2`$ and $`D_s`$, $`\overline{D_s}`$ production cross section $`=2.6\mu \mathrm{b}`$ to calculate $`\nu _\tau `$ flux, because $`\nu _\tau `$ beam flux has to be produced and estimated by $`D_s`$, $`\overline{D_s}`$ production. Method(C) is the most direct. Grotch and Robinett combined the results from ASP, MAC and CELLO experiments well below the $`Z^0`$ resonance and set the limits at the 90% confidence level on the magnetic moment and the charge radius of the tau neutrino, $`\mu _{\nu _\tau }<4\times 10^6\mu _B`$ and $`r^2<2\times 10^{31}\mathrm{cm}^2`$, respectively. The other is from the experiments at the $`Z^0`$ resonance. They give $`\mu _{\nu _\tau }<4.4\times 10^6\mu _B`$ and $`3.3\times 10^6\mu _B`$ . At energies well below $`Z^0`$, the dominant contribution to the process $`e^+e^{}\nu \overline{\nu }\gamma `$ involves the exchange of a virtual photon . The dependence on the magnetic moment comes from its direct coupling to the virtual photon, and the observed photon is the result of the initial-state Bremsstrahlung. While the results of the TRISTAN experiments (VENUS, TOPAZ and AMY collaborations) have been used to set limits on supersymmetric particles, we will make use of them here to set limits on the tau neutrino magnetic moment and charge radius. The standard expression for the cross section for the process $`e^+e^{}\nu \overline{\nu }\gamma `$ due to $`Z^0`$ and $`W`$ exchange (Fig.1(a)) is $`{\displaystyle \frac{d\sigma }{dxdy}}`$ $`=`$ $`{\displaystyle \frac{G_{F}^{}{}_{}{}^{2}\alpha }{6\pi ^2}}`$ (1) $`\left\{{\displaystyle \frac{M_Z^4\left\{N_\nu (g_V^2+g_A^2)+2(g_V+g_A)[1s(1x)/M_Z^2]\right\}}{[s(1x)M_Z^2]^2+(M_Z\mathrm{\Gamma }_Z)^2}}+2\right\}`$ $`{\displaystyle \frac{s}{x(1y^2)}}\left[(1x)(1x/2)^2+x^2(1x){\displaystyle \frac{y^2}{4}}\right],`$ where $`x=E_\gamma /E=2E_\gamma /\sqrt{s}`$ is the photon energy in units of the incident beam energy, $`y=\mathrm{cos}\theta _\gamma `$ is the direction cosine of the photon momentum with respect to the incident beam direction, $`\alpha `$ is the fine-structure constant, $`s`$ is the square of the center of mass energy, $`N_\nu `$ is the number of low-mass neutrino generations, $`M_Z`$ is the mass of the $`Z^0`$, $`\mathrm{\Gamma }_Z`$ is the total width of $`Z^0`$, $`g_V=1/2+2\mathrm{sin}^2\theta _W`$($`\theta _W`$ is the weak mixing angle) and $`g_A=1/2`$. It is worth noting that the $`(g_V^2+g_A^2)`$ term of equation (1) arises from the square of the s-channel $`Z^0`$ amplitude, the ‘2’ term from the square of the t-channel $`W`$-exchange amplitude, and the $`(g_V+g_A)`$ term from $`Z^0W`$ interference. We now allow for a neutrino electromagnetic interaction given by the vertex $`ie(\gamma _\mu F_1(q^2)+(\kappa /2m_e)\sigma _{\mu \rho }q^\rho )`$, where we express $`F_1(q^2)`$ as $`q^2r^2/6`$ in order to extract a limit on a possible charge radius<sup>1</sup><sup>1</sup>1We consider here only Dirac neutrinos since Majorana neutrinos are well known to have quite different electromagnetic properties, in particular, they cannot possess a magnetic moment.. We will include such a contribution only for the tau neutrino because the limits which were already obtained for $`\nu _e`$ and $`\nu _\mu `$ are more stringent than the limit we will be obtaining for the tau neutrino. We obtain the additional contributions from the diagram of Fig.1(b) to the cross section, $`{\displaystyle \frac{d\sigma }{dxdy}}`$ $`=`$ $`{\displaystyle \frac{\alpha ^3}{3}}\{{\displaystyle \frac{2r^2^2s(1x)}{9}}+{\displaystyle \frac{\kappa ^2}{m_e^2}}`$ (2) $`{\displaystyle \frac{g_Vr^2M_Z^2s(1x)(1s(1x)/M_Z^2)}{3\mathrm{sin}^2\theta _W\mathrm{cos}^2\theta _W[s(1x)M_Z^2]^2+(M_Z\mathrm{\Gamma }_Z)^2}}\}`$ $`{\displaystyle \frac{[(1x/2)^2+x^2y^2/4]}{x(1y^2)}}.`$ We have integrated (2) over the relevant range given in Table 1 for each experiment . In Table 1, $`x_T(=E_{T\gamma }/E)`$ is the photon transverse energy normalized to the beam energy, and $`ϵ`$ is the overall efficiency for each data sample. For instance, integrating (2) over the VENUS kinematical region, we changed the variable from $`x`$ to $`x_T`$ with Jacobian, then integrated it over the region $`0.13<x_T<1`$ and $`\mathrm{cos}130.3^{}y\mathrm{cos}50.0^{}`$. We also applied a similar method to the other experiments. Table 2 is a summary of the number of single-photon candidates for each experimental result. It is worth noting that there is no interference between (1) and (2), since the anomalous contribution given in equation (2) flips helicity, but the standard model contribution given in equation (1) does not . We obtained the upper limits on the number of signal events for the observed events and the expected background using two methods: One is Poisson statistics , $$1\alpha =1\frac{e^{(n_B+N)}{\displaystyle \underset{n=0}{\overset{n_0}{}}}{\displaystyle \frac{(n_B+N)^n}{n!}}}{e^{n_B}{\displaystyle \underset{n=0}{\overset{n_0}{}}}{\displaystyle \frac{n_B^n}{n!}}},$$ (3) where $`n_0`$ is the number of the single-photon candidates which each experiment has obtained, $`n_B`$ is the mean for the sum of all backgrounds and $`N`$ is the desired upper limit on the unknown mean for the signal with confidence coefficient $`\alpha `$. The other is the Unified Approach . We applied $`n_B=n_0`$ to both methods because each $`n_0`$ in Table 2 could be explained by the sum of the number of physically expected events and that of non-physical backgrounds. Then, we required $$N_i>\sigma _i\times ϵ_i\times L_i𝑑t$$ (4) for each experiment, where $`N_i`$ is the upper limit at 90% C.L. on the number of signal events, $`\sigma _i`$ is the cross section obtained from integration of (2) over the relevant ranges, $`ϵ_i`$ is the overall efficiency, and $`L_i𝑑t`$ is the integrated luminosity for each experiment, i.e., $`i`$ means VENUS, TOPAZ, and so on. Table 2 contains the upper limits of the tau neutrino magnetic moment for each experiment. We combined the bounds on the magnetic moment at the 90% C.L. at TRISTAN using Poisson statistics, $$\kappa <9.1\times 10^6,$$ (5) and also derived the bounds on the charge radius of the tau neutrino for the experiments (Table 2) and combined them at the 90% C.L., $$r^2<3.1\times 10^{31}\mathrm{cm}^2.$$ (6) In addition, using the Unified Approach, we obtained the combined bound on the magnetic moment at the 90% C.L., $$\kappa <8.0\times 10^6$$ (7) and charge radius $$r^2<2.7\times 10^{31}\mathrm{cm}^2.$$ (8) The obtained results (5)-(8) give upper limits comparable to those obtained from other $`e^+e^{}`$ colliding beam experiments. We have reported on the bound on the tau neutrino magnetic moment and charge radius from the TRISTAN experiments and have obtained the bound from single photon production cross section at TRISTAN, $`9.1\times 10^6\mu _B`$ and $`3.1\times 10^{31}\mathrm{cm}^2`$ at 90% C.L. using Poisson statistics, and $`8.0\times 10^6\mu _B`$ and $`2.7\times 10^{31}\mathrm{cm}^2`$ at 90% C.L. using Unified Approach. They are still far above what is predicted by the standard electroweak theory extended to include massive neutrinos although comparable to the results from other $`e^+e^{}`$ colliding experiments. We would like to thank G.J. Feldman and R.D. Cousins for providing the calculations for the upper end of the signal mean using Unified Approach.
warning/0002/nlin0002037.html	ar5iv	text	# Dispersionless sTB ## 1 Introduction: In recent years, dispersionless integrable models have received a lot of attention. They involve equations of hydrodynamic type \[1-5\] and include such systems as the Riemann equation \[5-7\], the the polytropic gas dynamics , the chaplygin gas and the Born-Infeld equation \[9-10\]. These are models which can be obtained from a “classical” limit of integrable models where the dispersive terms are absent. They have many interesting properties including the fact that, unlike their dispersive counterparts, each of them can be described by a Lax equation which involves a Lax function in the classical phase space and a classical Poisson bracket relation. Even more interesting and more difficult are the supersymmetric dispersionless models. In a recent paper , we gave, for the first time, the Lax description for the dispersionless supersymmetric KdV equation as well as the dispersionless Kupershmidt equation . Unlike the bosonic models, the Lax functions for the dispersionless supersymmetric models do not follow trivially from the Lax operator of the dispersive counterpart. Furthermore, while a lot of the interesting properties follow from the Lax description of the model, we also pointed out several open questions that arise. In this paper, we follow up on our earlier investigation and describe the Lax formulation for the dispersionless supersymmetric two boson (TB) hierarchy. The TB hierarchy as well as its supersymmetric counterpart are known to yield various other integrable models upon appropriate reduction. In this sense, the supersymmetric TB hierarchy is a more interesting model to study. In particular, it has a natural $`N=2`$ supersymmetry \[16-17\] and its dispersionless limit would lead to the first $`N=2`$ supersymmetric model of its kind. There are, in fact, two distinct supersymmetric generalizations of the TB hierarchy. The first is known as the sTB-B hierarchy , so called because it leads upon reduction to the supersymmetric KdV equation considered by the Beckers . The second supersymmetric generalization, on the other hand, leads upon reduction to the supersymmetric KdV equation considered by Manin and Radul . We call this the sTB hierarchy (although, a more appropriate name may be sTB-MR along the same lines). In section 2, we give the Lax description for the dispersionless sTB-B hierarchy, which is fairly straightforward, and bring out its properties as well as the open questions associated with this model. In section 3, we give the Lax description for the dispersionless sTB hierarchy. This is quite nontrivial and naturally reduces to the dispersionless sKdV equation upon appropriate restriction. However, unlike the sTB model, it does not lead to the dispersionless supersymmetric non-linear Schrödinger equation. In fact, even the dispersionless bosonic TB model does not quite give the dispersionless non-linear Schrödinger equation (at least, we have not succeeded in finding a field redefinition which would do this). We bring out various properties of the dispersionless sTB model as well as the open questions associated with this system. In section 4, we describe the dispersionless sTB system in a manifestly $`N=2`$ supersymmetric formulation. Finally, we make some brief observations in section 5. We have used REDUCE and special supersymmetric package in Reduce extensively in some of the algebraic calculations. ## 2 Dispersionless Limit of sTB-B Equation: The sTB hierarchy , like the TB hierarchy, is an integrable system in $`1+1`$ dimensions. The basic dynamical variables for this system are the two fermionic superfields $`\mathrm{\Phi }_0(t,x,\theta )`$ $`=`$ $`\psi _0+\theta J_0`$ $`\mathrm{\Phi }_1(t,x,\theta )`$ $`=`$ $`\psi _1+\theta J_1`$ (1) where $`\theta `$ represents a Grassmann coordinate and we are suppressing the space-time dependence on the right hand side for simplicity. The sTB-B hierarchy is described by the non-standard Lax equation $$\frac{L}{t_n}=[L,(L^n)_1]$$ (2) where $`n=1,2,\mathrm{}`$ and the Lax operator has the form $$L=D^2(D\mathrm{\Phi }_0)+D^2(D\mathrm{\Phi }_1),$$ (3) with the super-covariant derivative defined to be $$D=\frac{}{\theta }+\theta \frac{}{x},D^2=\frac{}{x}$$ (4) Explicitly, the first three flows of the sTB-B hierarchy have the form $$\frac{\mathrm{\Phi }_0}{t}=\mathrm{\Phi }_{0x},\frac{\mathrm{\Phi }_1}{t}=\mathrm{\Phi }_{1x},$$ (5) $`{\displaystyle \frac{\mathrm{\Phi }_0}{t}}`$ $`=`$ $`(D^4\mathrm{\Phi }_0)+D((D\mathrm{\Phi }_0)^2+2(D\mathrm{\Phi }_1))),`$ $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`(D^4\mathrm{\Phi }_1)+2D\left((D\mathrm{\Phi }_0)(D\mathrm{\Phi }_1)\right),`$ (6) $`{\displaystyle \frac{\mathrm{\Phi }_0}{t}}`$ $`=`$ $`\mathrm{\Phi }_{0xxx}D\left(6(D\mathrm{\Phi }_0)(D\mathrm{\Phi }_1)3(D\mathrm{\Phi }_0)(D\mathrm{\Phi }_{0x})+(D\mathrm{\Phi }_0)^3\right),`$ $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`\mathrm{\Phi }_{1xxx}3D\left((D\mathrm{\Phi }_1)^2+(D\mathrm{\Phi }_1)(D\mathrm{\Phi }_0)^2+(D\mathrm{\Phi }_{1x})(D\mathrm{\Phi }_0)2(D\mathrm{\Phi }_1)(D\mathrm{\Phi }_{0x})\right).`$ (7) The last two equations lead, under the reduction $`\mathrm{\Phi }_0=0`$, to the supersymmetric KdV equation considered by the Beckers (the so called sKdV-B equation). $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`\mathrm{\Phi }_{1xxx}3D\left((D\mathrm{\Phi }_1)^2\right).`$ (8) The Lax description for the dispersionless limit of the sTB-B hierarchy is quite straightforward, much like the dispersionless limit of the sKdV-B hierarchy . Consider the Lax function $$L=p(D\mathrm{\Phi }_0)+p^1(D\mathrm{\Phi }_1)$$ (9) where $`p`$ is the momentum variable of the classical phase space, satisfying the canonical PB relations $$\{x,p\}=1,\{x,x\}=0=\{p,p\}.$$ (10) Then, it is easily seen, with the standard canonical Poisson bracket relations ($`\{p,f\}=\frac{df}{dx}`$), that the Lax equation $$\frac{L}{t_n}=\{(L^n)_1,L\}$$ (11) leads to the dispersionless sTB-B hierarchy. Explicitly, the first three flows of this hierarchy are , $$\frac{\mathrm{\Phi }_0}{t}=\mathrm{\Phi }_{0x},\frac{\mathrm{\Phi }_1}{t}=\mathrm{\Phi }_{1x},$$ (12) $$\frac{\mathrm{\Phi }_0}{t}=2(D\mathrm{\Phi }_0)\mathrm{\Phi }_{0x}+2\mathrm{\Phi }_{1x},\frac{\mathrm{\Phi }_1}{t}=2D((D\mathrm{\Phi }_0)(D\mathrm{\Phi }_1)),$$ (13) $`{\displaystyle \frac{\mathrm{\Phi }_0}{t}}`$ $`=`$ $`D\left((D\mathrm{\Phi }_0)^36(D\mathrm{\Phi }_1)(D\mathrm{\Phi }_0)\right)`$ $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`3D\left((D\mathrm{\Phi }_1)^2+(D\mathrm{\Phi }_1)(D\mathrm{\Phi }_0)^2\right).`$ (14) Equation (14) allows the reduction ($`\mathrm{\Phi }_0=0`$) to the dispersionless supersymmetric KdV-B equation $$\frac{\mathrm{\Phi }_1}{t}=3D((D\mathrm{\Phi }_1)^2).$$ (15) whose properties we have studied earlier in detail . From the Lax description, the conserved quantities of the hierarchy can be easily obtained. In fact, the normalized conserved quantities can be written as $`H_n`$ $`=`$ $`{\displaystyle \frac{(1)^{n+1}}{n}}{\displaystyle 𝑑z\mathrm{Res}L^n}`$ (16) $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{m=0}{\overset{n}{}}}{}_{}{}^{n}C_{nm}^{nm}C_{m+1}{\displaystyle 𝑑z(D\mathrm{\Phi }_0)^{n2m1}(D\mathrm{\Phi }_1)^{m+1}}`$ Here $`dz=dxd\theta `$ represents the integration over the superspace. Explicitly, the first few conserved quantities are $`H_1`$ $`=`$ $`{\displaystyle 𝑑z(D\mathrm{\Phi }_1)}=0`$ $`H_2`$ $`=`$ $`{\displaystyle 𝑑z(D\mathrm{\Phi }_0)(D\mathrm{\Phi }_1)}`$ $`H_3`$ $`=`$ $`{\displaystyle 𝑑z[(D\mathrm{\Phi }_0)^2+(D\mathrm{\Phi }_1)](D\mathrm{\Phi }_1)}`$ and so on. As it stands, it is clear that these conserved quantities are fermionic. This is a peculiarity of the sTB-B system (for that matter any -B system) that the Hamiltonians are fermionic. Correspondingly, the Hamiltonian structures are odd and we note the first two structures for completeness, namely, $$𝒟_1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),𝒟_2=\left(\begin{array}{cc}2& D(D\mathrm{\Phi }_0)D^1\\ D^1(D\mathrm{\Phi }_0)D& D^1(D\mathrm{\Phi }_1)D+D(D\mathrm{\Phi }_1)D^1\end{array}\right)$$ (17) While the Hamiltonian structure $`𝒟_1`$ can be obtained from the Lax function as the standard Gelfand-Dikii bracket, we do not know how to obtain the second structure from the Lax function. Furthermore, the Jacobi identity for this structure is complicated and needs to be checked. However, these two Hamiltonian structures do lead to the recursion operator $$R=𝒟_1^1𝒟_2=\left(\begin{array}{cc}D^1(D\mathrm{\Phi }_0)D& D^1(D\mathrm{\Phi }_1)D+D(D\mathrm{\Phi }_1)D^1\\ 2& D(D\mathrm{\Phi }_0)D^1\end{array}\right)$$ (18) which can be easily checked to connect the successive Hamiltonians of the hierarchy. Supersymmetric integrable systems have conserved non-local charges \[20-21\] and the dispersionless sTB-B hierarchy also has conserved non-local charges. For example, it can be checked that $`Q_n`$ $`=`$ $`{\displaystyle \frac{(1)^{n+1}}{n}}{\displaystyle 𝑑z(D^1\mathrm{Res}L^n)}`$ (19) $`=`$ $`{\displaystyle \frac{1}{n}}{\displaystyle \underset{m=0}{\overset{n}{}}}{}_{}{}^{n}C_{nm}^{nm}C_{m+1}{\displaystyle 𝑑z(D^1((D\mathrm{\Phi }_0)^{n2m1}(D\mathrm{\Phi }_1)^{m+1}))}`$ with the first few charges $`Q_1`$ $`=`$ $`{\displaystyle 𝑑z\mathrm{\Phi }_1}`$ $`Q_2`$ $`=`$ $`{\displaystyle 𝑑z(D^1((D\mathrm{\Phi }_0)(D\mathrm{\Phi }_1)))}`$ $`Q_3`$ $`=`$ $`{\displaystyle 𝑑z(D^1((D\mathrm{\Phi }_0)^2(D\mathrm{\Phi }_1)+(D\mathrm{\Phi }_1)^2))}`$ (20) and so on, are conserved under the flow of the system. These are bosonic charges and, interestingly enough, these non-local charges are also related to one another by the same recursion operator $`R`$ in eq. (18), namely, $$\left(\begin{array}{c}\frac{\delta Q_{n+1}}{\delta \mathrm{\Phi }_0}\\ \frac{\delta Q_{n+1}}{\delta \mathrm{\Phi }_1}\end{array}\right)=R\left(\begin{array}{c}\frac{\delta Q_n}{\delta \mathrm{\Phi }_0}\\ \frac{\delta Q_n}{\delta \mathrm{\Phi }_1}\end{array}\right)$$ (21) It is also nice to see explicitly that if we set $`\mathrm{\Phi }_0=0`$ and $`\mathrm{\Phi }_1=\mathrm{\Phi }`$, then, all the even charges vanish and the odd charges, namely, $`H_{2n+1}`$ and $`Q_{2n+1}`$ coincide with the corresponding charges of the sKdV-B system , as they should. ## 3 Dispersionless Limit of sTB Equation: In terms of the same basic variables, $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_1`$, the sTB hierarchy is described by the Lax operator $$L=D^2(D\mathrm{\Phi }_0)+D^1\mathrm{\Phi }_1,$$ (22) and the non-standard Lax equation $$\frac{L}{t_n}=[L,(L^n)_1].$$ (23) Explicitly, the lowest order equations have the form $$\frac{\mathrm{\Phi }_0}{t}=\mathrm{\Phi }_{0x},\frac{\mathrm{\Phi }_1}{t}=\mathrm{\Phi }_{1x},$$ (24) $`{\displaystyle \frac{\mathrm{\Phi }_0}{t}}`$ $`=`$ $`(D^4\mathrm{\Phi }_0)+D((D\mathrm{\Phi }_0)^2+2(D\mathrm{\Phi }_1)),`$ $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`(D^4\mathrm{\Phi }_1)+2D^2\left((D\mathrm{\Phi }_0)\mathrm{\Phi }_1\right),`$ (25) $`{\displaystyle \frac{\mathrm{\Phi }_0}{t}}`$ $`=`$ $`(D^6\mathrm{\Phi }_0)+D\left(3\mathrm{\Phi }_1\mathrm{\Phi }_{0x}6(D\mathrm{\Phi }_1)(D\mathrm{\Phi }_0)(D\mathrm{\Phi }_0)^3+3(D\mathrm{\Phi }_0)(D\mathrm{\Phi }_{0x})\right),`$ $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`(D^6\mathrm{\Phi }_1)3D^2\left(\mathrm{\Phi }_1(D\mathrm{\Phi }_0)^2+(D\mathrm{\Phi }_0)\mathrm{\Phi }_{1x}+\mathrm{\Phi }_1(D\mathrm{\Phi }_1)\right).`$ (26) The last equation allows the reduction ($`\mathrm{\Phi }_0=0`$) to the supersymmetric KdV equation considered by Manin and Radul , namely, $$\frac{\mathrm{\Phi }_1}{t}=(D^6\mathrm{\Phi }_1)3D^2\left(\mathrm{\Phi }_1(D\mathrm{\Phi }_1)\right)$$ (27) The dispersionless limit of this Lax operator is not as straightforward. However, with some work, it can be determined that the Lax function $$L=p(D\mathrm{\Phi }_0)+p^1(D\mathrm{\Phi }_1)p^2\mathrm{\Phi }_{0x}\mathrm{\Phi }_1+p^3\mathrm{\Phi }_{1x}\mathrm{\Phi }_1$$ (28) and the classical Lax equation $$\frac{L}{t_n}=\{(L^n)_1,L\}$$ (29) give the dispersionless sTB hierarchy whose first three equations have the explicit forms $$\frac{\mathrm{\Phi }_0}{t}=\mathrm{\Phi }_{0x},\frac{\mathrm{\Phi }_1}{t}=\mathrm{\Phi }_{1x},$$ (30) $$\frac{\mathrm{\Phi }_0}{t}=2(D\mathrm{\Phi }_0)\mathrm{\Phi }_{0x}+2\mathrm{\Phi }_{1x},\frac{\mathrm{\Phi }_1}{t}=2D^2\left((D\mathrm{\Phi }_0)\mathrm{\Phi }_1\right)$$ (31) $`{\displaystyle \frac{\mathrm{\Phi }_0}{t}}`$ $`=`$ $`D\left(3\mathrm{\Phi }_1\mathrm{\Phi }_{0x}6(D\mathrm{\Phi }_1)(D\mathrm{\Phi }_0)(D\mathrm{\Phi }_0)^3\right),`$ $`{\displaystyle \frac{\mathrm{\Phi }_1}{t}}`$ $`=`$ $`3D^2\left(\mathrm{\Phi }_1(D\mathrm{\Phi }_0)^2+\mathrm{\Phi }_1(D\mathrm{\Phi }_1)\right).`$ (32) There are several things to note from here. First, eq. (32), upon setting $`\mathrm{\Phi }_0=0`$ and $`\mathrm{\Phi }_1=\mathrm{\Phi }`$, gives the dispersionless sKdV equation , namely, $$\frac{\mathrm{\Phi }}{t}=3D^2\left(\mathrm{\Phi }(D\mathrm{\Phi })\right).$$ (33) This, therefore, gives the non-standard representation of the dispersionless sKdV equation (as opposed to the standard representation given in ), and is analogous to the reduction of the sTB hierarchy to the sKdV hierarchy. Second, the normalized conserved quantities of this system can be easily determined to be $`H_n`$ $`=`$ $`{\displaystyle \frac{(1)^{n+1}}{n}}{\displaystyle 𝑑zD^1(\mathrm{Res}L^n)}`$ (34) $`=`$ $`{\displaystyle \underset{m=0}{\overset{m_{max}}{}}}{\displaystyle \frac{(nm)!}{m!(m1)!(n2m+1)!}}{\displaystyle 𝑑z(D\mathrm{\Phi }_0)^{n2m+1}(D\mathrm{\Phi }_1)^{m1}\mathrm{\Phi }_1}`$ where the upper limit $`m_{max}=\frac{n}{2}`$ if $`n`$ is even, but $`m_{max}=[\frac{n}{2}]+1`$ if $`n`$ is odd. The first few of these conserved charges have the explicit forms $`H_1`$ $`=`$ $`{\displaystyle 𝑑z\mathrm{\Phi }_1}`$ $`H_2`$ $`=`$ $`{\displaystyle 𝑑z(D\mathrm{\Phi }_0)\mathrm{\Phi }_1}`$ $`H_3`$ $`=`$ $`{\displaystyle 𝑑z((D\mathrm{\Phi }_0)^2+(D\mathrm{\Phi }_1))\mathrm{\Phi }_1}`$ (35) and so on. These conserved charges are all bosonic and it is clear that if we set $`\mathrm{\Phi }_0=0`$ and $`\mathrm{\Phi }_1=\mathrm{\Phi }`$, all the even charges vanish while the odd charges coincide with those of the dispersionless sKdV hierarchy as they should. We have not been able to derive the Hamiltonian structures for this system from the Gelfand-Dikii formalism (as is the case in the dispersionless sKdV also). However, the first Hamiltonian structure can be easily checked to be $$𝒟_1=\left(\begin{array}{cc}0& D\\ D& 0\end{array}\right)$$ (36) and this trivially satisfies the Jacobi identity. Notice that this Hamiltonian operator defines a closed skew symmetric two-form $$\mathrm{\Omega }(𝒟_1)(a,b)=𝑑z(a_1(Db_2)+a_2(Db_1))$$ (37) where a,b are arbitrary, two component (column matrix) bosonic superfields. We have also checked that it is impossible to construct the second Hamiltonian operator out of local operators alone. However, there is a possibility, which we have not checked, to construct such an operator out of local as well as non-local operators. Let us now discuss some of the outstanding questions associated with such a system. First, we have not been able to construct the non-local charges from this Lax function. By brute force, we have checked that charges, such as $`Q_n`$ $`=`$ $`{\displaystyle 𝑑z(D^1\mathrm{\Phi }_1)^n},n=1,2,3,\mathrm{}`$ $`Q_2^{}`$ $`=`$ $`{\displaystyle 𝑑z\left[\frac{3}{2}(D^1\mathrm{\Phi }_1)^2+\frac{1}{2}\mathrm{\Phi }_0\mathrm{\Phi }_1(D^1((D\mathrm{\Phi }_0)\mathrm{\Phi }_1))\right]}`$ (38) and so on, are conserved under the flow. Furthermore, under the substitution, $`\mathrm{\Phi }_0=0`$ and $`\mathrm{\Phi }_1=\mathrm{\Phi }`$, these reduce to the appropriate non-local charges of the dispersionless sKdV (In this limit, $`Q_2=\frac{2}{3}Q_2^{}`$). However, for lack of a systematic procedure for constructing these charges, we do not have a general expression for the $`n`$th charge of this set in terms of the Lax function. Furthermore, we do know that the dispersionless sKdV has a second set of non-local charges and since this system can be obtained from the dispersionless sTB, it would be natural to expect the dispersionless sTB also to have a second set of non-local charges. However, we do not know if such a set exists. This emphasizes the need for a systematic understanding of the non-local charges in such systems, described a classical Lax function. It is also worth recalling that the sTB equation yields the sNLS equation (supersymmetric non-linear Schrödinger equation) under the redefinition $$\mathrm{\Phi }_0=(D\mathrm{ln}(DQ))+(D^1(\overline{Q}Q)),\mathrm{\Phi }_1=\overline{Q}(DQ)$$ (39) However, we would like to note here that we have not been able to find a redefinition of fields which would take the dispersionless sTB equation to the dispersionless sNLS equation. In fact, it is worth pointing out that it is not a difficulty only for the supersymmetric system. Even the dispersionless TB equation (bosonic) does not appear to yield the dispersionless NLS equation under the corresponding redefinition. This question certainly needs further study. ## 4 $`N=2`$ Formulation for the dispersionless sTB Equation: The sTB hierarchy has a $`N=2`$ supersymmetry . However, it is not manifest in the description of the last section. In order to see the $`N=2`$ supersymmetry manifestly, we have to define the basic variables of the theory appropriately. Let us recall that in the conventional description of the system (as given in the previous section), the basic variables were two fermionic superfields which depended on the usual bosonic coordinates, $`(x,t)`$, but they also depended on an additional anti-commuting Grassmann variable $`\theta `$. The Taylor expansion of such superfields in the Grassmann coordinate is simple and has been given in eq. (1). In the presence of an $`N=2`$ supersymmetry, the superfields depend on $`(x,t)`$ as well as two anti-commuting variables $`\theta _1`$ and $`\theta _2`$. (Incidentally, an $`N=2`$ supersymmetric system can also be described on an $`N=1`$ superspace which is the description used in the earlier section.) Expanding the superfield $`\mathrm{\Phi }(x,t,\theta _1,\theta _2)`$ in a Taylor series in the anti-commuting variables, we obtain $$\mathrm{\Phi }=\varphi _1+\theta _1\chi _1+\theta _2\chi _2+\theta _2\theta _1\varphi _2,$$ (40) where $`\varphi _1,\varphi _2`$ are bosonic functions (or fermionic ) and $`\chi _1,\chi _2`$ are fermionic (or bosonic ) if $`\mathrm{\Phi }`$ is a bosonic (or fermionic) superfield. In this case, the super-covariant derivatives on the $`N=2`$ superspace are defined as $`D_1`$ $`=`$ $`{\displaystyle \frac{}{\theta _1}}+\theta _1{\displaystyle \frac{}{x}},D_2={\displaystyle \frac{}{\theta _2}}+\theta _2{\displaystyle \frac{}{x}},`$ $`D_1^2`$ $`=`$ $`D_2^2=={\displaystyle \frac{}{x}},D_2D_1=D_1D_2.`$ (41) Recently several methods have been proposed to obtain $`N=2`$ supersymmetric soliton equations \[22-25\]. In this paper, we consider two different supersymmetric Lax operators which generate two distinct sTB hierarchies. The first is connected with the $`N=2,a=1`$ supersymmetric KdV equation (namely, it reduces to this model upon appropriate redefinition) and is defined by the Lax operator $$L=+\stackrel{~}{\mathrm{\Phi }}_0+^1D_1D_2\stackrel{~}{\mathrm{\Phi }}_1.$$ (42) where $`\stackrel{~}{\mathrm{\Phi }}_0`$ and $`\stackrel{~}{\mathrm{\Phi }}_1`$ are two bosonic superfields (and we have used a tilde to avoid confusion with the superfields of the earlier sections). This Lax operator is related to the supersymmetric Lax operator considered in , through the gauge transformation $$L=e^gLe^g,$$ (43) where $`g=𝑑x𝑑\theta _1𝑑\theta _2G=𝑑ZG`$ with $`G`$ a bosonic superfield. It is straightforward to check that the Lax operator of eq. (42) leads to consistent dynamical equations through the nonstandard Lax relation $$\frac{L}{t_n}=[L,(L^n)_1].$$ (44) Explicitly, we can write the first three flows as $$\frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}=\stackrel{~}{\mathrm{\Phi }}_{0x},\frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}=\stackrel{~}{\mathrm{\Phi }}_{1x},$$ (45) $$\frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}=(\stackrel{~}{\mathrm{\Phi }}_{0x}\stackrel{~}{\mathrm{\Phi }}_0^22(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)+\stackrel{~}{\mathrm{\Phi }}_1^2),\frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}=(\stackrel{~}{\mathrm{\Phi }}_{1x}2\stackrel{~}{\mathrm{\Phi }}_1\stackrel{~}{\mathrm{\Phi }}_0)$$ (46) $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`(\stackrel{~}{\mathrm{\Phi }}_{0xx}3\stackrel{~}{\mathrm{\Phi }}_{ox}\stackrel{~}{\mathrm{\Phi }}_0\stackrel{~}{\mathrm{\Phi }}_0^33(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0+3\stackrel{~}{\mathrm{\Phi }}_1^2\stackrel{~}{\mathrm{\Phi }}_03(D_1D_2(\stackrel{~}{\mathrm{\Phi }}_1\stackrel{~}{\mathrm{\Phi }}_0))),`$ $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}}`$ $`=`$ $`(\stackrel{~}{\mathrm{\Phi }}_{1xx}3(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_1+\stackrel{~}{\mathrm{\Phi }}_1^3+3\stackrel{~}{\mathrm{\Phi }}_{1x}\stackrel{~}{\mathrm{\Phi }}_03\stackrel{~}{\mathrm{\Phi }}_1\stackrel{~}{\mathrm{\Phi }}_0^2).`$ (47) The last equation leads, upon the reduction $`\stackrel{~}{\mathrm{\Phi }}_0=0`$, to the $`N=2,a=1`$ supersymmetric KdV equation $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}}=(\stackrel{~}{\mathrm{\Phi }}_{1xx}3(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_1+\stackrel{~}{\mathrm{\Phi }}_1^3).`$ (48) We can also construct a second Lax operator of the form $$L=D_1+\overline{\mathrm{\Phi }}_0+^1D_2\overline{\mathrm{\Phi }}_1,$$ (49) where $`\overline{\mathrm{\Phi }}_0`$ is a fermionic superfield while $`\overline{\mathrm{\Phi }}_1`$ is a bosonic superfield. Note that the Lax operator, in the present case, is fermionic while that in eq. (42) was bosonic. Nonetheless, as in the previous case, it is easy to check that this Lax operator is gauge equivalent to the one considered in . This Lax operator leads to dynamical equations of the non-standard form $$\frac{L}{t_{2n}}=[L,(L^{2n})_1].$$ (50) Explicitly, the first three flows of this hierarchy have the forms $$\frac{\overline{\mathrm{\Phi }}_0}{t}=\overline{\mathrm{\Phi }}_{0x},\frac{\overline{\mathrm{\Phi }}_1}{t}=\overline{\mathrm{\Phi }}_{1x},$$ (51) $`{\displaystyle \frac{\overline{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`D_1\left((D_1\overline{\mathrm{\Phi }}_{0x})2(D_1D_2\overline{\mathrm{\Phi }}_1)2(D_1\overline{\mathrm{\Phi }}_0)^22\overline{\mathrm{\Phi }}_1^2\overline{\mathrm{\Phi }}_1(D_2\overline{\mathrm{\Phi }}_0)\right),`$ $`{\displaystyle \frac{\overline{\mathrm{\Phi }}_1}{t}}`$ $`=`$ $`D_1\left((D_1\overline{\mathrm{\Phi }}_{1x})2\overline{\mathrm{\Phi }}_{0x}\overline{\mathrm{\Phi }}_1(D_2\overline{\mathrm{\Phi }}_1)\overline{\mathrm{\Phi }}_12(D_1\overline{\mathrm{\Phi }}_1)(D_1\overline{\mathrm{\Phi }}_0)\right),`$ (52) $`{\displaystyle \frac{\overline{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`D_1((D_1\overline{\mathrm{\Phi }}_{0xx}3(D_1\overline{\mathrm{\Phi }}_{0x})(D_1\overline{\mathrm{\Phi }}_0)6(D_1D_2\overline{\mathrm{\Phi }}_1)(D_1\overline{\mathrm{\Phi }}_0)6\overline{\mathrm{\Phi }}_1^2(D_1\overline{\mathrm{\Phi }}_0)`$ $`6\overline{\mathrm{\Phi }}_1(D_2\overline{\mathrm{\Phi }}_0)(D_1\overline{\mathrm{\Phi }}_0)+3(D_2\overline{\mathrm{\Phi }}_1\overline{\mathrm{\Phi }}_{0x}+3(D_1\overline{\mathrm{\Phi }}_1)(D_1D_2\overline{\mathrm{\Phi }}_0)(D_1\overline{\mathrm{\Phi }}_0)^3),`$ $`{\displaystyle \frac{\overline{\mathrm{\Phi }}_1}{t}}`$ $`=`$ $`D_1((D_1\overline{\mathrm{\Phi }}_{1xx}+3\overline{\mathrm{\Phi }}_{0x}\overline{\mathrm{\Phi }}_{1x}6\overline{\mathrm{\Phi }}_{0x}\overline{\mathrm{\Phi }}_1(D_1\overline{\mathrm{\Phi }}_0)+3(D_2\overline{\mathrm{\Phi }}_1)\overline{\mathrm{\Phi }}_{1x}`$ (53) $`6(D_2\overline{\mathrm{\Phi }}_1)\overline{\mathrm{\Phi }}_1(D_1\overline{\mathrm{\Phi }}_0)+3(D_1\overline{\mathrm{\Phi }}_{1x})(D_1\overline{\mathrm{\Phi }}_0)3(D_1\overline{\mathrm{\Phi }}_1)(D_1\overline{\mathrm{\Phi }}_0)^2`$ $`3(D_1\overline{\mathrm{\Phi }}_1)(D_1D_2\overline{\mathrm{\Phi }}_1)6(D_1\overline{\mathrm{\Phi }}_1)\overline{\mathrm{\Phi }}_1^26(D_1\overline{\mathrm{\Phi }}_1)\overline{\mathrm{\Phi }}_1(D_2\overline{\mathrm{\Phi }}_0))`$ Note that when $`\overline{\mathrm{\Phi }}_0=0`$, the last equation reduces to the $`N=2,a=2`$ supersymmetric KdV equation. Unfortunately the Lax operator, which gives rise to the sTB hierarchy that contains the $`N=2,a=4`$ supersymmetric KdV equation is not known as yet. Furthermore, since fermionic Lax operators do not lend easily to a dispersionless limit, we will not discuss this system any further. In order to obtain the dispersionless sTB hierarchy, we now have to introduce the concept of the fermionic momenta on the $`N=2`$ superspace. There will be two such fermionic momenta defined by $$\mathrm{\Pi }_1=(p_{\theta _1}+\theta _1p),\mathrm{\Pi }_2=(p_{\theta _2}+\theta _2p),$$ (54) We can now assume the “commutation” rules for the functions $`\mathrm{\Pi }_i`$ as $$\{\mathrm{\Pi }_i,\mathrm{\Pi }_j\}=2p\delta _{ij}$$ (55) Note that the $`\mathrm{\Pi }_i`$’s generate covariant differentiation through the PB relation (for any superfield $`A`$) $$\{\mathrm{\Pi }_i,A\}=(D_iA),i=1,2.$$ (56) In trying to obtain the dispersionless $`N=2`$ sTB hierarchy, we start from the Lax operator for the $`N=2`$ sTB hierarchy in eq. (42), and assume the Lax function for the dispersionless system to have the general form $$L=p+k_1\mathrm{\Pi }_1\mathrm{\Pi }_2+\stackrel{~}{\mathrm{\Phi }}_0+\underset{s=1}{\overset{3}{}}p^s\left(F_0^s+\mathrm{\Pi }_1F_1^s+\mathrm{\Pi }_2F_2^s+\mathrm{\Pi }_1\mathrm{\Pi }_2F_3^s\right).$$ (57) where $`k_1`$ is an arbitrary coefficient and $`F_k^s,k=0,1,2,3,s=1,2,3`$ are arbitrary functions of $`\stackrel{~}{\mathrm{\Phi }}_0`$ and $`\stackrel{~}{\mathrm{\Phi }}_1`$. We have checked that the classical analogue of eq. (44), namely, $$\frac{L}{t_n}=\{(L^n)_1,L\}$$ (58) with the projection $`1`$ defined as $`\left({\displaystyle \underset{s=\mathrm{}}{\overset{\mathrm{}}{}}}p^s\left(F_0^s+\mathrm{\Pi }_1F_1^s+\mathrm{\Pi }_2F_2^s+\mathrm{\Pi }_1\mathrm{\Pi }_2F_3^3\right)\right)_1=`$ $`\mathrm{\Pi }_1F_1^s+\mathrm{\Pi }_2F_2^s+\mathrm{\Pi }_1\mathrm{\Pi }_2F_3^s+{\displaystyle \underset{s=1}{\overset{\mathrm{}}{}}}p^s\left(F_0^s+\mathrm{\Pi }_1F_1^s+\mathrm{\Pi }_2F_2^s+\mathrm{\Pi }_1\mathrm{\Pi }_2F_3^3\right).`$ (59) leads to two possible solutions with $`k_1=0`$ in either case. The first Lax function that leads to consistent equations has the form $$L=p+\stackrel{~}{\mathrm{\Phi }}_0+p^1\mathrm{\Pi }_1\mathrm{\Pi }_2\stackrel{~}{\mathrm{\Phi }}_1$$ (60) The first three flows of this hierarchy have the explicit forms $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`\stackrel{~}{\mathrm{\Phi }}_{0x},{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}}=\stackrel{~}{\mathrm{\Phi }}_{1x},`$ $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`(\stackrel{~}{\mathrm{\Phi }}_0^2),{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}}=2(\stackrel{~}{\mathrm{\Phi }}_1\stackrel{~}{\mathrm{\Phi }}_0)`$ $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`(\stackrel{~}{\mathrm{\Phi }}_0^3),{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}}=3(\stackrel{~}{\mathrm{\Phi }}_1\stackrel{~}{\mathrm{\Phi }}_0^2)`$ (61) However, this hierarchy appears trivial since the equations do not have any explicit dependence on supersymmetric covariant derivatives. Therefore, we will not consider this hierarchy any further. The second Lax function does not contain the fermionic functions $`\mathrm{\Pi }_i`$ and has the form $$L=p+\stackrel{~}{\mathrm{\Phi }}_0+p^1(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)p^2(D_2\stackrel{~}{\mathrm{\Phi }}_1)(D_1\stackrel{~}{\mathrm{\Phi }}_0)+p^3(D_2\stackrel{~}{\mathrm{\Phi }}_{1x})(D_2\stackrel{~}{\mathrm{\Phi }}_1)$$ (62) This, on the other hand, produces an interesting supersymmetric hierarchy whose first three flows have the forms $$\frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}=\stackrel{~}{\mathrm{\Phi }}_{0x},\frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}=\stackrel{~}{\mathrm{\Phi }}_{1x},$$ (63) $$\frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}=(2(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)+\stackrel{~}{\mathrm{\Phi }}_0^2),\frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}=2D_2((D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0)$$ (64) $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_0}{t}}`$ $`=`$ $`(\stackrel{~}{\mathrm{\Phi }}_0^3+6(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_03(D_2\stackrel{~}{\mathrm{\Phi }}_1)(D_1\stackrel{~}{\mathrm{\Phi }}_0)),`$ (65) $`{\displaystyle \frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}}`$ $`=`$ $`3D_2((D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0^2+(D_2\stackrel{~}{\mathrm{\Phi }}_1)(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)).`$ (66) It is interesting to note that, when $`\stackrel{~}{\mathrm{\Phi }}_0=0`$, the last equation becomes $$\frac{\stackrel{~}{\mathrm{\Phi }}_1}{t}=3D_2((D_2\stackrel{~}{\mathrm{\Phi }}_1)(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)).$$ (67) which, in fact, reduces to eq. (33) with the substitution $`(D_2\stackrel{~}{\mathrm{\Phi }}_1)\|_{\theta _2=0}=\mathrm{\Phi }_1`$. Therefore, eq. (66) can be considered as the $`N=2`$ generalization of the dispersionless $`N=1`$ supersymmetric KdV-MR equation. Let us also note that the Lax function in eq. (62) as well as the resulting hierarchy coincide with those of the previous section with the identifications $$(D_2\stackrel{~}{\mathrm{\Phi }}_1)\|_{\theta _2=0}=\mathrm{\Phi }_1,\stackrel{~}{\mathrm{\Phi }}_0\|_{\theta _2=0}=(D_1\mathrm{\Phi }_0)=(D\mathrm{\Phi }_0).$$ (68) Namely, the redefinition also involves $`\theta _1\theta `$ or equivalently, $`D_1D`$. Notice that these transformations are highly nontrivial reductions. Consequently, we do not expect the conserved quantities of eqs. (34) to define conserved quantities of the $`N=2`$ supersymmetric hierarchy and we have verified this. On the other hand, we can construct conserved quantities for the $`N=2`$ hierarchy directly and we have found two such series of conserved quantities, by brute force, using the computer. The first set consists of bosonic conserved charges of the form $`H_n`$ $`=`$ $`{\displaystyle 𝑑Z(\stackrel{~}{\mathrm{\Phi }}_1)^n},n=1,2,3,4,\mathrm{},`$ $`\stackrel{~}{H}_1`$ $`=`$ $`{\displaystyle 𝑑Z\stackrel{~}{\mathrm{\Phi }}_0}`$ $`\stackrel{~}{H}_2`$ $`=`$ $`{\displaystyle 𝑑Z\stackrel{~}{\mathrm{\Phi }}_1\stackrel{~}{\mathrm{\Phi }}_0}`$ (69) where $`dZ=dxd\theta _1d\theta _2`$. The second series of conserved charges is fermionic of the form $`H_{5/2}`$ $`=`$ $`{\displaystyle 𝑑Z(D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0}`$ $`H_{7/2}`$ $`=`$ $`{\displaystyle 𝑑Z(D_2\stackrel{~}{\mathrm{\Phi }}_1)(\stackrel{~}{\mathrm{\Phi }}_0^2+(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1))},`$ $`H_{9/2}`$ $`=`$ $`{\displaystyle 𝑑Z(D_2\stackrel{~}{\mathrm{\Phi }}_1)(3(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0+\stackrel{~}{\mathrm{\Phi }}_0^3)},`$ $`H_{11/2}`$ $`=`$ $`{\displaystyle }dZ(4(D_2\stackrel{~}{\mathrm{\Phi }}_{1x})\stackrel{~}{\mathrm{\Phi }}_{1x}\stackrel{~}{\mathrm{\Phi }}_14(D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_{1xx}\stackrel{~}{\mathrm{\Phi }}_1+12(D_1\stackrel{~}{\mathrm{\Phi }}_{1x})(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_1`$ (71) $`+18(D_2\stackrel{~}{\mathrm{\Phi }}_1)(D_1D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0^2+3(D_2\stackrel{~}{\mathrm{\Phi }}_1)\stackrel{~}{\mathrm{\Phi }}_0^4)`$ Here, we have labelled the conserved quantities by the weights of the integrand ($`[\stackrel{~}{\mathrm{\Phi }}_0]=1=[\stackrel{~}{\mathrm{\Phi }}_1]`$ and $`[D_1]=\frac{1}{2}=[D_2]`$). Note that the second set of charges in eq. (69) follows from the Lax function, up to non-essential normalization, as $$H_{n+\frac{1}{2}}=𝑑ZD_1^1\mathrm{Res}L^n$$ (72) However, we do not know how to obtain the first set (except for the lowest one) from the Lax function, nor is it clear that these exhaust all the conserved charges of the system. Comparing with the discussion of the previous section and particularly from the form of the non-local conserved charges in eq. (38), we see that we can write non-local conserved charges for the $`N=2`$ system as $$\stackrel{~}{Q}_n=𝑑Z(D_1^1D_2\stackrel{~}{\mathrm{\Phi }}_1)^nn=1,2,3,4\mathrm{}$$ (73) It is, in fact, quite straightforward to check that they are conserved. Similarly, we note that $$\stackrel{~}{Q}_2^{}=𝑑Z(D_1^1(\stackrel{~}{\mathrm{\Phi }}_0(D_2\stackrel{~}{\mathrm{\Phi }}_1))\frac{1}{2}(D_1^1\stackrel{~}{\mathrm{\Phi }}_0)(D_2\stackrel{~}{\mathrm{\Phi }}_1)\frac{3}{2}\stackrel{~}{\mathrm{\Phi }}_1^2),$$ (74) also represents a conserved charge (compare with the second of the charges in eq. (38)). However, we do not know how to obtain these from the Lax function directly. The first Hamiltonian structure for the dispersionless $`N=2`$ hierarchy of eqs. (63)-(69) is easily seen to be $$\stackrel{~}{𝒟}_1=\left(\begin{array}{cc}0& D_2\\ D_2& 0\end{array}\right)$$ (75) and is trivially seen to satisfy the Jacobi identity. This Hamiltonian operator defines the closed skew symmetric two-form $$\mathrm{\Omega }(\stackrel{~}{𝒟}_1)(a,b)=𝑑Z(a_1(Db_2)a_2(Db_1))$$ (76) where in contrast to the $`N=1`$ case $`a`$ and $`b`$ are arbitrary, two component fermionic superfields. ## 5 Conclusion In this paper, we have studied the dispersionless limits of the sTB-B, the sTB as well as the $`N=2`$ sTB hierarchies in detail. We have obtained the Lax descriptions in terms of classical Lax functions, obtained conserved local as well as some of the non-local charges and brought out various other features associated with such systems. We have also tried to point out various open questions associated with such systems, the most pressing of which is a systematic understanding of the construction of non-local charges for such systems starting from the Lax description as well as a generalization of the Gelfand-Dikii procedure for construction of Hamiltonian structures for such systems.. ## Acknowledgments A.D. acknowledges support in part by the U.S. Dept. of Energy Grant DE-FG 02-91ER40685 while Z.P. is supported in part by the Polish KBN Grant 2 P0 3B 136 16.
warning/0002/astro-ph0002257.html	ar5iv	text	# Neural networks and separation of Cosmic Microwave Background and astrophysical signals in sky maps ## 1 Introduction Maps produced by large area surveys aimed at imaging primordial fluctuations of the Cosmic Microwave Background (CMB) contain a linear mixture of signals by several astrophysical and cosmological sources (Galactic synchrotron, free-free and dust emissions, both from compact and diffuse sources, extragalactic sources, Sunyaev-Zeldovich effect in clusters of galaxies or by inhomogeneous re-ionization, in addition to primary and secondary CMB anisotropies) convolved with the spatial and spectral responses of the antenna and of the detectors. In order to exploit the unique cosmological information encoded in the CMB anisotropy patterns as well as the extremely interesting astrophysical information carried by the foregound signals, we need to accurately separate the different components. A great deal of work has been carried out in recent years in this area (see de Oliveira-Costa & Tegmark 1999, and references therein; Tegmark et al. 2000). The problem of map denoising has been tackled with the wavelets analysis on the whole sphere \[Tenorio et al. 1999\] and on sky patches \[Sanz et al. 1999b\]. Algorithms to single out the CMB and the various foregrounds have been developed \[Bouchet et al. 1999, Hobson et al. 1998, \]. In these works, Wiener filtering (WF) and the maximum entropy method (MEM) have been applied to simulated data from the Planck satellite, taking into account the expected performances of the instruments. Assuming a perfect knowledge of the frequency dependence of all the components, as well as priors for the statistical properties of their spatial pattern, these algorithms are able to recover the the strongest components, at the best Planck resolution. We adopt a rather different approach, considering denoising and deconvolution of the signals on one side and component separation on the other as separate steps in the data analysis process, and focus here on the latter step only, presenting a ’blind separation’ method, based on ’Independent Component Analysis’ (ICA) techniques. The method does not require any a priori assumption on spectral properties and on the spatial distribution of the various components, but only that they are statistically independent and all but at most one have a non-Gaussian distribution. It is important to note that this is in fact the physical system we have to deal with: surely all the foregrounds are non-Gaussian, while the CMB is expected to be a nearly Gaussian fluctuation field for most of the candidate theories of the early universe. The paper is organized as follows. In Section 2 we introduce the relevant formalism and briefly review methods applied in previous works. In Section 3 we outline the ICA algorithm in a rather general framework, since it may be useful for a variety of astrophysical applications. In Section 4 we describe our simulated maps. In Section 5 we give some details on our analysis and present the results. In Section 6 we draw our conclusions and list some future developments. ## 2 Formalism and previous approaches We assume that the frequency spectrum of radiation components (referred to as sources) is independent of the position in the sky. Since we deal here with relatively small patches of the sky, we adopt Cartesian coordinates, $`(\xi ,\eta )`$. The function describing the i-th source then writes $$\stackrel{~}{s}_i(\xi ,\eta ,\nu )=s_i(\xi ,\eta )_i(\nu )i=1,\mathrm{},N$$ (1) where $`N`$ is the number of independent sources and $`_i(\nu )`$ is the emission spectrum. The signal received from the point $`(\xi ,\eta )`$ in the sky is $$\stackrel{~}{x}(\xi ,\eta ,\nu )=\underset{i=1}{\overset{N}{}}s_i(\xi ,\eta )_i(\nu )$$ (2) Suppose that the instrument has $`M`$ channels, with spectral response functions $`t_j(\nu )`$, $`j=1,\mathrm{}M`$ centered at different frequencies, and that the beam patterns are independent of frequency within each passband. Let beam patterns be described by the space-invariant PSF’s $`h_j(\xi ,\eta )`$, so that the maps are produced by a linear convolutional mechanism. (Note that this is an additional simplifying assumption since in real experiments a position dependent defocussing related to the chosen scanning strategy may occur.) Then, the map yielded by j<sup>th</sup> channel is: $$x_j(\xi ,\eta )=h_j(\xi x,\eta y)t_j(\nu )$$ $$\underset{i=1}{\overset{N}{}}s_i(x,y)_i(\nu )dxdyd\nu +ϵ_j(\xi ,\eta )=$$ $$=\stackrel{~}{x}_j(\xi ,\eta )h_j(\xi ,\eta )+ϵ_j(\xi ,\eta ),j=1,\mathrm{},M,$$ (3) where: $$\stackrel{~}{x}_j(\xi ,\eta )=\underset{i=1}{\overset{N}{}}a_{ji}s_i(xi,\eta ),j=1,\mathrm{},M,$$ (4) $$a_{ji}=_i(\nu )t_j(\nu )𝑑\nu ,j=1,\mathrm{},M;i=1,\mathrm{},N,$$ (5) $``$ denotes linear convolution and $`ϵ_j(\xi ,\eta )`$ represents the instrumental noise. Eq. (4) can also be written in matrix form: $$\stackrel{~}{\text{x}}(\xi ,\eta )=A\text{s}(\xi ,\eta )$$ (6) where the entries of the $`M\times N`$ matrix $`A`$ are given by Eq. (5). The unknowns of our problem are the $`N`$ functions $`s_i(\xi ,\eta )`$, and the data set is made of the $`M`$ maps $`x_j(\xi ,\eta )`$ in Eq. (3). Besides the measured data, we also know the instrument beam-patterns $`h_j(\xi ,\eta )`$, and, more or less approximately (depending on the specific source), the coefficients $`a_{ji}`$ in Eq. (4). Eq. (3) can be easily rewritten in the Fourier space: $$X_j(\omega _\xi ,\omega _\eta )=\underset{i=1}{\overset{N}{}}R_{ji}(\omega _\xi ,\omega _\eta )S_i(\omega _\xi ,\omega _\eta )+_j(\omega _\xi ,\omega _\eta ),$$ (7) where the capital letters denote the Fourier transforms of the corresponding lowercase functions, and $$R_{ji}(\omega _\xi ,\omega _\eta )=_j(\omega _\xi ,\omega _\eta )a_{ji},$$ (8) $`_j`$ being the Fourier transform of the beam profile $`h_j`$. Eq. (7) can thus be rewritten in matrix form: $$𝐗=R𝐒+.$$ (9) The above equation must be satisfied by each Fourier mode $`(\omega _\xi ,\omega _\eta )`$, independently. The aim is to recover the true signals $`S_i(\omega _\xi ,\omega _\eta )`$ constituting the column vector $`𝐒`$. If the matrix $`A`$ in Eq. (6) is known exactly then, in the absence of noise, the problem reduces to a linear inversion of Eq. (9) for each Fourier mode. In practice, however, $`_j`$ vanishes for some Fourier mode. For these modes the entire j-th row of the matrix $`R`$ also vanishes, and $`R`$ may become a non-full-rank matrix. An inversion based on statistical approaches built on a priori knowledge is thus needed. In the following two subsections we briefly describe two such approaches, and in the third one we briefly recall a technique so far mostly exploited for the denoising problem and for extraction of extragalactic sources. ### 2.1 The maximum entropy approach The Maximum Entropy Method (MEM) for the reconstruction of images is based on a Bayesian approach to the problem (Skilling 1988, 1989; Gull 1988). Let $`𝐗`$ be a vector of $`M`$ observations whose probability distribution $`P(𝐗\|𝐒)`$ depends on the values of $`N`$ quantities $`𝐒=S_1,\mathrm{},S_N`$. Let us $`P(𝐒)`$ be the prior probability distribution of $`𝐒`$, telling us what is known about $`𝐒`$ without knowledge of the data. Given the data $`𝐗`$, Bayes’ theorem states that the conditional distribution of $`𝐒`$ (the posterior distribution of $`𝐒`$) is given by the product of the likelihood of the data, $`P(𝐗\|𝐒)`$, with the prior: $$P(𝐒\|𝐗)=zP(𝐗\|𝐒)P(𝐒),$$ (10) where $`z`$ is a normalization constant. An estimator $`\widehat{𝐒}`$ of the true signal vector can be constructed by maximizing the posterior probability $`P(𝐒\|𝐗)P(𝐗\|𝐒)P(𝐒)`$. However, while the likelihood in Eq. (10) is easily determined once the noise and signal covariance matrices are known, the appropriate choice of the prior distribution for the model considered is a major problem in the Bayesian approach: since Bayes’ theorem is simply a rule for manipulating probabilities, it cannot by itself help us to assign them in the first place, so one has to look elsewhere. The MEM is a consistent variational method for the assignment of probabilities under certain types of constraints that must refer to the probability distribution directly. The Maximum Entropy principle states that if one has some information $`I`$ on which the probability distribution is based, one can assign a probability distribution to a proposition $`i`$ such that $`P(i\|I)`$ contains only the information $`I`$ that one actually possesses. This assignment is done by maximizing the Entropy $$H\underset{i=1}{\overset{N}{}}P(i\|I)logP(i\|I)$$ (11) It can be seen that when nothing is known except that the probability distribution should be normalized, the Maximum Entropy principle yields the uniform prior. In our case the proposition $`i`$ represents S, and the information $`I`$ is the assumption of signal statistical independence. The standard application of the method considered strictly positive signals (Skilling 1988, 1989; Gull 1988); the extension to the case of CMB temperature fluctuations, which can be both positive and negative, was worked out by Hobson et al. (1998). The construction of the entropic prior requires, in general, the knowledge of the frequency dependence of the components to be recovered as well as of the signal covariance matrix $`𝐂(𝐤)=<𝐒(𝐤)𝐒^{}(𝐤)>`$, with the average taken on all the possible realizations. ### 2.2 The multifrequency Wiener filtering If a Gaussian prior is adopted, the Bayesian approach gives the multifrequency Wiener filtering (WF) solution \[Bouchet et al. 1999\]. In in this case too an estimator of the signal vector is obtained by maximizing the posterior probability in Eq. (10), given the signal covariance matrix $`𝐂(𝐤)`$. The Gaussian prior probability distribution for the signal has the form $$P(𝐒)\mathrm{exp}(𝐒^{}𝐂^1𝐒).$$ (12) The estimator $`\widehat{\mathrm{S}}`$ is linearly related to the data vector $`\widehat{\mathrm{X}}`$ through the Wiener matrix $`𝐖(𝐂^1+𝐑^{}𝐍^1𝐑)^1`$, where $`𝐑`$ corresponds to the matrix in (9) and $`𝐍(𝐤)=<ϵ(𝐤)ϵ^{}(𝐤)>`$ is the noise covariance matrix: $$\widehat{𝐒}=\mathrm{𝐖𝐗}.$$ (13) The $`𝐖`$ matrix has the role of a linear filter; again, its construction requires an a priori knowledge of the spectral behavior of the signals. This method is endangered by the clear non-Gaussianity of foregrounds. ### 2.3 Wavelet methods The development of wavelet techniques for signal processing has been very fast in the last ten years \[see, e.g., Jawerth et al. 1994\]. The wavelet approach is conceptually very simple: whereas the Fourier transform is highly inefficient in dealing with the local behavior, the wavelet transform is able to introduce a good space-frequency localization, thus providing information on the contributions coming from different positions and scales. In one dimension, we can define the analyzing wavelet as $`\mathrm{\Psi }(x;R,b)R^{1/2}\psi [(xb)/R]`$, dependent on two parameters, dilation ($`R`$) and translation ($`b`$); $`\psi (x)`$ is a one-dimensional function satisfying the following conditions: a) $`_{\mathrm{}}^{\mathrm{}}𝑑x\psi (x)=0`$, b) $`_{\mathrm{}}^{\mathrm{}}𝑑x\psi ^2(x)=1`$ and c) $`_{\mathrm{}}^{\mathrm{}}𝑑k\|k\|^1\psi ^2(k)<\mathrm{}`$, where $`\psi (k)`$ is the Fourier transform of $`\psi (x)`$. The wavelet $`\mathrm{\Psi }`$ operates as a mathematical microscope of magnification $`R^1`$ at the space point $`b`$. The wavelet coefficients associated to a one-dimensional function $`f(x)`$ are: $$w(R,b)=𝑑xf(x)\mathrm{\Psi }(x;R,b).$$ (14) The computationally faster algorithms for the wavelet analysis of 2-dimensional images are those based on Multiresolution analysis \[Mallat 1989\] or on 2D wavelet analysis \[Lemarié & Meyer 1986\], using tensor products of one-dimensional wavelets. The discrete Multiresolution analysis entails the definition of a one-dimensional scaling function $`\varphi `$, normalized as $`_{\mathrm{}}^{\mathrm{}}𝑑x\varphi (x)=1`$ \[Ogden et al. 1997\]. Scaling functions act as low-pass filters whereas wavelet functions single out one scale. The 2D wavelet method \[Sanz et al. 1999b\] is based on two scales, providing therefore more information on different resolutions (defined by the product of the two scales) than the Multiresolution one. Recently, wavelet techniques have been introduced in the analysis of CMB data. Denoising of CMB maps has been performed on patches of the sky of $`12^{}.8\times 12^{}.8`$ using either multiresolution techniques \[Sanz et al. 1999a\] and 2D wavelets \[Sanz et al. 1999b\], as well as on the whole celestial sphere \[Tenorio et al. 1999\]. As a first step, maps with the cosmological signal plus a Gaussian instrumental noise have been considered. Denoising of CMB maps has been carried out by using a signal–independent prescription, the SURE thresholding method \[Donoho & Johnstone 1995\]. The results are model independent and only a good knowledge of the noise affecting the observed CMB maps is required, whereas nothing has to be assumed on the nature of the underlying field(s). Moreover, wavelet techniques are highly efficient in localizing noise variations and features in the maps. The wavelet method is able to improve the signal-to-noise ratio by a factor of 3 to 5; correspondingly, the error on $`C_{\mathrm{}}`$’s derived from denoised maps is about 2 times lower than that obtained with the WF method. Wavelets were also successfully applied to the detection of point sources in CMB maps in the presence of the cosmological signal and of instrumental noise \[Tenorio et al. 1999\]; more recently, successfull results on source detection have also been obtained in presence of diffuse galactic foregrounds \[Cayón et al. 2000\]. The results are comparable to those obtained with the filtering method presented by \[Tegmark & de Oliveira-Costa (1998)\] which, however, relies on the assumption that all the underlying fields are Gaussian. ## 3 The ICA approach We present here a rather different approach, characterized by the capability of working ‘blindly’ i.e. without prior knowledge of spectral and spatial properties of the signals to be separated. The method is of interest for a broad variety of signal and image processing applications, i.e. whenever a number of source signals are detected by multiple transducers, and the transmission channels for the sources are unknown, so that each transducer receives a mixture of the source signals with unknown scaling coefficients and channel distortion. In this exploratory study we confine ourselves to the case of simple linear combinations of unconvolved source signals \[Amari & Chichocki 1998, Bell & Sejnowski 1995\]. The problem can be stated as follows: a set of $`N`$ signals is input to an unknown frequency dependent multiple-input-multiple-output linear instantaneous system, whose $`M`$ outputs are our observed signals. We use the term instantaneous to denote a system whose output at a given point only depends on the input signals at the same point. Our objective is to find a stable reconstruction system to estimate the original input signals with no prior assumptions either on the signal distributions or on their frequency scalings. The problem in its general form is normally unsolvable, and a “working hypothesis” must be made. The hypothesis we make is the mutual statistical independence of our source signals, whatever their actual distributions are. Several solutions have been proposed for this problem, each based on more or less sound principles, not all of which are typical of classical signal processing. Indeed, information theory, neural networks, statistics and probability have played an important part in the development of these techniques. We do not consider here specific instrumental features like beam convolution and noise contamination, leaving the specialization of the ICA method to specific experiments to future work; this allows us to highlight the capabilities of this approach, able to work in conditions where other algorithms would not be viable. Therefore, we adopt Eq. (6) as our data model, just dropping the tilde accent on vector $`𝐱`$. Also, the instrumental noise term in Eq. (7) will be neglected. It can be proved that, to solve the problem described above, the following hypotheses should be verified \[Amari & Chichocki 1998, Comon 1994\]: * All the source signals are statistically independent; * At most one of them has a Gaussian distribution; * $`MN`$; * Low noise. The last two assumptions can be somewhat relaxed by choosing suitable separation strategies. As far as independence is concerned, roughly speaking, we may say that the search for an ICA model from non-ICA data (i.e. data not coming from independent sources) should give the most ‘interesting’ (namely, the most structured) projections of the data \[Hyvärinen & Oja 1999, Friedman 1987\]. This is not equivalent to say that separation is achieved; however, we have seen from our experiments that a good separation can be obtained even for sources that are not totally independent. The second assumption above tells us that Gaussian sources cannot be separated. More specifically, they can only be separated up to an orthogonal transformation. In fact, it can be shown that the joint probability of a mixture of Gaussian signals is invariant to orthogonal transformations. This means that if independent components are found from Gaussian mixtures, then any orthogonal transformation of them gives mutually independent components. Many strategies have been adopted to solve the separation problem on the basis of the above hypotheses, all based on looking for a set of independent signals, which can be shown to be the original sources. A formal criterion to test independence, from which all the separating strategies can be derived, is described later in this section. In order to recover the original source signals from the observed mixtures, we use a separating scheme in the form of a feed-forward neural network. The observed signals are input to an $`N\times M`$ matrix $`W`$, referred to as the the synaptic weight matrix, whose adjustable entries, $`w_{ij},i=1,\mathrm{}N,j=1,\mathrm{}M`$, are updated for every sample of the input vector $`\text{x}(\xi ,\eta )`$ (at step $`\tau `$) following a suitable learning algorithm. The output of matrix $`W`$ at step $`\tau `$ will be: $$\text{u}(\xi ,\eta ,\tau )=W(\tau )\text{x}(\xi ,\eta ).$$ (15) $`W(\tau )`$ is expected to converge to a true separating matrix, that is, a matrix whose output is a copy of the inputs, for every point $`(\xi ,\eta )`$. Ideally, this final matrix $`W`$ should be such that $`WA=I`$, where $`I`$ is the $`N\times N`$ identity. As an example, if $`M=N`$, we should have $`W=A^1`$. There are, however, two basic indeterminacies in our problem: ordering and scaling. Even if we are able to extract $`N`$ independent sources from $`M`$ linear mixtures, we cannot know a priori the order in which they will be arranged, since this corresponds to unobservable permutations of the columns of matrix $`A`$. Moreover, the scales of the extracted signals are unknown, because when a signal is multiplied by some scalar constant, the effect is the same as of multiplying by the same constant the corresponding column of the mixing matrix. This means that $`W(\tau )`$ will converge, at best, to a matrix $`W`$ such that: $$WA=PD,$$ (16) where $`P`$ is any $`N\times N`$ permutation matrix, and $`D`$ is a nonsingular diagonal scaling matrix. From Eqs. (6), (15) and (16) we thus have: $$𝐮=W𝐱=WA𝐬=PD𝐬.$$ (17) That is, as anticipated, each component of $`𝐮`$ is a scaled version of a component of $`𝐬`$, not necessarily in the same order. This is not a serious inconvenience in our application, since we should be able to recover the proper scales for the separated sources from other pieces of information, for example matching with independent lower resolution observations like those of COBE on the case of MAP and Planck. If $`A`$ was known, the performance of the separation algorithm could be evaluated by means of the matrix $`WA`$. If the separation is perfect, this matrix has only one nonzero element for each row and each column. In any non-ideal situation each row and column of $`WA`$ should contain only one dominant element. In all the cases treated here we assume $`MN`$, but we consider the case where $`N`$, although smaller than $`M`$, is not known. The mutual statistical independence of the source signals can be expressed in terms of a separable joint probability density function $`q(𝐬)`$: $$q(𝐬)=\underset{j=1}{\overset{N}{}}q_j(s_j)$$ (18) where $`q_j(s_j)`$ is the marginal probability density of the $`j^{th}`$ source. Various algorithms can be used to learn the matrix $`W`$. All these algorithms can be derived from a unified principle based on the Kullback–Leibler (KL) divergence between the joint probability density of the output vector $`𝐮`$, $`p_U(𝐮)`$, and a function $`q(𝐮)`$, which should be suitably chosen among the ones of the type of Eq. (18). The KL divergence between the two functions mentioned above may be written as a function of the matrix $`W`$, and can be considered as a cost function in the sense of Bayesian statistics: $$R(W)=p_U(𝐮)\mathrm{log}\frac{p_U(𝐮)}{q(𝐮)}𝑑𝐮.$$ (19) It can be proved that, under mild conditions on $`q(𝐮)`$, $`R(W)`$ has a global minimum where $`W`$ is such that $`WA=PD`$. The different possible choices for $`q(𝐬)`$ lead to the different particular learning strategies proposed in the literature \[Amari & Chichocki 1998, Yang & Amari 1997, Bell & Sejnowski 1995\]. The uniform gradient search method, which is a gradient-type algorithm, takes into account the Riemannian metric structure of our objective parameter space, which is the set of all nonsingular matrices $`W`$ \[Amari & Chichocki 1998\]. In a general case, where the number $`N`$ of sources is only known to be smaller than the number of observations, the following formula is derived: $$W(\tau +1)=W(\tau )+$$ $$+\alpha (\tau )[\mathrm{\Lambda }𝐮(\tau )𝐮^T(\tau )𝐟(𝐮(\tau ))𝐮^T(\tau )]W(\tau ),$$ (20) where $`\mathrm{\Lambda }`$ is a $`M\times M`$ diagonal matrix: $$\mathrm{\Lambda }=\mathrm{diag}[(u_1+f_1(u_1))u_1]\mathrm{}[(u_M+f_M(u_M))u_M].$$ (21) Pixel by pixel, the $`M\times M`$ matrix $`W`$ is multiplied by the M–vector x, and gives vector $`𝐮`$ as its output. This output is transformed through the nonlinear vector function $`𝐟(𝐮)`$, and the result is combined with $`𝐮`$ itself to build the update to matrix $`W`$, through Eq. (20). The process has to be iterated by reading the data maps several times. If $`N`$ is strictly smaller than $`M`$, then $`MN`$ outputs can be shown to rapidly converge to zero, or to pure noise functions. The convergence properties of this iterative formula are shown to be independent of the particular matrix $`A`$, so that, even a strongly ill-conditioned system does not affect the convergence of the learning algorithm. In other words, even when the contributions from some components are very small, there is no problem to recover them. This property is called the equivariant property since the asymptotic properties of the algorithm are independent of the mixing matrix. The $`\tau `$-dependent parameter $`\alpha `$ is the learning rate; its value is normally decreased during the iteration. As far as the choice of $`\alpha (\tau )`$ is concerned, a strategy to learn it and its annealing scheme is given in Amari et al. (1998); we have chosen $`\alpha (\tau )`$ decreasing from $`10^3`$ to $`10^4`$ linearly with the number of iterations. The final problem is how to choose the function $`𝐟(𝐮)`$. If we know the true source distributions $`q_j(u_j)`$, the best choice is to make $`f_j^{}(u_j)=q_j(u_j)`$, since this gives the maximum likelihood estimator. However, the point is that when $`q_j(u_j)`$ are specified incorrectly, the algorithm gives the correct answer under certain conditions. In any case, the choice of $`𝐟(𝐮)`$ should be made to ensure the existence of an equilibrium point for the cost function and the stability of the optimization algorithm. These requirements can be satisfied even though the nonlinearities chosen are not optimal. A suboptimal choice for sub-Gaussian source signals (negative kurtosis), is: $$f_i(u_i)=\beta u_i+u_i\|u_i\|^2,$$ (22) and, for super-Gaussian source signals (positive kurtosis): $$f_i(u_i)=\beta u_i+\mathrm{tanh}(\gamma u_i),$$ (23) where $`\beta 0`$ and $`\gamma 2`$; if one source is Gaussian, the above choices remain viable as well. In our case, we verified that all the source functions except CMB are super-Gaussian, and thus we implemented the learning algorithm following Eq.(20), with the nonlinearities in Eq. (23), and $`\beta =0`$, $`\gamma =2`$. As already stated, the mean of the input signal at each frequency is subtracted. In previous works \[Yang & Amari 1997\] the initial matrix was chosen as $`WI`$; in that analysis, the image data consisted of a set of components with nearly the same amplitude. The initial guess for $`W`$ affects the computation time, as well as the scaling of the reconstructed signals and their order. Interestingly, we found that adjusting the diagonal elements so that they roughly reflect the different weights of the components in the mixture can speed-up the convergence. For the problem at hand, the results shown in § 5 have been obtained starting from $`W=`$diag, for the case of a $`4\times 4W`$-matrix, and using only 20 learning steps: the time needed was about 1 minute on a UltraSparc machine, equipped with an 300 MHz UltraSparc processor, 256 MBytes RAM, running down SUN Solaris 7 Operating System, compiling the FORTRAN 90 code using SUN Fortran Workshop 5.0 ## 4 Simulated maps We produced simulated maps of the antenna temperature distribution with 3’.5 pixel size of a $`15^{}\times 15^{}`$ region centered at $`l=90^{}`$, $`b=45^{}`$, at the four central frequencies of the Planck/LFI channels \[Mandolesi et al. 1998\], namely 30, 44, 70 and 100 GHz (Fig. 1). The HEALPix pixelization scheme \[see Górski et al. 1999\] was adopted. The maps include CMB anisotropies, Galactic synchrotron and dust emissions, and extragalactic radio sources. CMB fluctuations correspond to a flat Cold Dark Matter (CDM) model ($`\mathrm{\Omega }_{CDM}=.95`$, $`\mathrm{\Omega }_b=.05`$, three massless neutrino species), normalized to the COBE data \[see Seljak & Zaldarriaga 1996\]. As it is well known, the CMB spectrum, in terms of antenna temperature, writes: $$s_{CMB}^{antenna}(\xi ,\eta ,\nu )=s_{CMB}^{thermod.}(\xi ,\eta )\frac{\stackrel{~}{\nu }^2e^{\stackrel{~}{\nu }}}{(e^{\stackrel{~}{\nu }}1)^2},$$ (24) where $`\stackrel{~}{\nu }=\nu /56.8`$ and $`\nu `$ is the frequency in GHz; $`s_{CMB}^{thermod.}(\xi ,\eta )`$ is frequency independent \[Fixsen et al. 1996\]. As for Galactic synchrotron emission, we have extrapolated the 408 MHz map with about 1 degree resolution \[Haslam et al. 1982\], assuming a power law spectrum, in terms of antenna temperature: $$_{syn}\stackrel{~}{\nu }^{n_s},$$ (25) with spectral index $`n_s=2.9`$. The dust emission maps with about 6’ resolution constructed by Schlegel et al. (1998) combining IRAS and DIRBE data have been used as templates for Galactic dust emission. The extrapolation to Planck/LFI frequencies was done assuming a grey-body spectrum: $$_{dust}\frac{\stackrel{~}{\nu }^{m+1}}{e^{\stackrel{~}{\nu }}1},$$ (26) with $`m=2`$, $`\stackrel{~}{\nu }=h\nu /kT_{dust}`$, $`T_{dust}`$ being the dust temperature. Although, in general, $`T_{dust}`$ varies across the sky, it turns out to be approximately constant at about $`18`$K in the region considered here; we have therefore adopted this value in the above equation. Because of the lack of a suitable template, we have ignored here free-free emission, which may be important particularly at 70 and 100 GHz. This component needs to be included in future work. The model by Toffolatti et al. (1998) was adopted for extragalactic radio sources, assumed to have a Poisson distribution. An antenna temperaure spectral index $`n_{\mathrm{rs}}=1.9`$ was adopted $`(_{\mathrm{rs}}\stackrel{~}{\nu }^{n_{\mathrm{rs}}})`$. ## 5 Blind analysis and results As it is well known, the strongest signals at the Planck/LFI frequencies come from the CMB and from radio sources (although the latter show up essentially as a few high peaks), whereas synchrotron emission and thermal dust are roughly 1 or 2 orders of magnitude lower, depending on frequency. Thus we are testing the performances of the ICA algorithm with four signals exhibiting very different spatial patterns, frequency dependences and amplitudes. Since we are interested in the fluctuation pattern, the mean of the total signal (sum of the four components) is set to zero at each frequency. We adopt a “blind” approach: no information on either the spatial distribution or the frequency dependence of the signals is provided to the algorithm. The reconstructed maps of the the four components are shown in Fig. 2. Several interesting features may be noticed. The order of the plotted maps is permuted with respect to the input maps in Fig. 1, reflecting the order of the ICA outputs: the first output is synchrotron, the second represents radio sources, the third is CMB and the fourth is dust. All the output maps look very similar to the true ones; even synchrotron lower resolution pixels have been reproduced. In Figs. 3, 5, 4 and 6 we analyze the goodness of the separation by comparing power spectra and showing scatter plots between the inputs and the outputs. ### 5.1 Signal reconstruction For each map, we have computed the angular power spectrum, defined by the expansion coefficients $`C_{\mathrm{}}`$ of the two point correlation function in Legendre polynomials. As is well known, it can conveniently be expressed in terms of the coefficients of the expansion of the signal $`S`$ into spherical harmonics, $`S(\theta ,\varphi )=_\mathrm{}ma_\mathrm{}mY_\mathrm{}m(\theta ,\varphi )`$: $$C_{\mathrm{}}=\frac{1}{2\mathrm{}+1}\underset{m}{}\|a_\mathrm{}m\|^2.$$ (27) Such coefficients are useful because from elementary properties of the Legendre polynomials it can be seen that the coefficient $`C_{\mathrm{}}`$ quantifies the amount of perturbation on the angular scale $`\theta `$ given by $`\theta 180/\mathrm{}`$ degrees. The panels on the top of Figs. 3, 4, 5, 6 show the power spectra of the input (left) and output (right) signals. The CMB exhibits the characteristic peaks on sub-degree angular scales due to acoustic oscillations of the photon-baryon fluid at decoupling; the dashed line represents the theoretical model from which the map was generated, while the solid line is the power spectrum of our simulated patch: the difference between the two curves is due to the sample variance corresponding to the CMB Gaussian statistics. Radio sources are completely different, having all the power on small scales with the typical shot noise spatial pattern; dust and synchrotron emissions have power decreasing on small scales roughly as a power law, as expected \[Mandolesi et al. 1998, Puget et al. 1998\]. The left-hand side panels on the bottom show the quality factor, defined as the ratio between true and reconstructed power spectrum coefficients, for each multipole $`\mathrm{}`$. Due to the limited size of the analyzed region, the power spectrum can be defined on scales below roughly $`2^{}`$. The bottom right-hand side panels are scatter plots of the ICA results: for each pixel of the maps, we plotted the value of the reconstructed image vs. the corresponding input value. The reconstructed signals have zero mean and are in unit of the constant $`d`$ multiplying each output map, produced during the separation phase, as mentioned in §$`\mathrm{\hspace{0.17em}3}`$: the scale of each signal is unreproducible for a blind separation algorithm like ICA. Nevertheless, a lot of information is encoded into the spatial pattern of each signals, and ultimately its overall normalization could be recovered exploiting data from other experiments. Therefore, the relation between each true signal and its reconstruction is $$s_i^{in}=ds_i^{out}+b,i=1,\mathrm{},N_{pixels},$$ (28) where $`b`$ represents merely the mean of the input signal, that is zero for the CMB and some positive value for the foregrounds. To quantify the quality of the reconstruction, we have recovered $`d`$ and $`b`$ by performing a linear fit of output to input maps (s<sup>in</sup>,s<sup>out</sup>) for each signal: $$d=\frac{\underset{i}{}s_i^{in}s_i^{out}\overline{s}^{in}\underset{i}{}s_i^{out}}{_i(s_i^{out})^2\overline{s}^{out}_is_i^{out}},b=\overline{s}^{in}d\overline{s}^{out},$$ (29) where the sums run over all the pixels, and the bar indicates the average value over the patch; the values of $`d`$ and $`b`$, as well as the linear fits (dashed lines), are indicated for all the signals in the scatter plot panels. Also, in the same panels we show the standard deviation of the fit, that is $$\sigma =\left[\frac{1}{N_{pixels}}\underset{i}{}(s_i^{in}ds_i^{out}b)^2\right]^{1/2}.$$ (30) A comparison of such quantity with the input signals (bottom right-hand side panels) gives an estimate of the goodness of the reconstruction. CMB and radio sources are recovered with percent and $`0.1\%`$ precision, respectively, while the accuracy drops roughly to $`10\%`$ for the (much weaker) Galactic components, synchrotron and dust. Also, the latter appear to be slightly mixed; this is likely due to the fact that they are somewhat correlated so that the hypothesis of statistical independence is not properly satisfied. We have also tested to what extent the counts of radio sources are recovered. This was done in terms of the relative flux $$\mathrm{\Delta }s=s/s_{\mathrm{max}},$$ (31) $`s_{\mathrm{max}}`$ being the flux of the brightest source. In Fig. 7 we show the cumulative number of input (dashed) and output (solid line) pixels exceeding a given value of $`\mathrm{\Delta }s`$. The algorithm correctly recovers essentially all sources with $`\mathrm{\Delta }s2\times 10^2`$, corresponding to a signal of $`T_s50\mu `$K, or to a flux density $`S=(2k_BT_s/\lambda ^2)\mathrm{\Delta }\mathrm{\Omega }15`$mJy, where $`k_B`$ the Boltzmann constant, $`\lambda `$ the wavelength and $`\mathrm{\Delta }\mathrm{\Omega }`$ the solid angle covered by a pixel, that is $`3.5^{}\times 3.5^{}10^6`$sr. At fainter fluxes the counts are overestimated; this is probably due to the contamination from the other signals. In any case, the flux limit for source detection is surprisingly low, even lower of the rms CMB fluctuations ($`\sigma _{CMB}70\mu `$K at the resolution limit of our maps), substantially lower or at least comparable to that achieved with other methods which require stronger assumptions \[Cayón et al. 2000, Hobson et al. 1998\]. This high efficiency in detecting point sources illustrates the ability of the method in taking the maximum advantage of the differences in frequency and spatial properties of the various components. On the other hand, we stress that our approach is idealized in a number of aspects: beam convolution and instrumental noise have not been taken into account, and the same frequency scaling has been assumed for all radio sources. Therefore more detailed investigations are needed to estimate a realistic source detection limit. Finally note that the quality of the separation is similar on all scales, as shown by the bottom left-hand side panels of Figs. 3, 4, 5, 6. The exception are radio sources, whose true power spectrum goes to zero at low $`\mathrm{}`$’s more rapidly than the reconstructed one. ### 5.2 Reconstruction of the frequency dependence Another asset of this technique is the possibility of recovering the frequency dependence of individual components. The outputs can be written as $`𝐮=W𝐱`$, where $`𝐱=A𝐬`$. As previously mentioned, in the ideal case $`WA`$ would be a diagonal matrix containing the constants $`d`$ for all the signals, multiplied by a permutation matrix. It can be easily seen that, if this is true, the frequency scalings of all the components can be obtained by inverting the matrix $`W`$ and performing the ratio, column by column, of each element with the one corresponding to the row corresponding to a given frequency. However, as pointed out in §$`\mathrm{\hspace{0.17em}3}`$, if some signals are much smaller than others the above reasoning is only approximately valid. This is precisely what is happening in our case: we are able to accurately recover the frequency scaling of the strongest signals, CMB and radio sources, while the others are lost (see Table LABEL:frequencyin). ## 6 Concluding remarks and future developments We have developed a neural network suitable to implement the Independent Component Analysis technique for separating different emission components in maps of the sky at microwave wavelengths. The algorithm was applied to simulated maps of a $`15^{}\times 15^{}`$ region of sky at 30, 44, 70, 100 GHz, corresponding to the frequency channels of Planck’s Low Frequency Instrument (LFI). Simulations include the Cosmic Microwave Background, extragalactic radio sources and Galactic synchrotron and thermal dust emission. The various components have markedly different angular patterns, frequency dependences and amplitudes. The technique exploits the statistical independence of the different signals to recover each individual component with no prior assumption either on their spatial pattern or on their frequency dependence. The great virtue of this approach is the capability of the algorithm to learn how to recover the independent components in the input maps. The price of the lack of a priori information is that each signal can be recovered multiplied by an unknown constant produced during the learning process itself. However this is not a substantial limitation, since a lot of physics is encoded in the spatial patterns of the signals, and ultimately the right normalization of each component can be obtained by resorting to independent observations. The results are very promising. The CMB map is recovered with an accuracy at the 1% level. The algorithm is remarkably efficient also in the detection of extragalactic radio sources: almost all sources brighter then 15 mJy at 100 GHz (corresponding to $`0.7\sigma _{CMB}`$, $`\sigma _{CMB}`$ being the rms level of CMB fluctuations on the pixel scale) are recovered; on the other hand, it must be stressed that is not directly indicative of what can be achieved in the analysis of Planck/LFI data because the adopted resolution ($`3^{}.5\times 3^{}.5`$) is much better than that of the real experiment, instrumental noise has been neglected and the same spectral slope was assumed for all sources. Also the frequency dependences of the strongest components are correctly recovered (error on the spectral index of 1% for the CMB and extragalactic sources). Maps of subdominant signals (Galactic synchrotron and dust emissions) are recovered with rms errors of about 10%; their spectral properties cannot be retrieved by our technique. The reconstruction has equal quality on all the scales of the input maps, down to the pixel size. All this indicates that this technique is suitable for a variety of astrophysical applications, i.e. whenever we want to separate independent signals from different astrophysical processes occurring along the line of sight. Of course, much work has to be done to better explore the potential of the ICA technique. It has to be tested under more realistic assumptions, taking into account instrumental noise and the effect of angular response functions as well as including a more complete and accurate characterization of foregrounds. In particular, the assumption that the spectral properties of each foreground component is independent of position will have to be relaxed to allow for spectral variations across the sky. Also, it will be necessary to deal with the fact that Galactic emissions are correlated. The technique is flexible enough to offer good prospects in this respect. In the learning stage, the ICA algorithm makes use of non-linear functions that, case by case, are chosen to minimize the mutual information between the outputs; improvements could be obtained by specializing the ICA inner non-linearities to our specific needs. Also, it is possible to take properly into account our prior knowledge on some of the signals to recover, still taking advantage as far as possible of the ability of this neural network approach to carry out a “blind” separation. Work in this direction is in progress. We warmly thank Luigi Danese for original suggestions. We also thank Krzysztof M. Górski and all the people who collaborated to build the HEALPix pixelization scheme extensively used in this work. Work supported in part by ASI and MURST. LT acknowledges financial support from the Spanish DGES, projects ESP98–1545–E and PB98–0531–C02–01.
warning/0002/cond-mat0002261.html	ar5iv	text	# Fluctuations in mixtures of lamellar- and nonlamellar-forming lipids ## I Introduction It is well known that the biological, lipid bilayer membrane contains a mixture of lipids some of which, alone in aqueous solution, do not form bilayers at all. This presents the interesting question as to what possible function in the membrane these lipids can be serving. There are two lines of thought on the answer. The first is that the nonlamellar-forming lipids, roughly characterized as having a small head group and long, splayed tails, can with these tails increase significantly the local lateral pressure distribution in the interior of the bilayer. This increased pressure may be necessary for a protein or peptide to function. The second is that the nonlamellar-forming lipids can make less expensive the energetic cost to form a small, nonlamellar region which may be needed in processes such as membrane fusion. Experiment shows that there is the expected correlation between the presence of nonlamellar-forming lipids and the ability of the lamellar phase to make a transition to a nonlamellar, inverted-hexagonal one. There is also a correlation between the concentration of such lipids and the ability to make other nonlamellar structures, such as occur in membrane fusion. In an attempt to shed some light on the role in membranes of nonlamellar-forming lipids, we first introduced a relatively simple and tractable model of a lipid. We showed that it exhibited the polymorphism observed in experiment as a function of a single architectural parameter, the fraction of the total lipid volume occupied by the head group. We also showed that the phase diagram of a model aqueous lipid whose head group fraction was similar to that of dioleoylphosphatidylethanolamine was in good agreement with the measured diagram. We then considered a mixture of two lipids with the same headgroup but different tail length. Thus they were distinguished by the different relative volume fractions of their headgroups such that one formed lamellar, $`L_\alpha `$, phases while the other formed inverted-hexagonal, $`H_{II}`$, ones. We examined the density profiles of the tail segments from the two different lipids in each of the two phases. In particular we determined that, in the $`H_{II}`$ phase of the mixture, the density of the tail segments of the nonlamellar-forming lipid varied around the Brillouin zone, with its density being largest in the direction between next-nearest-neighbor cores, as expected. We were able to make quantitative this expected variation, and found the relative difference in tail segment densities to be on the order of a few percent. The calculations in Ref. and were carried out within self-consistent field theory, (SCFT), and therefore ignored fluctuations. In this paper, we consider the effect of Gaussian fluctuations about the SCFT solution. This enables us to determine the effect of fluctuations, within Gaussian order, on the transition from the $`L_\alpha `$ to the $`H_{II}`$ phase, and to observe the initial stages of the path between them. In order to highlight the effect of fluctuations, we examine the metastable lamellar phase near its spinodal in a region of the phase diagram in which the $`H_{II}`$ phase is the stable one. We find that undulations occur in the lamellae occupied by the tails, and that the lamellae occupied by the head groups pinch off to form the tubes of the $`H_{II}`$ phase. This is similar to scenarios proposed earlier by Hui et al. and Caffrey and which arose from their experimental observations. The hexagonal phase that is initiated by the lowest energy fluctuation mode is characterized by a lattice parameter which is within 1% of the stable $`H_{II}`$ phase, and has the same orientation: one in which the tubes are coplanar with the lamellae from which they were formed. This coplanarity is in agreement with experiment. We observe that the nonlamellar-forming lipids dominate the undulations in the tail region, while the lamellar-forming lipids tend to fill in the regions between the cylinders as they pinch off. To make quantitative the effect of the nonlamellar-forming lipids, we examine three different measures of their role in the transition; the correlation between the density fluctuation and the order parameter, the ratio of the susceptibilities of the two different lipids to fluctuations at the critical wavevector, and the overlap of the order parameter fluctuations of each lipid with the equilibrium $`H_{II}`$ structure. All these measures show the nonlamellar-forming lipids to play a significant role in bringing about the transition. Having calculated the fluctuations in Gaussian approximation, we are able to calculate all structure factors. We show here the order parameter-order parameter structure factor, which is the one most readily measured, as well as the density-density structure factor of either lipid. In the next section we review the theory of Gaussian fluctuations about an ordered phase and extend it to our model of lipid mixtures. Results are presented next, and we conclude with a brief summary. ## II Theory We consider an anhydrous mixture of $`n_1`$ lipids of type 1 and $`n_2`$ lipids of type 2 in a volume $`V`$. Below we shall choose their architecture so that type 1 lipids form lamellar phases while type 2 lipids form $`H_{II}`$ phases. All lipids consist of the same head group of volume $`v_h`$ and two equal-length tails. Thus we model mixtures of lipids drawn from a homologous series, such as the phosphatidylethanolamines studied by Seddon et al.. The local density of the headgroups of lipid $`L=1,2`$, measured in units of the density $`1/v_h`$, is denoted $`\mathrm{\Phi }_{h,L}(𝐫)`$. Each tail of lipid $`L`$ consists of $`N\alpha _L`$ segments of volume $`v_t`$. The local density of lipid tails of lipid $`L`$, again measured in units of $`1/v_h`$, is denoted $`\mathrm{\Phi }_{t,L}(𝐫)`$. The local volume fraction of these tail segments is $`(2Nv_t/v_h)\mathrm{\Phi }_{t,L}\gamma _t\mathrm{\Phi }_{t,L}`$. The sole architectural parameter, $`f_Lv_h/(v_h+2N\alpha _Lv_t)=1/(1+\alpha _L\gamma _t)`$, which characterizes each lipid is the relative volume fraction of its headgroup. The only interaction in the system is between the headgroups and the tails. The interaction energy $`E`$ is $$\frac{1}{kT}E[\mathrm{\Phi }_{h,1}+\mathrm{\Phi }_{h,2},\mathrm{\Phi }_{t,1}+\mathrm{\Phi }_{t,2}]=\frac{2N\chi }{v_h}[\mathrm{\Phi }_{h,1}(𝐫)+\mathrm{\Phi }_{h,2}(𝐫)][\mathrm{\Phi }_{t,1}(𝐫)+\mathrm{\Phi }_{t,2}(𝐫)]𝑑𝐫,$$ (1) where $`\chi `$ is the strength of the interaction, and $`T`$ is the temperature. The partition function of the system can, without approximation, be written, $$𝒵=\underset{L=1}{\overset{2}{}}𝒟\mathrm{\Phi }_{h,L}𝒟W_{h,L}𝒟\mathrm{\Phi }_{t,L}𝒟W_{t,L}\mathrm{exp}[\mathrm{\Omega }/kT],$$ (2) where $`𝒟`$ denotes a functional integral, and the grand potential $`\mathrm{\Omega }`$ is given by $`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{kT}{v_h}}{\displaystyle \underset{L=1}{\overset{2}{}}}\left\{z_L𝒬_L[W_{h,L},W_{t,L}]+{\displaystyle 𝑑𝐫[W_{h,L}\mathrm{\Phi }_{h,L}+W_{t,L}\mathrm{\Phi }_{t,L}]}\right\}`$ (3) $`+`$ $`E[\mathrm{\Phi }_{h,1}+\mathrm{\Phi }_{h,2},\mathrm{\Phi }_{t,1}+\mathrm{\Phi }_{t,2}],`$ (4) with $`𝒬_L[W_{h,L},W_{t,L}]`$ the partition function of a single lipid of type $`L`$ in the external fields $`W_{h,L}`$ and $`W_{t,L}`$. The integrals over the densities $`\mathrm{\Phi }_{h,L}`$ and $`\mathrm{\Phi }_{t,L}`$ could be carried out as the densities appear only quadratically, but the integrals over the fields $`W_{h,L}`$ and $`W_{t,L}`$ can not be. Thus some approximate evaluation must be made. The SCFT consists in replacing the exact free energy, $`kT\mathrm{ln}𝒵`$ by the extremum of $`\mathrm{\Omega }`$. In addition, we extremize this potential subject to an incompressibility constraint that the sum of all head and tail densities is unity everywhere. We denote the values of the $`\mathrm{\Phi }_{h,L}`$, $`W_{h,L}`$, $`\mathrm{\Phi }_{t,L}`$ and $`W_{t,L}`$ which extremize $`\mathrm{\Omega }`$ subject to this constraint by lower case letters. They are obtained from the following set of self-consistent equations: $`\varphi _{h,L}(𝐫)`$ $`=`$ $`z_L{\displaystyle \frac{\delta 𝒬_L}{\delta w_{h,L}(𝐫)}},L=1,2,`$ (5) $`\varphi _{t,L}(𝐫)`$ $`=`$ $`z_L{\displaystyle \frac{\delta 𝒬_L}{\delta w_{t,L}(𝐫)}},L=1,2,`$ (6) $`w_{h,L}(𝐫)`$ $`=`$ $`2\chi N{\displaystyle \underset{L^{}}{}}\varphi _{t,L^{}}(𝐫)\xi (𝐫),L=1,2,`$ (7) $`w_{t,L}(𝐫)`$ $`=`$ $`2\chi N{\displaystyle \underset{L^{}}{}}\varphi _{h,L^{}}(𝐫)\gamma _t\xi (𝐫),L=1,2,`$ (8) $`1`$ $`=`$ $`{\displaystyle \underset{L}{}}\varphi _{h,L}(𝐫)+\gamma _t{\displaystyle \underset{L}{}}\varphi _{t,L}(𝐫),`$ (9) where $`\xi (𝐫)`$ is the Lagrange multiplier used to enforce the incompressibility constraint. The free energy in this approximation, $`\mathrm{\Omega }_{scf}`$, is obtained from Eq. 3 by replacing the densities and fields by their self-consistent values. To include fluctuations about this self-consistent field result, one decomposes the fields, $`W_{h,L}`$ and $`W_{t,L}`$, and the densities, $`\mathrm{\Phi }_{h,L}`$ and $`\mathrm{\Phi }_{t,L}`$, into their self-consistent field values plus a fluctuating part; $`W_{h,L}=w_{h,L}+\delta W_{h,L}`$, etc. One then substitutes these decompositions into the free energy, Eq. 3. Terms linear in the deviations vanish by definition of the self-consistent field approximation. We consider all quadratic terms in this paper so that the fluctuations are accounted for within Gaussian approximation. Thus the exact partition function of Eq. 2 is approximated by $$𝒵\mathrm{exp}[\mathrm{\Omega }_{scf}/kT]\underset{L=1}{\overset{2}{}}𝒟\delta \mathrm{\Phi }_{h,L}𝒟\delta W_{h,L}𝒟\delta \mathrm{\Phi }_{t,L}𝒟\delta W_{t,L}\mathrm{exp}[\mathrm{\Omega }^{(2)}/kT].$$ (10) The expression for $`\mathrm{\Omega }^{(2)}`$ is most easily written in matrix form. Introduce the column matrices $$\delta W=\left(\begin{array}{c}\delta W_{h,1}\\ \delta W_{t,1}\\ \delta W_{h,2}\\ \delta W_{t,2}\end{array}\right)$$ (11) and $$\delta \mathrm{\Phi }=\left(\begin{array}{c}\delta \mathrm{\Phi }_{h,1}\\ \delta \mathrm{\Phi }_{t,1}\\ \delta \mathrm{\Phi }_{h,2}\\ \delta \mathrm{\Phi }_{t,2}\end{array}\right)$$ (12) and the square matrices $$J=\left(\begin{array}{cccc}0& 1& 0& 1\\ 0& 0& 0& 0\\ 0& 1& 0& 1\\ 0& 0& 0& 0\end{array}\right)$$ (13) $$K=\left(\begin{array}{cccc}z_1C_{hh,1}& z_1C_{ht,1}& 0& 0\\ z_1C_{th,1}& z_1C_{tt,1}& 0& 0\\ 0& 0& z_2C_{hh,2}& z_2C_{ht,2}\\ 0& 0& z_2C_{th,2}& z_2C_{tt,2}\end{array}\right)$$ (14) where $$C_{\alpha \beta ,L}(𝐫,𝐫^{})V\frac{\delta ^2𝒬_L[W_{h,L},W_{t,L}]}{\delta W_{\alpha ,L}(𝐫)\delta W_{\beta ,L}(𝐫^{})}\alpha ,\beta =\mathrm{h}\mathrm{or}\mathrm{t}.$$ (15) In terms of these matrices, the second order correction, $`\mathrm{\Omega }^{(2)}`$, to the self-consistent free energy is $`{\displaystyle \frac{v_h\mathrm{\Omega }^{(2)}}{VkT}}`$ $`=`$ $`{\displaystyle \frac{2\chi N}{V}}{\displaystyle 𝑑𝐫\delta \mathrm{\Phi }^T(𝐫)J\delta \mathrm{\Phi }(𝐫)}`$ (16) $``$ $`{\displaystyle \frac{1}{V}}{\displaystyle 𝑑𝐫\delta W^T(𝐫)\delta \mathrm{\Phi }(𝐫)}{\displaystyle \frac{1}{2V^2}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\delta W^T(𝐫)K(𝐫,𝐫^{})\delta W(𝐫^{})},`$ (17) where the superscript $`T`$ denotes transpose. The Gaussian integrals in Eq. 10 over the four fields $`\delta W_{h,L},\delta W_{t,L}`$, $`L=1,2`$, can now be carried out. This will leave $`\mathrm{\Omega }^{(2)}`$ a functional of the quadratic density fluctuations. These fluctuations are not integrated over in order to leave displayed their coefficients, which are the inverse of the desired density-density correlation functions. It is the inclusion of the field fluctuations to Gaussian order for given density fluctuations which constitutes the random phase approximation. Integration over the field variables in Eq. 10 yields $$𝒵𝒩_0\mathrm{exp}[\mathrm{\Omega }_{scf}/kT]\underset{L=1}{\overset{2}{}}𝒟\delta \mathrm{\Phi }_{h,L}𝒟\delta \mathrm{\Phi }_{t,L}\mathrm{exp}[\mathrm{\Omega }_{rpa}/kT],$$ (18) where the factor $`𝒩_0`$ is the reciprocal of the square root of a determinant which is of no interest here, and $$\mathrm{\Omega }_{rpa}=\frac{kT}{2v_hV}𝑑𝐫𝑑𝐫^{}\delta \mathrm{\Phi }^T(𝐫)\left[4\chi NVJ\delta (𝐫𝐫^{})+K^1(𝐫,𝐫^{})\right]\delta \mathrm{\Phi }(𝐫^{}).$$ (19) In addition to the random phase approximation, we now impose the constraint of incompressibility, which reduces the number of fluctuating fields from four, $`\delta \mathrm{\Phi }_{h,L},\delta \mathrm{\Phi }_{t,L}`$, $`L=1,2`$, to three. We choose these three to be natural order parameters. Define the total local volume fractions of each lipid $$\mathrm{\Phi }_L\mathrm{\Phi }_{h,L}+\gamma _t\mathrm{\Phi }_{t,L},L=1,2,$$ (20) and the local difference between head and tail densities of each lipid $$\mathrm{\Psi }_L\mathrm{\Phi }_{h,L}\frac{1}{\alpha _L}\mathrm{\Phi }_{t,L},L=1,2.$$ (21) Note that as defined, the integral of $`\mathrm{\Psi }_L`$ over the whole system vanishes. In terms of these densities $`\mathrm{\Phi }_{h,L}`$ $`=`$ $`f_L\mathrm{\Phi }_L+(1f_L)\mathrm{\Psi }_L,`$ (22) $`\mathrm{\Phi }_{t,L}`$ $`=`$ $`\alpha _Lf_L[\mathrm{\Phi }_L\mathrm{\Psi }_L].`$ (23) We choose the three independent fluctuating quantities to be $`\delta \mathrm{\Psi }_1`$ and $`\delta \mathrm{\Psi }_2`$, the fluctuations in the difference in head and tail densities of each lipid, and $`\delta \mathrm{\Phi }_1`$, the fluctuation of the total density of lipid 1. Because of the incompressibility constraint, $`\delta \mathrm{\Phi }_1=\delta \mathrm{\Phi }_2`$. The reduction in variables is easily written in terms of the column matrix $$\delta \mathrm{\Theta }=\left(\begin{array}{c}\delta \mathrm{\Phi }_1\\ \delta \mathrm{\Psi }_1\\ \delta \mathrm{\Psi }_2\end{array}\right)$$ (24) and the $`4\times 3`$ matrix $$U=\left(\begin{array}{ccc}f_1& 1f_1& 0\\ \alpha _1f_1& \alpha _1f_1& 0\\ f_2& 0& 1f_2\\ \alpha _2f_2& 0& \alpha _2f_2\end{array}\right)$$ (25) so that $$\delta \mathrm{\Phi }=U\delta \mathrm{\Theta }.$$ (26) Then the correction, $`\mathrm{\Omega }_{rpa}`$ to the self-consistent free energy can be written $`\mathrm{\Omega }_{rpa}`$ $`=`$ $`{\displaystyle \frac{kT}{2v_hV}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\delta \mathrm{\Theta }^T(𝐫)U^T\left[4\chi NJ\delta (𝐫𝐫^{})+K^1(𝐫,𝐫^{})\right]U\delta \mathrm{\Theta }(𝐫^{})},`$ (27) $``$ $`{\displaystyle \frac{kT}{2v_hV}}{\displaystyle 𝑑𝐫𝑑𝐫^{}\delta \mathrm{\Theta }^T(𝐫)K_{rpa}^1(𝐫,𝐫^{})\delta \mathrm{\Theta }(𝐫^{})}.`$ (28) There remains the calculation of the partition function $`𝒬_L[W_{h,L},W_{t,L}]`$ of a single lipid in the external fields $`W_{h,L}`$ and $`W_{t,L}`$. Before carrying this out, we must specify further the way in which the lipid tails are modeled. They are treated as being completely flexible, with radii of gyration $`R_{g,L}=(N_La^2/6)^{1/2}`$ for each tail. The statistical segment length is $`a`$. The configuration of the $`l`$’th lipid of type $`L`$ is described by a space curve $`𝐫_{l,L}(s)`$ where $`s`$ ranges from 0 at one end of one tail, through $`s=\alpha _L/2`$ at which the head is located, to $`s=\alpha _L`$, the end of the other tail. Because the tails are completely flexible, one can define the propagator $`q_L(𝐫,s\|𝐫^{})`$ which gives the relative probability of finding the zero’th segment of a tail of a type $`L`$ lipid at $`𝐫^{}`$ and the segment $`s`$ of the same tail at $`𝐫`$: $$q_L(𝐫,s\|𝐫^{})=_{𝐫_{l,L}(0)=𝐫^{}}^{𝐫_{l,L}(s)=𝐫}\stackrel{~}{𝒟}𝐫_{l,L}(s)\mathrm{exp}\left[_0^s𝑑tW_{t,L}(𝐫_{l,L}(t))\right]0s\alpha _l/2.$$ (29) In this expression, $`\stackrel{~}{𝒟}𝐫_{l,L}(s)`$ denotes a functional integral over the possible configurations of the tail of the lipid of type $`L`$ and in which, in addition to the Boltzmann weight, the path is weighted by the factor $`𝒫[𝐫_{l,L};0,s]`$ with $$𝒫[𝐫,s_1,s_2]=𝒩\mathrm{exp}\left[\frac{1}{8R_g^2}_{s_1}^{s_2}𝑑s\|\frac{d𝐫(s)}{ds}\|^2\right],$$ (30) where $`𝒩`$ is an unimportant normalization constant and $`R_g(Na^2/6)^{1/2}`$ is the radius of gyration of a tail of length $`N`$. This propagator is of central importance, both within the SCFT and the random phase approximation. Because the tails are flexible and execute a random walk in the presence of the potential $`W_{t,L}`$, the propagator can be obtained from the solution of the diffusion equation $$\frac{q_L(𝐫,s\|𝐫^{})}{s}=2R_g^2^2q_L(𝐫,s\|𝐫^{})W_{t,L}(𝐫)q_L(𝐫,s\|𝐫^{}),$$ (31) subject to the initial condition $$q_L(𝐫,0\|𝐫^{})=\delta (𝐫𝐫^{}).$$ (32) The integrated propagator, or endpoint distribution function, $$\stackrel{~}{q}_L(𝐫,s)𝑑𝐫^{}q_L(𝐫,s\|𝐫^{})=𝑑𝐫^{}q_L(𝐫^{},s\|𝐫),$$ (33) is also useful. From the propagator, the partition function of a single lipid of type $`L`$ is obtained via $$𝒬_L[W_{h,L},W_{t,L}]=𝑑𝐫_1𝑑𝐫_2𝑑𝐫_3q_L(𝐫_3,\alpha _L/2\|𝐫_2)\mathrm{exp}[W_{h,L}(𝐫_2)]q_L(𝐫_2,\alpha _L/2\|𝐫_1).$$ (34) Functional derivatives, such as those required in Eqs. 5, 6, and 15, are obtained from the above and $$\frac{\delta q_L(𝐫,s\|𝐫^{})}{\delta W_{t,L}(𝐫_1)}=_0^s𝑑tq_L(𝐫,st\|𝐫_1)q_L(𝐫_1,t\|𝐫^{}).$$ (35) The second functional derivatives of the single lipid partition function are expressed in terms of the propagator; $`C_{hh,L}(𝐫,𝐫^{})`$ $`=`$ $`V\delta (𝐫𝐫^{})\stackrel{~}{q}_L(𝐫,\alpha _L/2)e^{W_{h,L}(𝐫)}\stackrel{~}{q}_L(𝐫,\alpha _L/2),`$ (36) $`C_{ht,L}(𝐫,𝐫^{})`$ $`=`$ $`2Ve^{W_{h,L}(𝐫)}\stackrel{~}{q}_L(𝐫,\alpha _L/2){\displaystyle _0^{\alpha _L/2}}𝑑sq_L(𝐫,\alpha _L/2s\|𝐫^{})\stackrel{~}{q}_L(𝐫^{},s),`$ (37) $`C_{th,L}(𝐫^{},𝐫)`$ $`=`$ $`C_{ht,L}(𝐫,𝐫^{}),`$ (38) $`C_{tt,L}(𝐫,𝐫^{})`$ $`=`$ $`2[F_L(𝐫,𝐫^{})+F_L(𝐫^{},𝐫)+G_L(𝐫,𝐫^{})],`$ (39) where $`F_L(𝐫,𝐫^{})`$ $`=`$ $`V{\displaystyle _0^{\alpha _l/2}}ds{\displaystyle _0^s}dt\stackrel{~}{q}_L(𝐫,\alpha _L/2s)q_L(𝐫,st\|𝐫^{})\times `$ (41) $`{\displaystyle 𝑑𝐫_2q_L(𝐫^{},t\|𝐫_2)e^{W_{h,L}(𝐫_2)}\stackrel{~}{q}_L(𝐫_2,\alpha _L/2)},`$ and $`G_L(𝐫,𝐫^{})`$ $`=`$ $`V{\displaystyle _0^{\alpha _L/2}}ds{\displaystyle _0^{\alpha _L/2}}dt\stackrel{~}{q}_L(𝐫,\alpha _L/2s)\stackrel{~}{q}_L(𝐫^{},\alpha _L/2t)\times `$ (43) $`{\displaystyle 𝑑𝐫_2q_L(𝐫,s\|𝐫_2)e^{W_{h,L}(𝐫_2)}q_L(𝐫_2,t\|𝐫^{})}.`$ Although Eqs. 2943 are exact for arbitrary fields $`W_{h,L}(𝐫)`$ and $`W_{t,L}(𝐫)`$, only the mean-field solutions $`w_{h,L}(𝐫)`$ and $`w_{t,L}(𝐫)`$ are needed in the random phase approximation. Correlation functions in real space can now be calculated. Because there are three independent order parameters, defined in Eqs. 20 and 21, there are six independent correlation functions, $`<\delta \mathrm{\Theta }_i(𝐫)\delta \mathrm{\Theta }_j(𝐫^{})>`$, $`i,j=1,2,3`$, where the brackets indicate an average using the partition function evaluated within the random phase approximation. We will be interested in correlation functions which are various linear combinations of these, so we define a general fluctuation $`\delta \theta _a(𝐫)`$ $``$ $`a_1\delta \mathrm{\Phi }_1(𝐫)+a_2\delta \mathrm{\Psi }_1(𝐫)+a_3\delta \mathrm{\Psi }_2(𝐫),`$ (44) $`=`$ $`(a_1a_2a_3)\left(\begin{array}{c}\delta \mathrm{\Phi }_1(𝐫)\\ \delta \mathrm{\Psi }_1(𝐫)\\ \delta \mathrm{\Psi }_2(𝐫)\end{array}\right)`$ (48) $``$ $`𝐚\delta \mathrm{\Theta }(𝐫),`$ (49) and consider the correlation of two such general fluctuations $`𝐚^T<\delta \mathrm{\Theta }(𝐫)\delta \mathrm{\Theta }^T(𝐫^{})>𝐛`$. When the ensemble average is carried out, one obtains $$𝐚^T<\delta \mathrm{\Theta }(𝐫)\delta \mathrm{\Theta }^T(𝐫^{})>𝐛=\frac{v_h}{V}𝐚^TK_{rpa}(𝐫,𝐫^{})𝐛,$$ (50) where $`K_{rpa}(𝐫,𝐫^{})`$ is the inverse of $`K_{rpa}^1(𝐫,𝐫^{})`$ defined in Eq. 28. Because the phases of interest to us are ordered, it is advantageous to exploit the space-group symmetry and, following Shi et al., to expand all functions of position in terms of Bloch states; $$\psi _{n𝐤}(𝐫)=e^{i𝐤𝐫}\underset{𝐆}{}u_{n𝐤}(𝐆)e^{i𝐆𝐫}.$$ (51) Here $`𝐤`$ is a wavevector in the first Brillouin zone, $`n`$ is the band index, the vectors $`𝐆`$ are the reciprocal lattice vectors of the ordered phase, and the wavefunctions $`u_{n𝐤}(𝐆)`$ satisfy the Schrödinger equation $$2R_g^2(𝐤+𝐆)^2u_{n𝐤}(𝐆)+\underset{𝐆^{}}{}\widehat{W}_t(𝐆𝐆^{})u_{n𝐤}(𝐆^{})=ϵ_{n𝐤}u_{n𝐤}(𝐆).$$ (52) The potentials, $`\widehat{W}_t(𝐆)`$, are just the coefficients of the expansion of the periodic field $`w_t(𝐫)`$, Eq. 8 $$w_{t,1}(𝐫)=w_{t,2}(𝐫)=\underset{𝐆}{}e^{i𝐆𝐫}\widehat{W}_t(𝐆).$$ (53) The expansion of functions of two spatial coordinates $`K(𝐫,𝐫^{})`$, such as the correlation functions, takes the form $$K(𝐫,𝐫^{})=\underset{n𝐤;n^{}𝐤^{}}{}\left[\widehat{K}\right]_{n𝐤,n^{}𝐤^{}}\psi _{n𝐤}(𝐫)\psi _{n^{}𝐤^{}}^{}(𝐫^{})$$ (54) Experimentally measured structure factors are Fourier transforms of the correlation functions $`S_{ab}(𝐪)`$ $``$ $`𝐚^T{\displaystyle \frac{1}{V^2}}{\displaystyle 𝑑𝐫𝑑𝐫^{}e^{i𝐪(𝐫𝐫^{})}<\delta \mathrm{\Theta }(𝐫)\delta \mathrm{\Theta }^T(𝐫^{})>𝐛}`$ (55) $`=`$ $`𝐚^T{\displaystyle \frac{v_h}{V}}{\displaystyle \underset{n,n^{}}{}}\left[\widehat{K}_{rpa}\right]_{n𝐪𝐆;n^{}𝐪𝐆}u_{n𝐪𝐆}(𝐆)u_{n^{}𝐪𝐆}^{}(𝐆)𝐛,`$ (56) where $`𝐪𝐆`$ lies in the first Brillouin zone. All calculations are carried out in the basis of Bloch functions. Therefore the important propagators $`q_L(𝐫,s\|𝐫^{})`$ must be expanded in them, and the procedure outlined above in real-space is followed in the space of Bloch functions. This is straightforward, but tedious. Many of the details can be found in the Appendix of , and we shall not repeat them here. We turn, instead, to our results. ## III Results We have chosen to model two lipids with the same head group, but with different length tails; lipid 1 characterized by $`\alpha _1=1`$ so that its tails are of length N, and lipid 2 characterized by $`\alpha _2=1.5`$ so that its tails are of length 1.5 N. Because the headgroups are identical, the ratio $`\gamma _t2Nv_t/v_h`$ is the same for each lipid, and we have chosen $`\gamma _t=2.5`$. With these parameters, the volume of the head group relative to that of the entire lipid is, for lipid 1, $`f_11/(1+\alpha _1\gamma _t)=0.2857`$, while that of lipid 2 is $`f_21/(1+\alpha _2\gamma _t)=0.2105`$. For comparison, the relative head group volume of dioleoylphosphatidylethanolamine calculated from volumes given in the literature is $`f=0.254`$. From our previous work, we know that anhydrous lipid 1 forms a lamellar phase, while anhydrous lipid 2 forms an inverted hexagonal phase. For orientation, we reproduce in Fig. 1 the phase diagram of the anhydrous system of the two lipids calculated earlier. The temperature, $`T^{}`$, is defined in terms of the interaction strength $`T^{}1/2\chi N_1`$, and $`\mathrm{\Theta }`$ is the volume fraction of lipid 1, the lamellar-forming lipid. Although we know of no phase diagrams for anhydrous mixtures of lamellar- and non lamellar-forming lipids with the same headgroup but different length tails, the system for which we have made our calculation, there are results for mixtures of lamellar-forming phosphatidylcholine and hexagonal-forming phosphatidylethanolamine. Our phase diagram is similar to those observed in anhydrous mixtures of dilinoleoylphosphatidylethanolamine and palmitoyloleoylphosphatidylcholine, and for mixtures of dioleoylphosphatidylcholine and dioleoylphosphatidylethanolamine and 10% water by weight. Each show a significant region of cubic phase between the inverted hexagonal phase, which dominates at low concentrations of the lamellar-forming lipid, and the lamellar phase, which dominates at high concentration. The dashed line in Fig. 1 is the calculated spinodal of the lamellar phase; that is, for volume fractions of the lamellar-forming lipid which are to the left of (i.e. less than) this line, the lamellar phase is absolutely unstable. To the right of this line, but in the regions in which either the inverted hexagonal, $`H_{II}`$ or inverted gyroid, $`G_{II}`$, is stable, the lamellar, $`L_\alpha `$, phase is metastable. At high temperatures, the system is disordered, $`D`$. We have indicated on this phase diagram the two points at which we have calculated structure factors; very near the lamellar spinodal line at $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$, at which point the fluctuations will be large, and at $`T^{}=0.0500`$, $`\mathrm{\Theta }=0.469`$, far from the spinodal, where fluctuations will be much reduced. At the former point, the stable phase is $`H_{II}`$, while at the latter it is the $`G_{II}`$. The structure factor most readily measured experimentally is the Fourier transform of the order parameter-order parameter correlation function $`<\delta \mathrm{\Psi }(𝐫)\delta \mathrm{\Psi }(𝐫^{})>`$, with $`\delta \mathrm{\Psi }(𝐫)\delta \mathrm{\Psi }_1(𝐫)+\delta \mathrm{\Psi }_2(𝐫)`$ the total order parameter. This order parameter is proportional to the difference in the densities of headgroups and tails. The structure factor is obtained by setting the vectors $`𝐚=𝐛`$, defined in Eq. 44, to be $`(0,1,1)`$. It has been calculated at $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$, and is shown in Fig. 2. The normal to the lamellar planes is in the $`z`$ direction. The wave vectors $`k_z`$ and $`k_x`$ are measured in units of $`2\pi /D`$, where $`D`$ is the equilibrium lattice spacing of the lamellar phase. As in the block copolymer system, there is a large response about the positions of the Bragg peaks $`(k_x,k_z)=(0,\pm 1)`$, $`(0,\pm 2)`$, and $`(0,\pm 3)`$. In addition there are four satellite peaks at $`(\pm 0.846,\pm 1/2)`$. As shown below, these fluctuation modes lead to undulations in the planes of head-groups causing them to pinch off, and to form tubes. The array of tubes produced by these fluctuations has a centered rectangular symmetry, c2mm. Its appearance is close to hexagonal, but the structure lacks the $`\pi /3`$ rotational symmetry elements. If a hexagonal structure were formed from the wavevectors $`(0,1)`$ and $`(\sqrt{3}/2,1/2)`$, the ratio of the hexagonal spacing, $`D_h`$, to the lamellar spacing $`D`$ would be $`2/\sqrt{3}1.15`$. We find that the actual ratio of the lattice parameter of the stable, (non-aqueous), hexagonal phase to that of the metastable lamellar phase at the same temperature and composition to be 1.16. Thus the fluctuations do seem to be responsible for driving the lamellar phase to an ordered phase in which cylinders are oriented parallel to the lamellae, and with the correct spacing. The complete hexagonal symmetry, however, cannot be obtained from Gaussian fluctuations alone. We note that because the lamellae are isotropic in the $`xy`$ plane, all structure factors are isotropic in the $`k_xk_y`$ plane. For comparison, the same order parameter-order parameter structure factor, but now calculated at $`T^{}=0.05`$, $`\mathrm{\Theta }=0.469`$, at which point the $`G_{II}`$ phase is stable, is shown in Fig. 3. While the structure about the Bragg peaks remains, there are no satellite peaks here, and therefore no tendency to break the lamellae of headgroups into cylinders. The density-density structure factor, the Fourier transform of $`<\delta \mathrm{\Phi }_1(𝐫)\delta \mathrm{\Phi }_1(𝐫^{})>`$, with $`\delta \mathrm{\Phi }_1=\delta \mathrm{\Phi }_2`$ the fluctuation in the total density of lipid 1, is obtained by setting the vectors $`𝐚=𝐛`$, defined in Eq. 44, to be $`(1,0,0)`$. As our two lipids are drawn from a homologous series and differ only in the length of the tails, this structure factor would not be easily measured. It is shown in Fig. 4 at $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$, and in Fig. 5 at $`T^{}=0.05`$, $`\mathrm{\Theta }=0.469`$. This structure factor is similar to the order parameter-order parameter structure factor. (One discernable difference is the value at zero wavevector which, for the density-density structure factor, is related to the osmotic compressibility. This small peak probably reflects the fact that it is rather easy to replace one lipid by the other because they have identical headgroups, and differ only by the length of their tails.) That the two different structure factors of Figs. 2 and 4 are so similar means that one of the two lipids tends to form cylinders, while the other must form their complement due to the constraint of incompressibility. Which lipid does which can not be determined from this correlation function, but it can from the cross correlation functions, as we now show. Consider the correlation between the density and order parameter fluctuations, $`<\delta \mathrm{\Phi }_1\delta \mathrm{\Psi }>`$. Because the headgroups of all lipids are identical, we expect that changes in the order parameter reflect, for the most part, changes in the densities of the tails. If the non lamellar-forming lipid 2 dominates the order parameter, then we expect that a positive fluctuation in the total tail density should be correlated with a positive change in the density of lipid 2. As the order parameter is defined as the density of heads minus the density of tails, a positive fluctuation in the tail density corresponds to a negative fluctuation in the order parameter. In addition, a positive fluctuation in in the density of lipid 2 corresponds to a negative value of $`\delta \mathrm{\Phi }_1=\delta \mathrm{\Phi }_2`$. As a consequence of the two minus signs, we expect $`<\delta \mathrm{\Phi }_1\delta \mathrm{\Psi }>`$ to be positive if the fluctuation in the non lamellar-forming tails dominate the order parameter. If they do not, it will be negative. A limit on this correlation can be obtained by noting that $$<(\delta \mathrm{\Phi }_1\pm \delta \mathrm{\Psi })^2>0.$$ (57) From this it follows that if $$R_1\frac{2<\delta \mathrm{\Phi }_1\delta \mathrm{\Psi }>}{<(\delta \mathrm{\Phi }_1)^2>+<(\delta \mathrm{\Psi })^2>},$$ (58) then $$1R_11.$$ (59) A value 1 of this ratio means that the density of lipid 2 is completely correlated with the order parameter, a value $`1`$ that the density of lipid 1 is completely correlated with the order parameter, and a zero value that there is no correlation between the density of either lipid and the order parameter. We have evaluated this quantity utilizing only the least stable fluctuation mode, i.e. that mode with non-zero wavevector whose energy will vanish at the spinodal. At $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$, near the lamellar spinodal, this wave vector is again $`k_z=1/2,k_x\sqrt{3}/2`$ . We find $`R_10.27,`$ which indicates that it is the nonlamellar forming lipid whose density is correlated with the order parameter, and that the magnitude of this correlation is appreciable. As a second measure of the importance of the nonlamellar forming lipid in the transition, we have considered the susceptibilities of each lipid to a disturbance with a wavevector equal to that of the least stable mode. At the spinodal, all responses at this wavevector will diverge, in general, but the amplitudes at which they do so will vary. Thus we have evaluated $$R_2\frac{<\delta \widehat{\mathrm{\Psi }}_2(𝐤)\delta \widehat{\mathrm{\Psi }}_2(𝐤)>}{<\delta \widehat{\mathrm{\Psi }}_1(𝐤)\delta \widehat{\mathrm{\Psi }}_1(𝐤)>}$$ (60) at the wavevector of the least stable mode, where $`\delta \widehat{\mathrm{\Psi }}_L(𝐤)`$ is the Fourier component of $`\delta \mathrm{\Psi }_L(𝐫)`$. At the same $`T^{}`$ and $`\mathrm{\Theta }`$ near the lamellar spinodal, we find $`R_2=2.05`$. Thus the susceptibility of the nonlamellar forming lipid to the perturbation at this wavelength is twice that of the lamellar forming lipid. A third measure of the relative importance of the two lipids is obtained by comparing the order parameter fluctuations of the two lipids to the order parameter in the $`H_{II}`$ phase itself. We do this as follows. Consider the first star of wavevectors of the equilibrium hexagonal phase which is the stable one at the temperature and composition at which the lamellar phase is metastable. It consists of six wavevectors, $`𝐤_i`$, $`i=1`$ to 6, of magnitude $`k^{(h)}=4\pi /\sqrt{3}D_h`$, where $`D_h`$ is the lattice parameter of the $`H_{II}`$ phase. The first four can be written as $`(k_x,k_z)=k^{(h)}(\pm \sqrt{3}/2,\pm 1/2)`$, and the other two as $`(k_x,k_z)=k^{(h)}(0,\pm 1)`$. As noted earlier, if $`D_h/D`$ were equal to $`2/\sqrt{3}`$, the first four wavevectors would lie on the edge of the Brillouin zone of the lamellar phase, while the others would correspond to reciprocal lattice vectors of the lamellar phase. We can determine the Fourier components with these six wavevectors of the order parameter fluctuations $`\delta \mathrm{\Psi }_2`$ and $`\delta \mathrm{\Psi }_1`$. This tells us how large is the overlap of these fluctuations with the equilibrium hexagonal structure. We denote this Fourier component $`\delta \widehat{\mathrm{\Psi }}_L(𝐤_i)`$, $`L=1,2.`$ In general, it is complex so we calculate $`\delta \widehat{\mathrm{\Psi }}_L(𝐤_i)+\delta \widehat{\mathrm{\Psi }}_L(𝐤_i)`$, which is real. At $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$, we find $`R_3{\displaystyle \frac{\delta \widehat{\mathrm{\Psi }}_2(𝐤_i)+\delta \widehat{\mathrm{\Psi }}_2(𝐤_i)}{\delta \widehat{\mathrm{\Psi }}_1(𝐤_i)+\delta \widehat{\mathrm{\Psi }}_1(𝐤_i)}}`$ $`=`$ $`1.43i=1\mathrm{to}4`$ (61) $`=`$ $`1.47i=5,6.`$ (62) If the fluctuations had hexagonal symmetry, the two values would be equal. We see that the amplitudes of the hexagonal fluctuations of the non lamellar-forming lipid are almost 50% larger than those of the lamellar-forming lipid. We next display the way in which the fluctuations of the lamellar phase near its spinodal indicate the beginning of the path to the stable $`H_{II}`$ phase. We examine the effect of the fluctuation mode with lowest energy on the order parameter. In the absence of fluctuations, the order parameter, which is the difference in the total head and tail segment density, is $`\varphi _h^{(1)}{\displaystyle \frac{\varphi _t^{(1)}}{\alpha _1}}+\varphi _h^{(2)}{\displaystyle \frac{\varphi _t^{(2)}}{\alpha _2}}`$ We add the contribution to this order parameter from the real part of the fluctuation mode with the lowest energy, $`Re\{\delta \mathrm{\Psi }(𝐫)\}`$. We add it with a variable amplitude, $`ϵ`$ to better visualize its effect. Fig. 6(a) shows the order parameter in mean-field theory, with no fluctuation contribution, $`ϵ=0.`$ The point in the phase diagram is $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$. The center of the headgroup region is at $`z=0`$. The lamellae are in the $`xy`$ plane, and the coordinates are measured in units of the lattice spacing, $`D`$. The grey scale divides the values from $`0.6`$ to $`0.6`$ into ten shades. The darkest values are most positive, and correspond to a dominance of head groups. Values in (a) range from $`0.214`$ to $`0.251`$. In (b), (c) , and (d), the fluctuation contribution is turned on with amplitudes $`ϵ=0.06`$, $`0.12`$, and $`0.18`$ respectively. The extreme values in (d) are $`0.301`$ and $`0.572`$. One clearly sees the undulation of the tail region, and the pinching off of the head group region, leading to an $`H_{II}`$ phase in which the location of the tubes is coplanar with the lamellar phase. Again, the lattice spacing of the stable $`H_{II}`$ is within 1% of the spacing expected if there were a perfect epitaxy between the phases. The mechanism we see indicated here by the fluctuations from the lamellar phase is very similar to those proposed by Hui et al. and Caffrey. The coplanarity which is a direct result of it has been observed by Gruner et al. with oriented photoreceptor membranes. We also note the possible significance for this problem of an observation made in Ref. . Due to isotropy in the $`xy`$ plane, the fluctuation modes are infinitely degenerate in the $`k_xk_y`$ plane, that is, all fluctuation modes with the same magnitude $`k_x^2+k_y^2`$ have the same energy of excitation. If one particular direction is preferentially excited, the mode leads to the formation of ordered cylinders as we have seen above. But if several directions are excited simulataneously, then the ripples in the corresponding directions will interfere with one another. One pattern that can be obtained is a periodic array of holes, in the lamellae, producing a perforated lameller phase. These perforations are similar to pores in the membrane, and to stalks. Finally, we display the effect of the fluctuation mode with the lowest energy on the difference in the tail distribution of the two lipids. The tail distribution of lipid $`L`$ as calculated in the self-consistent field theory, without fluctuations, is $`\varphi _{t,L}(𝐫)`$. We add the contribution to this distribution from the fluctuation mode with the lowest energy, $`ϵRe\{\delta \mathrm{\Phi }_{t,L}(𝐫)\}`$ where $`ϵ`$ is, again, an amplitude which we can adjust. In order to compare the tail distributions of the two lipids, we must take account of the fact that, because the concentrations of the two lipids are unequal, they have different average tail segment densities $$\overline{\varphi }_{t,L}\frac{1}{V}𝑑𝐫\varphi _{t,L}(𝐫).$$ (63) Therefore we utilize the quantity $$P_{t,L}(ϵ,𝐫)\frac{\varphi _{t,L}(𝐫)+ϵRe\{\delta \mathrm{\Phi }_{t,L}(𝐫)\}}{\overline{\varphi }_{t,L}}$$ (64) and plot $$\mathrm{\Delta }P_t(ϵ,𝐫)P_{t,1}(ϵ,𝐫)P_{t,2}(ϵ,𝐫)$$ (65) for various values of $`ϵ`$. This shows us how the fluctuations about the lamellar phase tend to spatially arrange the two different lipids. Again we place ourselves at the point in the phase diagram at which $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$. Figure 7(a) shows this difference within the self-consistent field theory alone, that is with $`ϵ=0`$. The center of the headgroup region is at $`z=0`$. The grey scale is such that darkest values correspond to positive values of $`P_t`$ at which lipid 1 dominates, and lightest values correspond to negative values of $`P_t`$ at which the non lamellar-forming lipid 2 dominates. Thus the density of lipid 2 is largest at the center of the tail region, which is expected as lipid 2 has the longer tails. The maximum and minimum values on the grey scale of Fig. 7 correspond to $`0.16`$ and $`0.16`$, divided into ten shades of grey. The maximum and minimum values attained by $`\mathrm{\Delta }P_t(0,𝐫)`$ are $`0.089`$ and $`0.109`$. Thus the differences in local tail segment densities shown in Fig. 6 are on the order of $`10\%`$ of the average tail segment densities. In figs. 7(b), (c), and (d) we turn on the contribution to the density difference of the tail segments from the lowest energy fluctuation mode, and plot $`\mathrm{\Delta }P_t(0.06,𝐫)`$, $`\mathrm{\Delta }P_t(0.12,𝐫)`$, and $`\mathrm{\Delta }P_t(0.18,𝐫)`$. The largest and smallest values attained by $`\mathrm{\Delta }P_t(0.18,𝐫)`$ are 0.160 and $`0.147.`$ In this sequence, one clearly sees that in the fluctuations, the non lamellar-forming lipid dominates the undulations of the tail region, particularly in the direction of the second neighbors, the direction which is most difficult for the tails to fill. The tails of the lamellar-forming lipid are initially relegated to filling in the region between the cylinders which form within a given lamellae. In sum, we have employed a model system of a non-aqueous mixture of lamellar- and nonlamellar-forming lipids to examine the Gaussian fluctuations of the lamellar phase near its spinodal. We have calculated the structure factor which would be experimentally measurable were the lameller phase sufficiently oriented. We have found that the initial stage of the $`L_\alpha `$ to $`H_{II}`$ transition involves the formation of undulations in the tail-rich lamellae, undulations which are dominated by the nonlamellar-forming lipid. As we have noted, the lattice constant of the ordered, stable, $`H_{II}`$ phase is within 1% of that which would be produced by the lowest energy fluctuation mode, and the orientations are, of course, identical. There is also a coplanarity between one of the principal direction of the tubes in the $`H_\alpha `$ and the orientation of the $`L_\alpha `$ phase which produced it, a coplanarity observed in experiment. We have used several measures to show that the nonlamellar-forming lipid plays a dominant role in this transition in the mixture. Because we have examined the fluctuations very near the spinodal, we are presumably seeing the beginning of spinodal decomposition. Further from the spinodal, and in particular, near the $`L_\alpha `$ phase boundary which is generally far from the spinodal (see Fig. 1), it is to be expected that the transition proceeds by a nucleation and growth mechanism, which may well involve some form of lipid intermediate structure. Presumably by making a series of controlled quenches from the lamellar to the hexagonal phase, one should observe the progression from the spinodal decomposition mechanism of undulations to the nucleation and growth intermediates. Certainly both mechanisms seem to have been observed in the same system. In either form of the transition, structures which are nonlamellar must be created. As we have shown and made quantitative, the nonlamellar forming lipids tend to dominate the process of this kind which we have examined. This lends support to the line of argument that their role in the biological membrane is to facilitate the creation of such structures. This work was supported in part by the National Science Foundation under grant number DMR9876864. Phase diagram of an anhydrous mixture of lipids 1 and 2 in the plane of temperature, $`T^{}`$, and volume fraction, $`\mathrm{\Theta }`$, of lipid 1, the lamellar-forming lipid. Phases shown are inverted hexagonal, $`H_{II}`$, inverted gyroid, $`G_{II}`$, lamellar, $`L_\alpha `$, and disordered, $`D`$. The dashed line is the spinodal of the $`L_\alpha `$ phase. Order-parameter order-parameter structure factor shown in the $`k_zk_x`$ plane. The wavevectors are in units of $`2\pi /D`$, with $`D`$ the lamellar spacing. The point in the phase diagram, $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$, at which the factor is calculated, is near the spinodal of the lamellar phase. Same structure factor as in Fig. 2, but at the point in the phase diagram $`T^{}=0.05`$, $`\mathrm{\Theta }=0.469`$, far from the spinodal of the lamellar phase. Density density structure factor calculated at $`T^{}=0.0625`$, $`\mathrm{\Theta }=0.356`$. Density density structure factor calculated at $`T^{}=0.05`$, $`\mathrm{\Theta }=0.469`$. The total order parameter is shown in the $`xz`$ plane with different amplitudes, $`ϵ`$, of the contribution from the fluctuation mode with the lowest energy; (a) $`ϵ=0.`$, (b) $`ϵ=0.06,`$ (c) $`ϵ=0.12`$, (d) $`ϵ=0.18`$. Light areas indicate regions in which tail density is greatest, dark regions where head density is largest. The difference in the tail segment densities arising from lipid 1 and from lipid 2. The amplitude of the contribution from the fluctuation mode with the lowest energy is (a) $`ϵ=0.`$, (b) $`ϵ=0.06,`$ (c) $`ϵ=0.12`$, (d) $`ϵ=0.18`$. Light areas indicate regions in which tails of lipid 2, the nonlamellar former, dominate, dark regions the predominance of segments from the tails of lipid 1.
warning/0002/gr-qc0002095.html	ar5iv	text	# Spacetime as a Feynman diagram: the connection formulation ## 1 Introduction The spin foam formalism has emerged in the last few years as an elegant synthesis of several approaches to quantum gravity and diffeomorphism invariant theories more generally. It can be viewed as a “path integral” formulation corresponding to the canonical loop quantization framework and also as an extension of the simplicial framework for topological field theories , which allows more general, non-topological, field theories to be represented. Spin foams are coloured two dimensional complexes consisting of two dimensional faces (of arbitrary topology), joined on edges of valence $`3`$. Faces are coloured with non-trivial irreducible representations of a “gauge group” $`G`$<sup>1</sup><sup>1</sup>1 We will see that $`G`$ corresponds to the local gauge group or fibre bundle structure group commonly refered to as the “gauge group” in discussions of Yang-Mills theory, as opposed to the group of all gauge transformations of the fields in spacetime. while edges carry “intertwiners” - $`G`$ invariant tensors in the product representation formed by the representations on the incident faces. In a spin foam model the spin foams are regarded as histories of the physical system and assigned quantum amplitudes. In models of gravity a spin foam defines a discrete spacetime geometry ; Spin foam models of gravity incorporate the discreteness of the geometry of space first uncovered in loop quantized canonical theory into a spacetime sum over histories formalism naturally suited to a 4-diffeo invariant theory. This is very appealing because the discreteness, which is not “put in by hand” but arises naturally in this framework,<sup>2</sup><sup>2</sup>2 In using the spin foam formalism, or the canonical loop quantization formalism one is making an ansatz, so it is possible that the discreteness is not real because the whole framework is not the one used by Nature. In lattice spin foam models there is another possibility, namely that an infinite renormalization of the bare theory that is necessary to go to the continuum limit wipes out the discreteness of the geometry in this limit. Finally it should be noted that in a spin foam model of Lorentzian 2+1 general relativity developed by Freidel geometrical observables do not have an entirely discrete spectrum but there is a non-zero minimal length. This is a result of the non-compactness of the three dimensional Lorentz group $`SO(2,1)`$ and might be a general feature of Lorentzian models. promises to remove the ultraviolet divergences found in perturbative theories of quantum gravity and matter fields. A number of spin foam models of Euclidean<sup>3</sup><sup>3</sup>3 These models are not Euclidean field theories in the ordinary sense because essentially the exponential of i times the action is used to weight histories. We call these models of Euclidean quantum gravity because they are proposals for quantizations of classical general relativity with metric signature $`++++`$. and Lorentzian quantum gravity in four spacetime dimensions have been proposed. Three dimensional general relativity as well as four dimensional BF theory can also be given a spin foam formulation. Even hypercubic Yang-Mills theory can be expressed in this framework. In fact any lattice model in which a connection with a compact gauge group forms the boundary data can be translated into a spin foam model . Of the spin foam models of gravity referred to above all but and are simplicial lattice models. A shortcoming of such models is that their predictions depend on the simplicial complex chosen to represent spacetime. Topological field theories, such as three dimensional general relativity or BF theory, can be formulated on a triangulated manifold in such a way as to be independent of the particular triangulation of the manifold chosen, essentially because such theories have only a finite number of degrees of freedom associated to global topological features of the manifold. Four dimensional gravity models on the other hand are expected to have infinitely many degrees of freedom in any given spacetime topology, so any finite triangulation (in which each simplex carries a finite set of degrees of freedom as in the proposed models) necessarily truncates and thus misrepresents the gravitational field. Some sort of continuum limit must be taken. One approach is to sum over triangulations of the spacetime manifold. This seems very difficult to do in three or more dimensions since it is difficult to identify the topology of a given simplicial complex. A technically easier, but also more ambitous approach would be to sum over all simplicial manifolds<sup>4</sup><sup>4</sup>4 Simplicial manifolds are simplicial complexes that are also manifolds. To be a manifold a complex must satisfy several local requirements the most non-trivial of which is that the boundary of the union of the simplicies incident on a point is a sphere. and thus sum over both triangulations and spacetime topology. It is an old dream of gravity theorists to make topology dynamical, thus removing yet another element of a priori background structure. It would be interesting to sum even models formulated on a continous spacetime manifold, such as those found in and , over spacetime topologies. In the present paper we will consider another, even wider, summation. Simplicial spin foams live on the 2-skeleton of the cellular complex dual to the simplicial complex. In other words, the spin foam faces consist of dual 2-cells, and the spin foam edges of dual 1-cells. The dual 2-skeleton, the two dimensional complex formed by these cells (together with the dual 0-cells), thus is spacetime as far as spin foams are concerned. We therefore propose to sum spin foam models of four spacetime dimensional gravity over all “admissible” 2-complexes - 2-complexes that preserve some simple local features of dual 2-skeletons of triangulated 4-manifolds. Specifically we will require of admissible 2-complexes that their 0-, 1-, and 2- cells are topologically points, line segments, and disks respectively. Furthermore we require that the combinatorial structure at each 0-cell is that of a dual 2-skeleton to a triangulated 4-manifold. Thus each 0-cell has incident on it five endpoints of 1-cells and ten corners of 2-cells with each incident 2-cell corner bounded by two of the incident 1-cells. This corresponds to the fact that a 4-simplex (dual to a 0-cell) has in its boundary five 3-simplices (dual to 1-cells) and ten 2-simplices (dual to 2-cells). (See Fig. 1). In a sufficiently small neighborhood of a 0-cell the 2-complex is indistinguishable from a dual 2-skeleton of a 4-d simplicial complex. Note however that we do not forbid 1-cells and 2-cells from being incident several times on the same 0-cell, so some fairly strange 2-complexes are included in the sum. Our main result is that this summation of a spin foam model over 2-complex spacetimes may be realized as the sum over Feynman diagrams in a perturbation expansion of a quantum field theory: The spin foam model defines a (somewhat unusual) quantum field theory on the Cartesian product of four copies of the gauge group $`G`$ the Feynman diagrams of which are precisely the admissible 2-complexes, with the amplitude of each given, modulo symmetry factors, by the spin foam sum on that 2-complex. (Comparing with a familiar QFT such as scalar $`\lambda \varphi ^4`$ theory in Minkowski space we see that the spacetime 2-complex plays the role of the graph of a Feynman diagram, while the spin foams living on the the 2-complex play the roles of the values of the momenta on the edges of a Feynman diagram, which are integrated over to obtain the amplitude of the diagram). We propose to adopt the perturbation series of the field theory, with its symmetry factors, as the definition of the spin foam model summed over 2-complex spacetimes. To make things more concrete let’s consider the amplitude for given boundary data on the boundary of a spacetime region. This amplitude is the sum of the amplitudes of the histories that match the boundary data.<sup>5</sup><sup>5</sup>5 This amplitude can be thought of as the quantum probability amplitude of the boundary data in the Hartle-Hawking vacuum. Alternatively, if the boundary is divided into past and future parts then the amplitude can be interpreted as the transition amplitude from past to future data. In a spin foam model on a fixed 2-complex $`J`$ representing spacetime the boundary is a graph $`\mathrm{\Gamma }=J`$ and the sum runs over all spin foams living on $`J`$ matching the boundary data - irreducible representations on the edges of $`\mathrm{\Gamma }`$ and intertwiners on the nodes. Now we wish to sum the amplitude obtained by summing over spin foams also over all $`J`$ bounded by the same, fixed, $`\mathrm{\Gamma }`$,<sup>6</sup><sup>6</sup>6 In the sum we are contemplating $`\mathrm{\Gamma }`$ need not be the boundary of $`J`$ in quite the usual sense. We allow the possibility that parts of $`\mathrm{\Gamma }`$ are glued together instead of to $`J`$. The situation is similar to that of encountered with shrink wrapped foods in which part of the plastic covering is stuck to itself rather than the food item. Returning to a more theoretical setting, this peculiar feature of the boundaries we admit is physically quite natural. If we divide $`\mathrm{\Gamma }`$ into a future half $`\mathrm{\Gamma }_+`$ and a past half $`\mathrm{\Gamma }_{}`$ and interpret the amplitude of the boundary data as the transition amplitude from the past to the future boundary data, then there is no reason to exclude complexes representing “bubble” time evolution in which parts of $`\mathrm{\Gamma }_+`$ and $`\mathrm{\Gamma }_{}`$ coincide, and no reason to exclude such complexes from a sum over complexes interpolating between $`\mathrm{\Gamma }_+`$ and $`\mathrm{\Gamma }_{}`$. which is to be considered part of the boundary data. In the field theory picture this sum over 2-complexes is the Feynman diagram expansion of the expectation value of an observable that encodes the boundary data. The field theory formulation provides us with formal functional integral expressions for this expectation value, and any other quantitiy of interest, which are often easier to manipulate than the original sums. Moreover, a regularization scheme that renders well defined the functional integral would also provide a definition of the values of the sums. We shall return to this point. Our work generalizes the results of Ooguri and De Pietri et. al. who obtained field theory formulations for two particular models (summed over spacetime 2-complexes), namely $`SU(2)`$ BF theory and the Barrett-Crane model of Euclidean quantum gravity. Many of the underlying ideas go back to the work of Boulatov , on three dimensional gravity, and to matrix models of two dimensional theories . The remainder of the paper is organized as follows. In §2 we define precisely the class of spin foam models under consideration and then we present and verify our main result: an explicit formula for the action of the field theory that defines the sum over 2-complexes of any given spin foam model in this class. The Turaev-Ooguri-Crane-Yetter model of BF theory is discussed as an example. Throughout this section spin foam models are formulated as lattice gauge theories on a special type of lattice. In §3 we show how these spin foam models, and the field theory generating the sum over spacetime 2-complexes, are represented in terms spin foams. The issue of divergences, regularization and renormalization is also briefly touched on. In the last section we give an argument suggesting that the field theory is actually finite for regulated spin foam models (which is not at all obvious from the perspective of a sum over spacetime 2-complexes) and outline how the theory can be extended to accomodate more general spin foam models. ## 2 Field theory formulation of spin foam models with dynamical topology ### 2.1 Local spin foam models in a lattice connection formulation In the present section we obtain, explicitly, a field theory defining the sum over 2-complex spacetimes for a wide class of “local” spin foam models living on the dual 2-skeletons of triangulated 4-manifolds, or more generally on admissible 2-complexes. This class includes all Euclidean simplicial four dimensional gravity models except that of Iwasaki . Instead of working directly with spin foam sums we will use the connection formulation of the spin foam models because this provides the easiest route to our result. In §3 we will show how the result looks in the spin foam formulation and indicate how it may be obtianed within that framework. In the connection formulation a history consists of a connection specified by elements of $`G`$ defining parallel transport along “boundary edges” which run from the center of a 2-cell of $`J`$ to the center of one of its bounding 1-cells.<sup>7</sup><sup>7</sup>7 Where the centers of these cells are placed does not really matter. In fact the parallel transporters could simply be associated with pairs consisting of a 2-cell and a 1-cell in its boundary. As illustrated in Fig. 2 a) these edges cut $`J`$ into disjoint, topologically identical pieces which we shall call “atoms”, because they can be viewed as the fundamental building blocks from which any admissible 2-complex can be built. Each atom contains one 0-cell, five one dimensional cells called “spokes” (portions of 1-cells of $`J`$), and ten two dimensional cells called “wedges” (portions of 2-cells of $`J`$). It is bounded by a graph made of boundary edges that is homeomorphic to the 1-skeleton $`\mathrm{\Gamma }_5`$ of a 4-simplex, having five 4-valent nodes connected in all possible ways by chains consisting of two boundary edges. The 4-valent nodes are the centers of incident 1-cells while the 2-valent nodes in the chains are the centers of incident 2-cells. (Henceforth when we speak of the nodes of the atomic boundaries we shall mean only the 4-valent nodes unless the 2-valent nodes are explicitly included). When $`J`$ is the dual 2-skeleton of a four dimensional simplicial complex the boundary edges are precisely the edges where the boundaries of the 4-simplices cut the dual 2-skeleton, and the atoms that they cut $`J`$ into are the portions of $`J`$ inside each 4-simplex. Atoms are in this sense the dual 2-skeletons of 4-simplices. See Fig. 2 b). The models we shall consider have compact $`G`$ and they are local. Each is defined by a “vertex function” $`V`$, a gauge invariant function of a connection on $`\mathrm{\Gamma }_5`$. This function evaluated on the connection on the boundary of an atom gives the quantum mechanical amplitude for that connection. The amplitude for the whole connection on all of $`J`$ is then the product of the amplitudes for all of the atoms. If we number the 4-valent nodes of $`\mathrm{\Gamma }_5`$ from 1 to 5 and let the indices $`i,j,k,\mathrm{}`$ range over these numbers then we may indicate the oriented edge of $`\mathrm{\Gamma }_5`$ from node $`i`$ to node $`j`$ by $`l_{ij}`$ and the half edge from $`i`$ to the center of $`l_{ij}`$ (which corresponds to a boundary edge in the boundary of an atom) by $`e_i^j`$. $`V`$ is then a function of the parallel transporters $`g_{ij}=h_i{}_{}{}^{j}[h_j{}_{}{}^{i}]_{}^{1}`$ along the $`l_{ij}`$, where $`h_i{}_{}{}^{j}G`$ is the parallel transporter along $`e_i^j`$. In the model the amplitude for one atom, $`x`$, is thus $`V(h_{xi}{}_{}{}^{j}[h_{xj}{}_{}{}^{i}]_{}^{1})`$ where $`h_{xi}^j`$ is the parallel transporter along the boundary edge mapped to $`e_i^j`$, and the amplitude for the connection on the whole spacetime is $$w=\underset{x\text{atoms of}J}{}V(h_{xi}{}_{}{}^{j}[h_{xj}{}_{}{}^{i}]_{}^{1}).$$ (1) (Since each atom contains one 0-cell the product can also be viewed as a product over 0-cells of $`J`$).<sup>8</sup><sup>8</sup>8 A footnote on vector and matrix notation: In (1) $`h_{xi}^j`$ indicates the whole matrix of group elements $`[h_{xi}^j]_{1i<j5}`$ in a sort of abstract index notation. It should be clear from the context when such expressions denote a whole matrix or vector or just one element. On other occasions we will use boldface letters to denote matrices or vectors, so the matrix of $`h_{xi}^j`$s would be written as $`𝐡_x`$ and the second row of this matrix can be written as $`𝐡_{x\mathrm{\hspace{0.17em}2}}^.`$. Again information that is clear from the context will be left out, with $`𝐡`$ in one context refering to all $`h`$s in the whole 2-complex while in another only to the four $`h`$s associated with a given node on an atomic boundary. Finally we shall use the Einstein summation convention on gauge group tensor indices, so any repeated index of this type is summed over. If the model is a simplicial approximation to a spacetime manifold field theory, as all the gravity models aim to be, then $`V`$ is an approximation to $`\mathrm{exp}(i[\text{Effective action in 4-simplex}])`$ in terms of the connection on the boundary of the atom in the 4-simplex. The atomic boundary is a graph in the boundary of the simplex so the connection on it is a discrete approximation to the continuum connection on the boundary of the 4-simplex. As is proper for a local theory regions of spacetime (composed of whole atoms) communicate only via boundary data - the connection on the boundary. A model in this connection formalism is nothing other than a lattice gauge theory defined on an unusual lattice in order to bring boundary data and locality to the fore. For this reason this class of models was called “local lattice gauge theory” in . ### 2.2 The main result: the field theory that generates the sum over topologies We can now introduce the field theory that generates the sum over admissible 2-complexes of the model defined by the vertex function $`V`$. It is a real scalar field theory on $`G^4`$, the Cartesian product of four copies of the gauge group manifold, determined by the action $$I[\psi ]=I_0[\psi ]\lambda 𝒱[\psi ],$$ (2) where $$I_0[\psi ]=\frac{1}{24!}_{G^4}d^4h\psi ^2(h_1,h_2,h_3,h_4)$$ (3) and $$𝒱[\psi ]=\frac{1}{5!}_{G^{20}}d^{20}hV(h_i{}_{}{}^{j}[h_j{}_{}{}^{i}]_{}^{1})\psi (h_1{}_{}{}^{i})\psi (h_2{}_{}{}^{i})\psi (h_3{}_{}{}^{i})\psi (h_4{}_{}{}^{i})\psi (h_5{}_{}{}^{i}).$$ (4) $`I_0`$ is the kinetic term. It is quadratic, but it contains no derivatives. $`\lambda 𝒱`$ is a 5th order, non-local interaction term, with $`\lambda `$ the coupling constant in which we will expand to get the perturbation series. The scalar field $`\psi (h_1,h_2,h_3,h_4)`$ is required to be symmetric in its four arguments. A wavefunction $`\theta `$ of the connection on boundary $`\mathrm{\Gamma }`$ is represented by the observable<sup>9</sup><sup>9</sup>9 The observable, mentioned in the introduction, that represents a particular boundary connection $`𝐡_\mathrm{\Gamma }`$ is obtained by taking $`\theta `$ to be a delta distribution on the gauge equivalence class of $`𝐡_\mathrm{\Gamma }`$. $$\mathrm{\Theta }[\psi ]=\frac{1}{sym(\mathrm{\Gamma })}d𝐡_\mathrm{\Gamma }\theta (𝐡_\mathrm{\Gamma })\underset{a\text{nodes of}\mathrm{\Gamma }}{}\psi (h_a{}_{}{}^{b}).$$ (5) The arguments $`h_a^b`$ of $`\psi `$ are the parallel transporters along the boundary edges $`e_a^b`$ from the node $`a`$ toward the four neighboring nodes (indexed by $`b`$); $`𝐡_\mathrm{\Gamma }`$ represents the whole connection on $`\mathrm{\Gamma }`$, i.e. the vector of all the parallel transporters along boundary edges in $`\mathrm{\Gamma }`$; and $`sym(\mathrm{\Gamma })`$ is the number of symmetries of $`\mathrm{\Gamma }`$, i.e. the number of mappings of the set of nodes of $`\mathrm{\Gamma }`$ to itself that preserve the adjacency matrix. (Since the identity mapping is a symmetry $`sym(\mathrm{\Gamma })1`$). The main result of the present paper is the following: Theorem The formal perturbation series of the expectation value of $`\mathrm{\Theta }^{}`$, $$𝒟\psi e^{I_0[\psi ]+\lambda 𝒱[\psi ]}\mathrm{\Theta }^{}[\psi ],$$ (6) in powers of $`\lambda `$ is $$\underset{J𝐀_\mathrm{\Gamma }}{}\frac{\lambda ^{n(J)}}{sym(J)}𝑑𝐡\theta ^{}(𝐡_\mathrm{\Gamma })w(𝐡),$$ (7) the sum over the set $`𝐀_\mathrm{\Gamma }`$ of admissible 2-complexes $`J`$ bounded by $`\mathrm{\Gamma }`$ of the overlap of the state $`\theta `$ with the Hartle-Hawking vacuum of the spin foam model on $`J`$. $`n(J)`$ is the number 0-cells in $`J`$ and $`sym(J)`$ is the number of symmetries of $`J`$, i.e. of mappings of $`J`$ to itself that preserve the combinatorial structure of $`J`$.<sup>10</sup><sup>10</sup>10 The identity mapping is included in this set so $`sym(J)`$ is at least 1. $`𝒟\psi `$ is normalized so that $`𝒟\psi e^{I_0[\psi ]}=1`$. Proof: The proof is similar to the derivation the Feynman diagram expansion in local field theories in Minkowski space. The order $`\lambda ^n`$ term in (6) is $$\frac{\lambda ^n}{n!}𝒟\psi e^{I_0[\psi ]}𝒱^n[\psi ]\mathrm{\Theta }^{}[\psi ].$$ (8) The functional integral is Gaussian. It’s value is zero if the number of factors of $`\psi `$ in $`𝒱^n\mathrm{\Theta }^{}`$ is odd. If the number is even the value of the integral is a sum of terms associated with each possible partition of the factors of $`\psi `$ in $`𝒱^n\mathrm{\Theta }^{}`$ into pairs. Here we imagine that $`𝒱^n\mathrm{\Theta }^{}`$ has been written as $`n`$ factors of $`𝒱`$ followed by $`\mathrm{\Theta }^{}`$ and that the $`𝒱`$s and $`\mathrm{\Theta }^{}`$ have been written explicitly as the integrals (4) and (5) with the factors of $`\psi `$ in the integrals appearing in some definite order. The pairings being summed over are the distinct pairings of the elements in the ordered sequence of $`\psi `$s appearing in this explicit expression for $`𝒱^n\mathrm{\Theta }^{}`$. The term corresponding to a given pairing is obtained by replacing each pair $`\psi (𝐡)\psi (𝐡^{})`$ by the corresponding propagator $`\psi (𝐡)\psi (𝐡^{})_0`$ $``$ $`{\displaystyle 𝒟\psi e^{I_0[\psi ]}\psi (𝐡)\psi (𝐡^{})}`$ (9) $`=`$ $`{\displaystyle 𝒟\psi \psi (𝐡)\underset{\sigma S_4}{}\frac{\delta }{\delta \psi (\sigma [𝐡^{}])}e^{I_0[\psi ]}}`$ (10) $`=`$ $`{\displaystyle \underset{\sigma S_4}{}}{\displaystyle 𝒟\psi \frac{\delta \psi (𝐡)}{\delta \psi (\sigma [𝐡^{}])}e^{I_0[\psi ]}}`$ (11) $`=`$ $`{\displaystyle \underset{\sigma S_4}{}}\delta (h_1,h_{\sigma (1)}^{})\delta (h_2,h_{\sigma (2)}^{})\delta (h_3,h_{\sigma (3)}^{})\delta (h_4,h_{\sigma (4)}^{}).`$ (12) Here $`𝐡`$ and $`𝐡^{}`$ are the sequences of four group element arguments of $`\psi `$, and $`\sigma [𝐡^{}]=[h_{\sigma (1)}^{},\mathrm{},h_{\sigma (4)}^{}]`$ is $`𝐡^{}`$ reordered according to the permutation $`\sigma `$. Note that the symmetrization had to be introduced in (10) because the integral runs over symmetrized $`\psi `$ only. Thus the integration by parts is justified only for the symmetrized derivative.<sup>11</sup><sup>11</sup>11 Note also that the delta distributions in (12) are normalized Haar measure deltas. The definition (3) of $`I_0`$ is a Haar measure integral and the functional derivatives in (10) are Haar measure functional derivatives. $`𝒱^n\mathrm{\Theta }^{}`$ is an integral over the connections on the boundaries of $`n`$ separate atoms and on $`\mathrm{\Gamma }`$. When each of the pairs of $`\psi `$s is replaced by a particular term in the sum (12) for the propagator the delta distributions reduce the integral to one over matching connections on the $`n`$ atomic boundaries and $`\mathrm{\Gamma }`$, where in each term the nodes are matched up according to the pairing of the $`\psi `$s and the boundary edges attached to each node are matched according to the permutation $`\sigma S_4`$ of the particular term of (12) being employed. (Actually, to make the matching of boundary edges corresponding to a given $`\sigma `$ unambigous we need to adopt a convention to specify the unpermuted order of these edges on the atoms and on $`\mathrm{\Gamma }`$. Since boundary edges $`e_a^b`$ are indexed by pairs of nodes $`a,b`$ the ordering of the $`\psi `$s, which is equivalent to an ordering of the nodes, provides a natural choice for the ordering of the boundary edges attached to a particular node $`a`$.) Glueing together the atoms and $`\mathrm{\Gamma }`$ in this way produces an admissible 2-complex $`J`$ with boundary $`\mathrm{\Gamma }`$. The corresponding contribution to (8) is $$\frac{\lambda ^n}{n!}(\frac{1}{5!})^n\frac{1}{sym(\mathrm{\Gamma })}𝑑𝐡\theta ^{}(𝐡_\mathrm{\Gamma })w(𝐡).$$ (13) To evaluate (8) we must add up the contributions (13) corresponding to each glueing, i.e. to each to each pairing of $`\psi `$s with permutations $`\sigma S_4`$ associated to each pair. Since (13) depends only on the structure of $`J`$ we can first sum over glueings that lead to the same $`J`$ by multiplying (13) by the number of such glueings, and then sum over $`J`$s. In our expansion the atoms and the nodes in each atomic boundary as well as in $`\mathrm{\Gamma }`$ are numbered, in an order arising from the way in which the (8) was written explicitly as an integral of a product of $`\psi `$s. When the atoms and $`\mathrm{\Gamma }`$ are glued together to form a 2-complex $`J`$ this reference numbering defines a numbering of the atoms, atomic boundary nodes and $`\mathrm{\Gamma }`$ nodes of $`J`$. Conversely, if any admissible, $`n`$ atom 2-complex bounded by $`\mathrm{\Gamma }`$ is equiped with a “good” numbering then it defines uniquely a glueing of the atoms and $`\mathrm{\Gamma }`$ with their reference numbering, and thus a pairing of the sequence of $`\psi `$s used in the expansion and a set of associated $`S_4`$ permutations. Here a good numbering is an arbitrary numbering of the atoms and of the nodes in each atom together with a numbering of the nodes of $`\mathrm{\Gamma }`$ such that $`\mathrm{\Gamma }`$ has the same adjacency matrix as in the reference numbering. From this we conclude, firstly, that our expansion generates all admissible, $`n`$ atom 2-complexes bounded by $`\mathrm{\Gamma }`$, so the sum should run over all such complexes, and secondly, that all glueings that lead to the same 2-complex arise from good numberings of the atoms and nodes in that 2-complex. There are $`n!(5!)^nsym(\mathrm{\Gamma })`$ good numberings, but renumbering does not necessarily produce a new glueing. A good renumbering of $`J`$ that produces the same gluing is a symmetry of $`J`$, because it defines a mapping of the set of cells of $`J`$ to itself that preserves all incidences, and conversely. There are thus $`n!(5!)^nsym(\mathrm{\Gamma })/sym(J)`$ glueings that produce the complex $`J`$ so (8) equals $$\underset{\{J𝐀_\mathrm{\Gamma }\|n(J)=n\}}{}\frac{\lambda ^n}{sym(J)}𝑑𝐡\theta ^{}(𝐡_\mathrm{\Gamma })w(𝐡).$$ (14) If the number of factors of $`\psi `$ in (8) is odd then $`n`$ atoms and $`\mathrm{\Gamma }`$ together have an odd number of nodes and so cannot be matched up to form a 2-complex bounded by $`\mathrm{\Gamma }`$. (14) therefore covers this case as well, because the integral (8) is zero while (14) is zero because $`\{J𝐀_\mathrm{\Gamma }\|n(J)=n\}`$ is empty. The theorem is thus established. Note that we could have restricted $`\psi (h_1,\mathrm{},h_4)`$ to be invariant under the gauge transformation $$h_ngh_n$$ (15) without affecting the result. If one writes $`\psi =\overline{\psi }+\mathrm{\Delta }\psi `$ where $`\overline{\psi }=_G𝑑g\psi (g𝐡)`$, the gauge average of $`\psi `$, then one finds $`I_0[\psi ]=I_0[\overline{\psi }]+I_0[\mathrm{\Delta }\psi ]`$ while $`𝒱[\psi ]=𝒱[\overline{\psi }]`$ and $`\mathrm{\Theta }[\psi ]=\mathrm{\Theta }[\overline{\psi }]`$ because $`V`$ and $`\theta `$ are gauge invariant. It follows that the integral over $`\mathrm{\Delta }\psi `$ cancels between numerator and denominator in the expectation value $$\mathrm{\Theta }^{}=𝒟\psi e^{I_0[\psi ]+\lambda 𝒱[\psi ]}\mathrm{\Theta }^{}[\psi ]/𝒟\psi e^{I_0[\psi ]}$$ (16) and therefore this expectation value is unchanged if $`\psi `$ is replaced by $`\overline{\psi }`$, i.e. is restricted to functions invariant under (15). (In (6) the denominator of (16) was set to 1. We may simultaneously set $`𝒟\overline{\psi }e^{I_0[\overline{\psi }]}=1`$ and leave out the denominator also in the $`\overline{\psi }`$ formulation). In the perturbation series of the $`\overline{\psi }`$ field theory the propagator (12) is replaced by one that requires the connections to match only up to gauge. However, the gauge invariance of $`V`$ and $`\theta `$ imply that the integrals over the connection can be evaluated in a gauge in which the connections on the atomic boundaries and $`\mathrm{\Gamma }`$ really do match. The integrals over the new gauge degrees of freedom at nodes then simply contribute factors of 1 because the normalized Haar measure is being used. ### 2.3 An example: The Turaev-Ooguri-Crane-Yetter model of BF theory A very simple gauge theory to which we can apply our scheme is BF theory. In this theory the action for the boundary value problem in which the connection is fixed on the boundary is $`tr[BF]`$ where $`F`$ is the curvature of the connection, $`B`$ is a 2-form which takes values in the Lie algebra of the gauge group $`G`$, and the trace is taken in this Lie algebra. Extremization of this action with respect to $`B`$ requires $`F=0`$, that is, a flat connection and in particular a flat connection on the boundary. Indeed a naive path integral quantization of this problem gives as the amplitude for the boundary connection a delta distribution with support on flat connections. It is therefore natural to discretize BF theory as a simplicial local lattice gauge theory in which the amplitude for the connection on the boundary of each simplex is a delta distribution with support on flat connections or, more concretely, $`V`$ is a delta distribution with support on flat connections on $`\mathrm{\Gamma }_5`$. Here flatness on $`\mathrm{\Gamma }_5`$ means that the holonomy around any closed loop of $`\mathrm{\Gamma }_5`$ is trivial. The only such connections are the trivial connection and its gauge transforms, for which $$h_i{}_{}{}^{j}=p_iq^{ij}$$ (17) with $`p_iG`$ characterizing the gauge at the 4-valent node $`i`$ while $`q^{ij}=q^{ji}G`$ does so at the bivalent node in the chain connecting $`i`$ and $`j`$. Thus $`V_{BF}(h_i{}_{}{}^{j}[h_j{}_{}{}^{i}]_{}^{1})`$ $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle }dp_k{\displaystyle \underset{i<j}{}}\delta (h_i{}_{}{}^{j}[h_j{}_{}{}^{i}]_{}^{1},p_ip_j^1)`$ (18) $`=`$ $`{\displaystyle \underset{k}{}}{\displaystyle }dp_k{\displaystyle \underset{l<m}{}}{\displaystyle }dq^{lm}{\displaystyle \underset{ij}{}}\delta (h_i{}_{}{}^{j},p_iq^{ij}).`$ (19) Substituting $`V_{BF}`$ into the formula (4) for the interaction term in the field action we find $`𝒱_{BF}[\psi ]`$ $`=`$ $`{\displaystyle \frac{1}{5!}}{\displaystyle _{G^{20}}}d^{20}hV_{BF}(h_i{}_{}{}^{j}[h_j{}_{}{}^{i}]_{}^{1})\psi (h_1{}_{}{}^{i})\mathrm{}\psi (h_5{}_{}{}^{i})`$ (20) $`=`$ $`{\displaystyle \frac{1}{5!}}{\displaystyle \underset{k}{}}{\displaystyle 𝑑p_k\underset{l<m}{}𝑑q^{lm}\psi (p_1q^{1j})\mathrm{}\psi (p_5q^{5j})}`$ (21) $`=`$ $`{\displaystyle \frac{1}{5!}}{\displaystyle \underset{l<m}{}}{\displaystyle 𝑑q^{lm}\overline{\psi }(q^{1j})\mathrm{}\overline{\psi }(q^{5j})}.`$ (22) $`I_0[\overline{\psi }]\lambda 𝒱_{BF}[\overline{\psi }]`$ is precisely Ooguri’s action for generating BF theory summed over 2-complexes. ## 3 Spin foam formulation ### 3.1 From local lattice gauge theories to spin foams The results of the previous section can also be understood within the spin foam framework. This has the advantage, among others, of replacing the functional integrals of (6) by more concrete multiple integrals over infinite sequences of real numbers, suggesting avenues toward regulating these integrals. Any local lattice gauge theory with compact gauge group $`G`$ can be given a spin foam formulation. The spin foam formulation is conjugate to the connection formulation in the sense that if a connection is a history of the configuration variables of the model a spin foam is essentially a history of the momentum variables - the states in the spin foam picture are a sort of Fourier transform of the states in the connection picture. The spin foam sum for the amplitude $$𝑑𝐡\theta ^{}(𝐡_\mathrm{\Gamma })w(𝐡),$$ (23) of a boundary state $`\theta `$ on a fixed admissible 2-complex is obtained by expanding $`V(𝐡_x)`$ for each atom and $`\theta (𝐡_\mathrm{\Gamma })`$ on a basis of “spin network functions” , and then integrating out the connection in each term of the resulting expansion of $`\theta ^{}w`$. Each non-zero term in this sum for (23) is associated to a spin foam. The translation of integrals over spacetime connections to sums over spin foams is described in detail in . Here we only outline its main features. They key necessary element for any reformulation of this type is an orthonormal basis of the space of distributions on the gauge group manifold. Since $`G`$ is assumed to be compact the Peter-Weyl theorem assures us that the matrix elements of the unitary irreducible representations,<sup>12</sup><sup>12</sup>12 Only one irreducible representation is selected from each class of unitarily equivalent representations. The basis of functions on $`G`$ consists of the matrix elements of these selected representations. $`U^{(r)}{}_{m}{}^{}{}_{}{}^{n}(g)`$ (with $`g`$ ranging over $`G`$, $`r`$ labeling the representation, and $`(m,n)`$ the matrix elements) form an orthogonal basis of distributions. Indeed, since $$_Gdg[U^{(r_1)}{}_{m_1}{}^{}{}_{}{}^{n_1}(g)]^{}U^{(r_2)}{}_{m_2}{}^{}{}_{}{}^{n_2}(g)=\frac{1}{dimr_1}\delta _{r_1r_2}\delta _{m_2}^{m_1}\delta _{n_1}^{n_2}$$ (24) with $`dimr`$ the dimensionality of the representation $`r`$, the set $`\{\sqrt{dimr}U^{(r)}{}_{m}{}^{}{}_{}{}^{n}\}_{r,m,n}`$ is an orthonormal basis. Any function of the parallel propagators along the edges of a graph may thus be expanded in terms of tensor products of representation matrices of the parallel propagators. Notice that a representation matrix of the parallel transporter along an edge is a two point tensor, with one index living at each end of the edge. The subspace of functions that are gauge invariant at a given vertex is thus spanned by functions obtained from the tensor product of representation matrices by contracting the indices at that vertex with an invariant tensor, also called an “intertwiner”. (An example of an intertwiner when $`G=SU(2)`$ is the antisymmetric tensor $`ϵ_{ijk}`$, which corresponds to a trivalent vertex with each incident edge carrying the spin 1 representation). An orthonormal basis of fully gauge invariant functions can be found by first selecting at each vertex an orthonormal basis $`\{W_I^𝐫\}_I`$ of the space of intertwiners for each set $`𝐫`$ of incident representations that admits intertwiners<sup>13</sup><sup>13</sup>13 That is to say, each $`𝐫=[r_1,r_2,\mathrm{},r_v]`$ such that $`r_1\mathrm{}r_v`$ contains the trivial representation, so that non-zero invariant tensors exist. and then forming the basis functions by contracting factors $`\sqrt{dimr_e}U^{(r_e)}{}_{m}{}^{}{}_{}{}^{n}(g_e)`$ for the edges $`e`$ with intertwiners belonging to the chosen intertwiner bases at the nodes. The data labeling a basis function, namely the representation on each edge and the basis intertwiner at each node, define a “spin network” on the graph,<sup>14</sup><sup>14</sup>14 Strictly the spin network consists of the subgraph carrying carrying non-trivial representations - edges that carry trivial representations are left off - together with the “colouring” data consisting of the non-trivial representations and the intertwiners. In our present context of lattice gauge theory this point is of no importance but in the continuum in which there is no convenient graph of all possible edges it is essential. and the basis functions are called “spin network functions”. Now consider expanding $`V(𝐠)`$ into a sum of spin network basis functions. $$V(𝐠)=\underset{𝐫_{..},𝐈_.}{}A^{(𝐫_{..})𝐈_.}[\underset{i<j}{}\sqrt{dimr_{ij}}]\varphi _{𝐈_.}^{(𝐫_{..})}(𝐠).$$ (25) (Recall $`g_{ij}=h_i{}_{}{}^{j}[h_j{}_{}{}^{i}]_{}^{1}`$). $`\varphi _{𝐈_.}^{(𝐫_{..})}(𝐠)`$ is the (normalized) basis spin network function on $`\mathrm{\Gamma }_5`$ with intertwiners $`W_{I_i}`$ at the 4-valent nodes and representations $`r_{ij}`$ on the edges connecting the 4-valent nodes. $$\varphi _{𝐈_.}^{(𝐫_{..})}(𝐠)=\underset{i<j}{}\sqrt{dimr_{ij}}U^{(r_{ij})}{}_{m_{ij}}{}^{}{}_{}{}^{n_{ij}}(g_{ij})\underset{k}{}W_{I_kn_{1k}\mathrm{}n_{k1k}}{}_{}{}^{m_{kk+1}\mathrm{}m_{k5}}.$$ (26) $`𝐫_{..}`$ denotes the whole matrix of representations $`r_{ij}`$ and $`𝐈_.`$ similarly denotes the sequence of basis intertwiners $`I_i`$. The factor $`[_{i<j}\sqrt{dimr_{ij}}]`$ in (25) has been separated from the coefficient $`A^{(𝐫_{..})𝐈_.}`$ for later convenience. Each spin network on the boundary of an atom defines a spin foam on the body of the atom: The intertwiner $`I_i`$ on the 4-valent node $`i`$ is the spin foam intertwiner on the spoke connecting the node with the 0-cell at the center of the atom; The representation $`r_{ij}`$ on the chain connecting nodes $`i`$ and $`j`$ defines the representation on the wedge of the atom bounded by the chain. The orthogonality relation (24) and the orthogonality of the basis intertwiners then ensure that when $`\theta ^{}w`$ is integrated over the connection the only terms that remain are those in which these atomic spin foams match up, in the sense that the intertwiners on the two halves of a 1-cell of $`J`$ agree and the representations on all the wedges of a 2-cell agree.<sup>15</sup><sup>15</sup>15 They must agree in the sense that if two wedges of a 2-cell have matching (“coherent”) orientations then they must carry the same representation, whereas if they have opposite orientations they must have complex conjugate representations. These representations will be discussed more a little further on. Moreover, they ensure that the intertwiners and representations match at the boundary $`\mathrm{\Gamma }`$ where the spin foam meets the spin network in the expansion of $`\theta `$. The weight, or amplitude, of each spin foam $`S`$ in the sum for (23) that remains once the connection is integrated out is<sup>16</sup><sup>16</sup>16 This formula should be viewed as somewhat schematic some details have been ignored. The matching of basis intertwiners on coincident nodes is not between basis intertwiners in the same invariant tensor space but rather in dual spaces. If we use a raised index $`I`$ to label the dual basis intertwiners we see that the matching of intertwiner indices $`A`$ in the product (27) should always be between a downstairs and an upstairs index. For many groups one may choose the intertwiner bases so that raising or lowering the intertwiner index is trivial, leaving $`A`$ unchanged, but we are not sure that this is always true. Sign factors, not included in our formula, are also present for “psuedo-real” representations that some groups have which depend on the relative orientations of the coincident boundary edges on which neighboring atoms meet. The proper resolution of these issues is explained in outline at the end of this section. $$c_S^{}\underset{x\text{0-cells of}J}{}A^{(𝐫_{..})𝐈_.}\underset{y\text{2-cells of}J}{}dimr_y$$ (27) where $`c_S`$ proportional to the coefficient of the spin network $`S`$ (the boundary of the spin foam) in the expansion of $`\theta `$, being obtained from the coefficient by dividing by a factor of $`\sqrt{dimr}`$ for each edge of $`\mathrm{\Gamma }`$. ### 3.2 Divergences, regularization, and renormalization Spin foam models can have divergences when $`J`$ has closed 2-surfaces. For instance, in the model of BF theory that we have discussed the weight $`w`$ of a connection has redundant delta distributions, and is thus infinite, whenever $`J`$ contains topological 2-spheres. Other closed surfaces can also contribute singular factors<sup>17</sup><sup>17</sup>17 The singular factor associated with a closed surface in BF theory is $$\underset{rR}{}(dimr)^\chi +\underset{rPR}{}(dimr)^\chi +\kappa \underset{rC}{}(dimr)^\chi $$ (28) where $`\chi `$ is the Euler characteristic of the surface, $`\kappa =1`$ if the surface is orientable and $`0`$ if it is not, $`R`$ is the set of (unitary equivalence classes of) real irreducible representations, $`PR`$ is the set of psuedo-real irreps, and $`C`$ is the set of complex irreps (for definitions see ). When $`G=U(1)`$ we obtain a singular factor for every orientable closed surface, whereas for $`G=SU(2)`$ only the four surfaces with Euler characteristic $`\chi 0`$ create divergences. In the spin foam sum the infinities take the form of a divergence of the sum over representations and can thus be regulated by somehow cutting off this sum. This can be accomplished by replacing $`G`$ with the quantum group $`G_q`$ with $`q1`$ a root of unity, which has a finite set of inequivalent irreducible representations. The spin foams are coloured with the representations and intertwiners of $`G_q`$, and $`A`$ is replaced by its natural generalization to such representations yielding a finite sum. One can also simply cut off the sum, leaving the model otherwise unchanged. In lattice models of topological theories, or Yang-Mills theories for that matter, the continuum limit theory is usually defined not by a sum over lattices but by the limiting values of the observables of the theory as the lattice is refined. In the case of simplicial models the existence of this limit requires that the model be independent of the triangulation used (at least so long as the triangulation is very fine, i.e. has very many simplices). This requires that the regulated model be renormalized, because the number of divergences, and thus the number of the large, regulator dependent factors in the transition amplitudes, depends on the simplicial complex used. For BF theory on a triangulated manifold this means dividing out a factor $`\delta (0)_{\text{regulated}}=[_r(dimr)^2]_{\text{regulated}}`$ for each independent 2-sphere in $`J`$.<sup>18</sup><sup>18</sup>18 The expansion on the basis of matrix elements of the Haar measure delta distribution on $`G`$ is $$\delta (g)=\underset{r}{}dimrtrU^{(r)}.$$ (29) This is the origin of the factor $`\delta (0)_{\text{regulated}}^{N_0N_1}`$ (with $`N_0`$ the number of 0-simplices and $`N_1`$ the number of 1-simplices in the triangulation) in the weight of a history in simplicial BF theory . Renormalization is more subtle in the context of a general admissible 2-complex because the spheres do not necessarily span all the closed 2-surfaces (for instance one can make a 2-complex having a torus as it’s only closed surface). As far as we know the renormalization of the model has not been carried out in this wider context. In the present paper we are of course advocating a different approach to the continuum limit, in which one sums over 2-complex spacetimes. We have not studied the issue of renormalization in this approach. ### 3.3 The field theory in terms of spin foams Now let’s see how the field theory (2) can be expressed in terms of spin foams and networks. While this is qualitatively quite straightforward it turns out to be a little tricky in detail. The difficulty has to do with the fact that the spin network basis is not quite unique. Some conventions have to be introduced to define a particular basis. When translating from an integral over connections to a sum over spin foams on a particular, given, 2-complex one can choose the spin network bases on the atomic boundaries so as to simplify the calculations on that particular 2-complex. On the other hand, in order to express our field theory in spin foam terms, which boils down to expressing it in terms of the coefficients in spin network expansions of $`V`$, $`\theta `$, and $`\psi `$, we have to choose a particular spin network basis on the prototypical atomic boundary $`\mathrm{\Gamma }_5`$. This is then the basis on every atomic boundary in the complexes generated by the field theory. The resulting inflexibility in the choice of basis complicates the calculation of the amplitudes of Feynman diagrams. We have taken a compromise route. We present the field theory in spin foam language using a basis that makes the presentation as simple and symetrical as possible, obtaining a very neat result. Then we outline how the amplitudes of the Feynman diagrams can be calculated by exploiting the fact that one is free to change to adapted spin network bases independently for each Feynman diagram, making the calculation as simple as in a spin foam model on a fixed 2-complex spacetime. It may well be that the results could be presented more cleanly using a slightly more abstract “basis independent” approach to the spin network basis, in which the conventional choices that fix a particular basis are not made. Here we will stick with definite bases. Let us begin the translation of the field theory by expanding $`\psi `$ using the Peter-Weyl theorem: $$\psi (𝐡)=\underset{𝐫}{}b^{(𝐫)𝐦}{}_{𝐧}{}^{}\underset{t=1}{\overset{4}{}}\sqrt{dimr_t}U^{(r_t)}{}_{m_t}{}^{}{}_{}{}^{n_t}(h_t).$$ (30) The symmetry of $`\psi `$ under permutations of its arguments requires the coefficients $`b`$ to be similarly symmetric: $$b^{(𝐫)𝐦}{}_{𝐧}{}^{}=b^{(\sigma [𝐫])\sigma [𝐦]}{}_{\sigma [𝐧]}{}^{}\sigma S_4.$$ (31) The free action is easily expressed in terms of these coefficients if we write it as $`I_0[\psi ]=\frac{1}{24!}_{G^4}\psi ^{}\psi d^4h`$. This is valid because $`\psi `$ is real. We find, using (24), that $$I_0[\psi ]=\frac{1}{24!}\underset{𝐫}{}[b^{(𝐫)𝐦}{}_{𝐧}{}^{}]^{}b^{(𝐫)𝐦}{}_{𝐧}{}^{}.$$ (32) The corresponding propagator is<sup>19</sup><sup>19</sup>19 It can be evaluated by translating the field propagator $`\psi (𝐡)\psi (𝐡^{})_0`$ or directly from (32). In carrying out the Gaussian integral in the second approach the reality conditions satisfied by $`b`$ (spelled out further on in the text) must be taken into account. $$b^{(𝐫_1)𝐦_1}{}_{𝐧_1}{}^{}[b^{(𝐫_2)𝐦_2}{}_{𝐧_2}{}^{}]_{}^{}_0=\underset{\sigma S_4}{}\delta _{𝐫_1\sigma [𝐫_2]}\delta _{\sigma [𝐦_2]}^{𝐦_1}\delta _{𝐧_1}^{\sigma [𝐧_2]}$$ (33) with $`\delta _𝐚^𝐛=_{t=1}^4\delta _{a_t}^{b_t}`$. To express the interaction term $`𝒱`$ in terms of the coefficients $`b`$ in as clean a way as possible we shall introduce a further representation theoretic tool. If $`g_{12}G`$ defines parallel transport on the edge from node 1 to node 2 then $`g_{21}=[g_{12}]^1`$ defines parallel transport along the inverse edge from 2 to 1. A unitary representation matrix $`U_{m_1}^{(r)}{}_{}{}^{m_2}(g_{12})`$ can be expressed in terms of $`g_{21}`$ as $`[U^{(r)}(g_{21})]_{m_1}^1{}_{}{}^{m_2}=[U_{m_2}^{(r)}{}_{}{}^{m_1}(g_{21})]^{}`$. It is just the parallel propagator from 2 to 1 in the complex conjugate of the representation $`r`$. Thus a spin network function of the parallel transporters along the oriented edges of a graph can be expressed just as well in terms of parallel transporters of the edges with different orientations, provided the representations on the reversed edges are replaced by their complex conjugates. Now recall that in a spin network basis the representations $`U^{(r)}`$ that may be placed on edges are particular representatives chosen one out of each unitary equivalence class. Moreover, in order to keep the definition of the spin network basis simple we shall use the same set of representatives on all edges. In general the conjugate representation $`[U^{(r)}]^{}`$ will not be the representative $`U^{(r^{})}`$ of its own equivalence class, it will only be unitarily equivalent to it. For instance, all representations of $`SU(2)`$ are equivalent to their conjugates, but unitary irreducible representations of non-integer spin are not real and so are not equal to their conjugates. We therefore introduce $`\epsilon _{mn}^{(r)}`$, a unitary matrix such that $$[U^{(r)}]^{}=\epsilon ^{(r)}U^{(r^{})}\epsilon ^{(r)}$$ (34) and let $`\epsilon ^{(r)mn}=[\epsilon _{mn}^{(r)}]^{}`$, which is the inverse of $`\epsilon _{mn}^{(r)}`$ in the sense that $`\epsilon _{mx}^{(r)}\epsilon ^{(r)nx}=\delta _m^n`$. A few notes on $`\epsilon ^{(r)}`$: $``$ It follows from (34) that $`\epsilon _{mn}^{(r)}`$ is an invariant tensor (an intertwiner) with $`m`$ a covector index of the representation $`r^{}`$ and $`n`$ a covector index of the representation $`r`$. $``$ (34) determines $`\epsilon ^{(r)}`$ up to phase. Consequently $`\epsilon ^{(r^{})mn}`$ equals $`\epsilon ^{(r)nm}`$ up to a phase factor, and thus when $`r=r^{}`$, $`\epsilon ^{(r)}`$ is either purely symmetric or antisymmetric.<sup>20</sup><sup>20</sup>20 For example consider the spin $`1/2`$ and spin $`1`$ representations of $`SU(2)`$ in the standard eigenbasis of the generator $`J_z`$. In the spin $`1/2`$ representation $`\epsilon _{mn}^{(\frac{1}{2})}=ϵ_{mn}`$, while in the spin 1 representation $`\epsilon _{mn}^{(1)}=\delta _{m,n}`$, so it is antisymmetric for spin $`1/2`$ and symmetric for spin 1. Indeed one finds that $`\epsilon _{mn}^{(j)}`$ is symmetric for all integer $`j`$ and antisymmetric for all half odd integer $`j`$. When $`rr^{}`$ one is of course free to choose the representative of the equivalence class of the conjugate representation as one likes. One could simply choose the conjugate representation itself and use $`\epsilon _{mn}^{(r)}=\delta _{mn}`$. We are therefore free to set $`\epsilon ^{(r)}=\epsilon ^{(r^{})}`$ for all $`r`$.<sup>21</sup><sup>21</sup>21 When $`r=r^{}`$ and $`\epsilon ^{(r)}`$ is symmetric the representation $`r`$ is real; If $`r=r^{}`$ and $`\epsilon ^{(r)}`$ is antisymmetric $`r`$ is psuedo-real; Finally, if $`rr`$ $`r`$ is complex. See . $``$ $`\epsilon ^{(r)}`$ can be used as a (possibly antisymmetric) metric to turn a representation $`r`$ vector (carrying an upstairs index) into a representation $`r^{}`$ covector (carrying a downstairs index) according $`a_m^{(r^{})}=\epsilon _{mx}^{(r)}a^{(r)x}`$. The inverse, index raising, operation is $`a^{(r)x}=\epsilon ^{(r)xm}a_x^{(r^{})}`$. Of course complex conjugation also turns unitarily transforming $`r`$ vectors into $`r^{}`$ covectors and vice versa. The index positioning in all our equations is consistent with this fact. The relation (34) can be used to define a modified, more symmetric spin network expansion of $`V`$ which depends only minimally on the orientations chosen for the edges of $`\mathrm{\Gamma }_5`$. Notice that $`U_{m_1}^{(r)}{}_{}{}^{n}(h_1h_2^1)`$ $`=`$ $`U_{m_1}^{(r)}{}_{}{}^{x}(h_1)[U_n^{(r)}{}_{}{}^{x}(h_2)]^{}`$ (35) $`=`$ $`U_{m_1}^{(r)}{}_{}{}^{x}(h_1)\epsilon _{xy}^{(r)}U_{m_2}^{(r^{})}{}_{}{}^{y}(h_2)\epsilon ^{(r)nm_2}.`$ (36) Thus adopting the notation $`r_{ij}=r_{ji}^{}`$ for $`i>j`$ we may write $$V=\underset{𝐫_{..},𝐈_.}{}A^{(𝐫_{..})𝐈_.}[\underset{ij}{}\sqrt{dimr_{ij}}]U^{(r_{ij})}{}_{m_{ij}}{}^{}{}_{}{}^{x_{ij}}(h_i{}_{}{}^{j})\underset{k}{}W_{I_k}^{m_{k1}\mathrm{}m_{kk1}m_{kk+1}\mathrm{}m_{k5}}\underset{i<j}{}\epsilon _{x_{ij}x_{ji}}^{(r_{ij})}.$$ (37) The lower indices of the intertwiners appearing in (25) have been raised using $`\epsilon `$. The invariance of $`\epsilon `$ implies that the resulting tensors $`W_{I_k}^{𝐦_{k.}}`$ are also invariant and thus intertwiners for the incident representations $`r_{ki}`$. It is easy to verify that they form an orthonormal basis of such intertwiners. (37) is thus a spin network expansion of $`V`$ with the bivalent nodes assigned the intertwiners $`\epsilon _{x_{ij}x_{ji}}^{(r_{ij})}`$ (which is $`\sqrt{dimr_{ij}}`$ times the unique (mod phase) normalized bivalent intertwiner). In (37) the orientations of the edges manifest themselves only in the ordering of the indices in these $`\epsilon `$s. Evaluating $`𝒱`$ is now straightforward. If we use the reality of $`\psi `$ to replace $`\psi `$ by $`\psi ^{}`$ in (4) the integrals can be carried out using (24) and we obtain $$𝒱=\frac{1}{5!}\underset{𝐫_{..},𝐈_.}{}A^{(𝐫_{..})𝐈_.}\underset{k}{}[b^{(𝐫_k)𝐦_k}{}_{𝐱_k}{}^{}]^{}W_{I_k}^{𝐦_k}\underset{i<j}{}\epsilon _{x_{ij}x_{ji}}^{(r_{ij})}.$$ (38) $`𝒱`$ is even simpler when expressed in terms of the spin network coefficients in an expansion of the “gauge invariant” field $`\overline{\psi }(𝐡)=_G\psi (g𝐡)𝑑g`$. A spin network type expansion of this field can be obtained by inserting into the expansion (30) of $`\psi `$ the projector $$_G\underset{t=1}{\overset{4}{}}U^{(r_t)}{}_{m_t}{}^{}{}_{}{}^{n_t}(g)dg=\underset{I}{}[W_I^{(𝐫)m_1m_2m_3m_4}]^{}W_I^{(𝐫)n_1n_2n_3n_4}$$ (39) onto invariant $`r_1r_2r_3r_4`$ tensors. Thus $$\overline{\psi }(𝐡)=\underset{𝐫,I}{}c^{(𝐫)I}{}_{𝐧}{}^{}W_{I}^{(𝐫)𝐦}\underset{t=1}{\overset{4}{}}\sqrt{dimr_t}U^{(r_t)}{}_{m_t}{}^{}{}_{}{}^{n_t}(h_t)$$ (40) with $`c^{(𝐫)I}{}_{𝐧}{}^{}=b^{(𝐫)𝐦}{}_{𝐧}{}^{}[W_I^{(𝐫)𝐦}]_{}^{}`$. In terms of the coefficients $`c`$ the interaction term is $$𝒱=\frac{1}{5!}\underset{𝐫_{..},𝐈_.}{}A^{(𝐫_{..})𝐈_.}\underset{k}{}[c^{(𝐫_k)I_k}{}_{𝐱_k}{}^{}]^{}\underset{i<j}{}\epsilon _{x_{ij}x_{ji}}^{(r_{ij})},$$ (41) It’s just $`A`$ times the result of contracting together the indices of five $`c^{}`$ in the pattern of a 4-simplex. More precisely it can be obtained from (37) by replacing the intertwiners $`W_I`$ by $`[c^I]^{}`$ and setting all the parallel transporters to $`\mathrm{𝟏}`$. This prescription can also be applied to $`\theta `$ to express $`\mathrm{\Theta }`$ in terms of the coefficiants $`c^{}`$. $`\epsilon `$ also allows us to express the reality conditions satisfied by the coefficients $`b`$: $`\psi ^{}=\psi `$ and (34) implies that $$b^{(𝐫^{})𝐦}{}_{𝐧}{}^{}=[b^{(𝐫)𝐱}{}_{𝐲}{}^{}]^{}\epsilon ^{(𝐫)\mathrm{𝐦𝐱}}\epsilon _{\mathrm{𝐧𝐲}}^{(𝐫)}$$ (42) (with $`\epsilon ^{(𝐫)\mathrm{𝐦𝐧}}=_{t=1}^4\epsilon ^{(r_t)m_tn_t}`$). We are now in a position to rewrite the functional integral (6) as an ordinary integral over the coefficients $`b`$ (which are finite in number if the spin foam model has been regulated by cutting off the sum over representations). Moreover we may evaluate the Feynman diagrams of the perturbation series in $`\lambda `$ using the propagator (33) and the reality conditions (42). What we have given is just about the simplest, most symmetric statement of the spin foam formulation of the theory. However, it is not the most convenient form for actually evaluating the Feynman diagrams. In fact it is best to first select the Feynman diagram one wishes to evaluate, corresponding to a particular 2-complex $`J`$, and then choose the spin network basis for expanding $`\mathrm{\Theta }^{}`$ and the factor of $`𝒱`$ corresponding to each atom. To define this adapted set of spin network bases we choose an orientation for each 1-cell and 2-cell of $`J`$. The orientations on the 2-cells induce an orientation on each edge of $`\mathrm{\Gamma }`$ and each atomic boundary edge that matches the sense of circulation of the 2-cell that it cuts. The orientations on the 1-cells define a sign on each node, which is positive if the 1-cell is outgoing at the node and negative if it is incoming. Now we choose spin network basis functions corresponding to these orientations of the edges. We obtain an expansion of $`V`$, rather like that given in (25) with coherent orientations on the chain of two edges connecting a pair of nodes. However, we do not necessarily have two outgoing and two incoming edges at each node as in (25). The intertwiner bases are also adapted. At each positive node a basis $`\{W_I^{(𝐫)}\}`$ with indices suitably lowered for incoming edges, is used. At the corresponding negative node (the negative node that lies on the same 1-cell) the intertwiners $`[W_I^{(\sigma [𝐫])}]^{}`$ are used. Here $`\sigma S_4`$ is the mapping of the incident edges at the negative node to the corresponding incident edges at the positive node. These complex conjugate and permuted intertwiners are still orthonormal and have the index positions compatible with our convention for the orientations of the edges. With these conventions (and with the correspondingly adjusted coefficients $`A`$ and $`c_S`$) the amplitude of the Feynman diagram is relatively easy to obtain. The result is as described in (27): It is a sum over histories consisting of representations and basis intertwiners, with only histories in which every 2-cell $`y`$ carries a single representation $`r_y`$ \- the common representation on the atomic boundary edges that cut the 2-cell and on the edges of $`\mathrm{\Gamma }`$ that bound it - and each 1-cell carries a single intertwiner basis index $`I`$, the common value of the indices at the positive and negative nodes on the 1-cell. Such a history is a spin foam. The amplitude for each history is the product of the coefficients $`A`$ for the atoms times $`c_S^{}`$ and a factor $`dimr_y`$ for each 2-cell. This last factor results from the factors of $`\epsilon _{x_{ij}x_{ji}}^{(r_{ij})}`$ that appear in (38) and (41), associated to each wedge. When the orientations of the edges are changed to match those of the 2-cells these are replaced by Kronecker deltas. Then, when the propagators for $`b`$ are substituted in, the Kronecker deltas associated to wedges in a given 2-cell end up being contracted in a chain around around the 2-cell, thus contributing a factor $`dimr_y`$. ## 4 Some closing remarks ### 4.1 On the possible finiteness of the field theory corresponding to regulated spin foam models The number of admissible 2-complexes increases very rapidly with the number of atoms. In fact the number of complexes of $`n`$ atoms with a given boundary $`\mathrm{\Gamma }`$ having $`m`$ 4-valent nodes is given approximately by $$\underset{\{JgothA_\mathrm{\Gamma }\|n(J)=n\}}{}=(\frac{4!}{2})^{(5n+m)/2}\frac{(5n+m)!}{([5n+m]/2)!}\frac{1}{n!(5!)^n}\frac{1}{sym(\mathrm{\Gamma })}.$$ (43) This increases as $`(n!)^{3/2}`$ for large $`n`$, so unless the vertex function $`V`$ is very special indeed the radius of convergence of the series (7) is zero. What prospect is there of making sense of the functional integral (6)? In fact it is not unreasonable to hope that once the spin foam model is regulated by cutting off the sums over representations (as we discussed in the context of BF theory) that the integral (6) can be assigned a finite value by analytic continuation. We will suppose that the finite set of representations summed over in the cut off model includes the representation $`r^{}`$ whenever it includes $`r`$. Then $`\psi `$ in the cut off model is a real function on $`G^4`$ determined by a finite set of parameters. These paramenters can be the $`b^{(𝐫)𝐦}_𝐧`$, but we will use a set of linear combinations $`\{x_p\}_{p=1}^N`$ of these chosen so that they are real and $`I_0[\psi ]=𝐱^T𝐱`$. Then (6) becomes $$_{𝐑^N}d^Nxe^{𝐱^T𝐱}e^{\lambda 𝒱^{p_1\mathrm{}p_5}x_{p_1}\mathrm{}x_{p_5}}\mathrm{\Theta }^{q_1\mathrm{}q_m}x_{p_1}\mathrm{}x_{p_m},$$ (44) a function of $`\lambda `$ which we shall call $`F(\lambda )`$. For simplicity we will suppose that $`V(𝐠)`$ is real, as is the case in BF theory, and in the models of and , then also $`𝒱^{n_1\mathrm{}n_5}`$ is real and we see that the integral is convergent when $`\lambda `$ is pure imaginary - the integral is a Gaussian times a function of modulus $`O(\|x\|^m)`$. Can the result be analytically continued away from the imaginary axis? To see that the answer is quite possibly “yes” consider the simplest analog of the integral (44): $$f(\lambda )=_{\mathrm{}}^{\mathrm{}}𝑑xe^{x^2}e^{\lambda x^5}.$$ (45) The coefficients in the formal power series expansion of $`f`$ about $`\lambda =0`$ also diverge as $`(n!)^{3/2}`$ for large $`n`$ and, like (44) the integral converges for purely imaginary $`\lambda `$. Moreover, $`f`$ can be continued to the entire complex plane, with a fivefold branch cut extending from $`0`$ to $`\mathrm{}`$. To show this one first writes $`f`$ as $`f(\lambda )=h(\lambda )+h(\lambda )`$ where $$h(\lambda )=_0^{\mathrm{}}𝑑xe^{x^2}e^{\lambda x^5}.$$ (46) This integral is convergent when $`\mathrm{}\lambda 0`$. It may be reexpressed in terms of $`\lambda ^{1/5}`$ and the rescaled integration variable $`y=\lambda ^{1/5}x`$. From this form one obtains a series for $`h`$ in powers of $`\lambda ^{1/5}`$ that is convergent for all finite values of this parameter, and the claimed result follows. The first step in this argument can be repeated for the regulated functional integral (44). Let $`M`$ be the union of rays in $`𝐑^4`$ on which $`𝒱^{n_1\mathrm{}n_5}x_{n_1}\mathrm{}x_{n_5}0`$. Then $$F(\lambda )=H(\lambda )+H(\lambda )$$ (47) with $$H(\lambda )=_Md^Nxe^{𝐱^T𝐱}e^{\lambda 𝒱(𝐱)}\mathrm{\Theta }^{}(𝐱).$$ (48) Again we find that the integral converges for $`\mathrm{}\lambda 0`$. The question is now whether it can be continued beyond this domain. It seems likely that this would be the case for generic models, especially as it is finite on $`\mathrm{}\lambda =0`$. If it is continuable then $`F(\lambda )`$ will be defined and finite on an open region in the complex plane, possibly including $`\lambda =1`$ which is the value corresponding most closely to a simple sum of the spin foam model over all admissible 2-complexes. ### 4.2 Generalizations of the formalism Our formalism can easily be generalized to a wider class of 2-complexes. We have allowed only atoms which are dual 2-skeletons of 4-simplices, which means that they have five spokes each of which is four valent. The field theory can be extended so that it generates 2-complexes including atoms with any given number $`p`$ of 4-valent spokes (which are dual to 4-polyhedra bounded by $`p`$ 3-simplices). All that is necessary is to include a suitable interaction term in the action (2) formed from a vertex function for the new type of atom like $`𝒱`$ was formed from the vertex function $`V`$ and $`\mathrm{\Theta }`$ from the boundary state $`\theta `$. The expression (23) for the amplitude of the state $`\theta `$ on the boundary $`\mathrm{\Gamma }`$ of $`J`$ can be viewed as the partition function for a closed 2-complex consisting of $`J`$ and one copy of a new type of atom with boundary $`\mathrm{\Gamma }`$ ($`\mathrm{\Gamma }`$ with reversed orientation) and vertex function $`\theta ^{}`$. Adding $`\mathrm{\Theta }^{}`$ to the action would generate all complexes consisting of both this new type of atom and the original type. Even with the above generalization all atoms have spokes of valence four and thus all bounding graphs of atoms and of complexes have non-trivial nodes of valence four only. The theory may be generalized to acomodate other valences $`q`$ by introducing additional fields depending on $`q`$ group elements instead of four like $`\psi `$ does, and adding a quadratic free action to the total action for each such field. Any model having a finite set of types of atoms can be handled by these means, including of course models with spacetimes of any dimensionality - to the extent that this dimensionality can be captured in the combinatorial structure of the atoms. ### 4.3 Matrix models When our formalism is applied to two dimensional models with atoms dual to triangles we find that the analog of $`c^{(𝐫)I}_𝐧`$ is $`c^{(r_1)}{}_{n_1n_2}{}^{}b^{(r_1,r_2)m_1m_2}{}_{n_1n_2}{}^{}\frac{1}{\sqrt{dimr_1}}\delta _{r_1r_2^{}}\epsilon _{m_1m_2}^{(r_1)}`$ ($`r_1`$ has to equal $`r_2^{}`$ because $`r_1r_2`$ contains the trivial representation only in this case. $`\epsilon ^{(r)}/\sqrt{dimr}`$ is the only normalized bivalent intertwiner). $`c^{(r)}`$ is essentially the matrix of matrix models : The triangular matrix model is recovered by choosing the vertex function to be $`V(g_{ij})=A(dimr)^3/3!trU^{(r)}(g_{12}g_{23}g_{31})`$ so that only one representation $`r`$ appears in its spin network expansion. Then $`𝒱=A/3!c_{n_1}^{(r)}{}_{}{}^{n_2}c_{n_2}^{(r)}{}_{}{}^{n_3}c_{n_3}^{(r)}^{n_1}`$. ## Acknowledgements We thank Roberto DePietri for discussions and help, and M.R. would like to thank Rodolfo Gambini for a discussion on the subject of this paper. This work was partially supported by NSF grant PHY-95-15506, PPARC grant PPA/6/s/1998/00613 and a gift from the Jesse Phillips Foundation.
warning/0002/math0002236.html	ar5iv	text	# SOME COMMENTS ON CATEGORIES, PARTICLES AND INTERACTIONS ## 1 Introduction It is well–known that particles like baryons, nuclei, atoms or molecules are characterized by their own specific excitation spectrum. The existence of these spectra is one of the fundamental properties of the structure of matter. Suppose that we have a particle system with a collection of bound–state energy levels. There is a ground state and there are several excited states. It is also well–known that there are some transitions between different levels. They are results of interactions of the system with an external field or with other particle systems. These transitions which need more energy are known as excitation processes. They are the result of absorption of quanta of an external field. On the other hand there are transitions connected with an energy spending, they correspond to some decay processes. It is interesting that there are also processes with no energy change. Let us consider for instance these processes which can be described as sequences of vertex interactions of particle charges with an external quantum field. Charged particles are transformed under these interactions into a composite nonlocal discrete system which contains charges and quanta of the field. These systems are said to be dressed particles . We describe the structure of these dressed particles as a lattice with $`n`$ sites, $`n=1,2,\mathrm{}`$. Every lattice site is a center for a vertex interaction of charge with the external quantum field. We assume that there are $`N`$ elementary excitation states on every lattice sites. There are also collective excitation states. It should be interesting to consider the general formalism for the study of all possible collective excitations of such system. We can imagine our lattice as a $`d`$-dimensional space (a manifold) equipped with $`n`$ distinct points as lattice sites. One can consider a quantum dot, spin chains, or a set of vertex interactions of particles moving in two–dimensional space under influence of transversal magnetic field as examples of such systems corresponding for $`d=0,1,2`$, respectively. Our fundamental assumptions are that a collection of $``$ initial configurations of the system, $``$ elementary particle processes. is given. It seems to be natural to assume that every possible configurations of the system can be obtained as a result of certain physical processes. We also assume that every process can be described as a sequence of elementary ones. These elementary processes represent elementary acts of lattice interactions. If all final configurations for the system under consideration can be described in an unique way as a result of transformation of an initial configuration, then we say that the system is equipped with a category symmetry or coherent evolution. Every such transformation is said to be a evolution transformation. This means that our category symmetry is in fact a formalism for the description of particle interactions. The problem is to determine for a given system the smallest collection of symmetry transformations generating all others in an unique way. The classical notion of the concept of symmetry in physics is based on group theory. The role played by the group representation theory for the study of symmetries in particle physics is well–known. The construction of a tensor product of representations is essential for such study. For instance, it allows to built states for composite systems of particles from single particles ones. The unitary symmetry and corresponding quark model is here a good example. It is known that we need a comultiplication for a tensor product of representations. It is interesting that a comultiplication does exist for a large class of $`q`$–deformed universal enveloping algebras. Hence they provide new possibilities for the study of particles, fields and their interactions in mathematical physics. Also categories which contain a bifunctor $`:\times `$ called a monoidal operation provide some additional possibilities. Such categories are said to be monoidal. They can contain more structures and operations. The name ”monoidal” indicate, that there is one essential operation, just the monoidal one, other operations play an auxiliary role. The bifunctor $``$ plays a similar role like the usual tensor product of group representations. Note that a related subject has been studied previously by several authors and others. Formalism corresponding to the braided symmetry has been developed mainly by Majid . Categories in the context of quantum groups has been presented by Kassel . The application of categories in the topological quantum field theory has been considered by Sawin . Similar formalism has been also developed previously by the author . Note that all these studies can be included in our general scheme. One can also consider the $`q`$–extended supersymmetry concept as a particular example of our general formalism. Our considerations are mainly motivated by applications for the investigation of interacting systems of particles in low–dimensional spaces, but there are more different possible applications. It is known that the study of certain integrable models on a lattice leads to the investigation of some new formalism . Hence it is interesting to study all these additional possibilities for the developing of the formalism beyond of the quantum mechanics and field theory. In this paper we are going to study these additional possibilities in a general manner in terms of monoidal categories. All our considerations are on abstract algebraic level. We would like to consider the most fundamental algebraic structures suitable for the description of particle interactions. We would like also to discuss the physical application for the classification of interacting particle systems. One can further develop our concept in terms of quantum von–Neumann algebras and their representations . The paper is organized as follows. In Section 2 the general concept of particle interactions is considered in terms of monoidal categories. Particle processes are described as certain transformations of categories. The essential problem is the unique description, it is related to the coherence in categories . In Section 3 categories relevant for our goal are described in details. Commutation relations for creation and annihilation quantum processes are described as certain specific transformations. They lead to the system with generalized statistics . An introduction to the category theory is given in the Appendix. We believe that our approach can be useful for the deeper understanding of such new methods in quantum optics or both condensed matter and particle physics. ## 2 General considerations Let us consider a system of hard core particles moving on certain $`d+1`$–dimensional space–time manifold under influence of some external field. All our considerations are based on the assumption that there is the vacuum state $`\mathrm{𝟏}\|0`$, the lowest energy elementary excited states $`\{\|i\}_{i=1}^N`$, and their conjugated states $`\mathrm{𝟏}^{}0\|`$, and $`\{i\|\}_{i=N}^1`$ with the scalar product $`i\|j\text{}`$. We also assume that there are collective excitations of the system which can be described as a result of certain multiple product of elementary excitations. There is an energy gap between the vacuum state $`\mathrm{𝟏}\|0`$ and the lowest energy excited state $`\|i`$. A finite set of $`N`$ operators $`L=\{x^i\}_{i=1}^N`$ is given as a starting point for our considerations. Every such operator act on the Hilbert space of functions on $`d`$–dimensional space. We also assume that these operators transform functions representing the ground state of the system into states representing elementary excitations of the system, i.e. $`\|i:=x^i\|0`$. In this way these operators represent elementary excited states of our system. Hence these operators are said to be primary. For the description of other states representing for instance collective excitations we need a product of operators. Such product need not forms a closed algebra but it must be defined in an unique way. A product of $`n`$ arbitrary primary operators should represents a collective $`n`$–tuple excitation. It is an analog of a $`n`$ particle state. We would like to study such product in terms of the category theory. If $`x^i`$ is a primary operator, then there is a corresponding vector space $`𝒰=𝒰(x^i)`$. It is formally a –linear span of $`x^i`$, i. e. $`𝒰(x^i):=\{\alpha x^i;\alpha \text{}\}`$. The –linear span of the ground state $`\mathrm{𝟏}`$ is denoted by $`I`$. It said to be the unit object. If $`𝒰`$ and $`𝒱`$ are –linear spans of $`x^i`$ and $`x^j`$, respectively, then the linear span corresponding for certain product of these operators is denoted by $`𝒰𝒱`$ and is also said to be a product of $`𝒰`$ and $`𝒱`$. If for example $`𝒰`$ represents charged particles excitation and $`𝒱`$ some quanta, then the product $`𝒰𝒱`$ describes the composite system containing both particles and quanta. This means that the operation $`:𝒰\times 𝒱𝒰𝒱`$ describes the ”composition” process of states. Such process tell us how to built a space of composite quantum states of the system from elementary ones. Hence it can be also understood as a generalization of the usual tensor product of group representations. Observe that the arrow $`𝒰𝒰𝒱`$ describes the process of absorption and the arrow $`𝒰𝒱𝒰`$ describes the process of emission. Let us denote by the collection of all formal linear spans of primary operators, i.e. $`\text{}:=\{𝒰=𝒰(x^i):i=1,\mathrm{},N\}`$. The collection of complex conjugated spaces is denoted by $`\text{}^{}`$. We assume that an arbitrary sequence consists of the unit object $`I`$ or spaces from the collection or $`\text{}^{}`$ represents initial configuration of our system. These configurations can be transformed into some new ones by a set of certain transformations. These transformations represent certain physical processes like composition, emission, absorption, etc… It is obvious that these transformations can be coherent or not. Coherence for a set of transformations means path–independent construction of these transformations. Note that the coherence problem can be expressed graphically in terms of tangle tree operads . Our goal is the construction of a collection of transformations which transform in an unique way initial configurations into final ones – representing the result of interactions. We denote $``$ the generating set for these transformations. Let us consider some examples. Example 1. If $``$ contains only one operation, namely the product $`:\text{}\times \text{}\text{}\text{}`$, then starting from this product we can construct a set of multiproducts $`^n:\text{}^{\times n}\text{}^^n`$ such that $`^2`$ and every multiproduct $`^n`$ for $`n>2`$ can be calculated by an iteration procedure. Such procedure need not be unique. For instance for $`n=3`$ we obtain $`^3:=(\times id)`$, But we also obtain $`^3:=(id\times )`$. Hence for the uniqueness we need some additional assumptions like the associativity constraints, see the Appendix for more details. Example 2. For the (left) $``$–operation we use the standard relations $$𝒰^{}=𝒰,(𝒰𝒱)^{}=𝒱^{}𝒰^{}.$$ (1) In this case $`:=\{.\}`$. Example 3. We introduce a generating set $`g(\text{}):=\{g_𝒰:𝒰\text{}\}`$ of $`I`$–valued mappings $`g_𝒰:𝒰^{}𝒰I`$, where $$g_𝒰(x^ix^j)(x^i\|x^j):=i\|j$$ (2) for pairing $`g`$. The extension $`g_{𝒰𝒱}`$ of this pairing to the product $`𝒰𝒱`$ is a problem. We need here the following commutative diagram $$\begin{array}{ccc}& id_𝒱^{}g_𝒰id_𝒱& \\ 𝒱^{}𝒰^{}𝒰𝒱& & 𝒱^{}𝒱\\ & & \\ & & g_𝒱\\ & & \\ (𝒰𝒱)^{}(𝒰𝒱)& & I\\ & g_{𝒰𝒱}& \end{array}$$ (3) for the extension. We can introduce the set $`g(\text{}^{}):=\{g_𝒰:𝒰\text{}\}`$ of $`I`$–valued mappings $`g_𝒰^{}:𝒰𝒰^{}I`$ in a similar way. For the extension we use the diagram $$\begin{array}{ccc}& id_𝒰g_𝒱^{}id_𝒰^{}& \\ 𝒰𝒱𝒱^{}𝒰^{}& & 𝒰𝒰^{}\\ & & \\ & & g_𝒰^{}\\ & & \\ (𝒰𝒱)(𝒰𝒱)^{}& & I\\ & g_{(𝒰𝒱)^{}}& \end{array}$$ (4) In the way we can construct a category $`:=(\text{},)`$ whose objects are multiple products of object of and morphisms are obtained by iteration procedures applied to operations from the set $``$. The monoidal operation of $``$ can also be obtained by the proper iteration of the initial product $``$. Note that there is the uniqueness problem with the construction of the category $`=(\text{},)`$ over . One can construct many different categories $``$ for a given initial collection of spaces with different generating set $``$. We denote by $`\mathrm{𝐂𝐚𝐭}(\text{})`$ the class of all these categories. Let $`(\text{},)`$ and $`𝒩(\text{},^{})`$ be two such categories. An arbitrary functor $`:𝒩`$ which transform the set of operations $``$ into $`^{}`$ is said to be a generalized transmutation. In this case we say that the set $``$ is transmuted into $`^{}`$. The category $`𝒩`$ is then said to be functored over $``$ . If $`:=\{I,\}`$ and $`^{}:=\{I^{},\underset{¯}{}\}`$, then the corresponding generalized transmutation $`:𝒩`$ is a monoidal functor of categories $``$ and $`𝒩`$. This means that it is a triple $$:=\{,\phi _2,\phi _0\}:𝒩$$ which consists of a functor $`:𝒩`$, a natural isomorphism $$\phi :=\phi _{2,𝒰,𝒱}:𝒰\underset{¯}{}𝒱(𝒰𝒱),$$ and an isomorphism $`\phi _0:II=I^{}`$, such that the following diagrams $$\begin{array}{ccc}& \phi _2\underset{¯}{}id& \\ 𝒰\underset{¯}{}𝒱\underset{¯}{}𝒱& & (𝒰𝒱)\underset{¯}{}𝒲\\ & & \\ id\phi _2& & \phi _2\\ & & \\ 𝒰\underset{¯}{}(𝒱𝒲)& & (𝒰𝒱𝒲)\\ & \phi _2& \end{array}$$ (5) $$\begin{array}{ccc}& \phi _2& \\ \hfill I\underset{¯}{}𝒰& & (I𝒰)\\ & & \\ \hfill \phi _0\underset{¯}{}id& & (l_𝒰)\\ & & \\ \hfill 𝒰& & \end{array}$$ (6) $$\begin{array}{ccc}& \phi _2& \\ \hfill 𝒰\underset{¯}{}I& & (𝒰I)\hfill \\ & & \\ \hfill id\underset{¯}{}\phi _0& & (r_𝒰)\hfill \\ & & \\ \hfill 𝒰& & \end{array}$$ (7) are commutative. If $`\phi _2`$ and $`\phi _0`$ are identities, then $``$ is said to be strict. For the $``$–operation and pairing we have $$((𝒰))^{}=(𝒰^{}),g_𝒰=g_{(𝒰)}^{};$$ (8) respectively. A generalized transmutation $`:𝒩`$ is said to be cross symmetric if the following diagram $$\begin{array}{ccc}& \phi _2& \\ 𝒰^{}\underset{¯}{}𝒱& & (𝒰^{}𝒱)\\ & & \\ \mathrm{\Psi }^{}& & (\mathrm{\Psi })\\ & & \\ 𝒱\underset{¯}{}𝒰^{}& & (𝒱𝒰^{})\\ & \phi _2& \end{array}$$ (9) is commutative for every generating objects $`𝒰,𝒱`$ of $``$, where $`\mathrm{\Psi }`$ and $`\mathrm{\Psi }^{}`$ are the cross symmetries in $``$ and $`𝒩`$, respectively. In the case of a braid symmetries we have the following diagram $$\begin{array}{ccc}& \phi _2& \\ 𝒰\underset{¯}{}𝒱& & (𝒰𝒱)\\ & & \\ \mathrm{\Psi }^{}& & (\mathrm{\Psi })\\ & & \\ 𝒱\underset{¯}{}𝒰& & (𝒱𝒰)\\ & \phi _2& \end{array}$$ (10) Note that the functor $`:𝒩`$ generalizes the Majid concept of transmutation of braid statistics . If $``$ is the category of cobordisms of smooth manifolds and $`𝒩`$ is a category of vector spaces, that the generalized transmutation is known as the Topological Quantum Field Theory . Example 4. Let $`H`$ and $`H^{}`$ be two Hopf algebras. If $`h:HH^{}`$ is a Hopf algebra homomorphism, then we can introduce the transmutation $`:^H^H^{}`$ of categories of right comodules as follows. The functor $``$ acts as the identity functor on arbitrary comodule $`𝒰`$ but the coaction $`\rho _𝒰:𝒰𝒰H`$ transform into new one, namely into $`\rho _𝒰^{}:𝒰𝒰H^{}`$, where $$\rho _𝒰^{}:=(id_𝒰h)\rho _𝒰.$$ (11) For coquasitriangular Hopf algebras $`H`$ and $`H^{}`$ with coquasitriangular structures $`,:HHI`$ and $`,^{}:H^{}H^{}I`$, respectively, we obtain $$\begin{array}{c}\mathrm{\Psi }_{𝒰,𝒱}^{}(uv)=\mathrm{\Sigma }h(v_1),h(u_1)^{}v_0u_0,\end{array}$$ (12) where $`\rho (u)=\mathrm{\Sigma }u_0u_1𝒰H`$, $`\rho (v)=\mathrm{\Sigma }v_0v_1𝒱H`$ for every $`u𝒰,v𝒱`$, and $`k,l=h(k),h(l)^{}`$ for every $`k,lH`$. Example 5. Let $`H:=\text{}G`$ and $`H^{}:=\text{}G^{}`$ be group algebras, where $`G`$ and $`G^{}`$ are Abelian groups equipped with factors $`ϵ`$ and $`ϵ^{}`$, respectively. Then the transmutation $`:(𝖦,ϵ)(𝖦^{},ϵ^{})`$ is determined by a group homomorphism $`h:GG^{}`$ such that $$ϵ(\alpha ,\beta )=ϵ(h(\alpha ),h(\beta ))$$ (13) for $`\alpha ,\beta G`$. ## 3 Commutation relations Let us denote by $`\text{}(n)`$ the collection of all tensor products of the form $$\begin{array}{c}𝒰_{i_1}\mathrm{}𝒰_{i_n},\hfill \end{array}$$ (14) for all $`𝒰_{i_1},\mathrm{},𝒰_{i_n}\text{}`$. We introduce $`\text{}^{}(n)`$ in a similar way. We also introduce the collection $`(\text{}^{}\text{})(n)`$ of sequences of the form $$\begin{array}{c}𝒰_{j_n}^{}\mathrm{}𝒰_{j_1}^{}𝒰_{i_1}\mathrm{}𝒰_{i_n}.\hfill \end{array}$$ (15) The collection $`(\text{}\text{}^{})(n)`$ can be defined in an obvious way. We have the following examples. Example 4. Let us denote by $`:=(\text{},I,,,g)`$ the category $`:=(\text{},)`$, where $`:=\{I,,,g\}`$. We introduce two sets of transformations $`a^+:=\{a_𝒰^+:\text{}(n)\text{}(n+1)\}`$ and $`a^{}:=\{a_𝒰^{}^{}:\text{}(n)\text{}(n1)\}`$, where $$\begin{array}{c}a_𝒰^+(𝒰_{i_1}\mathrm{}𝒰_{i_n}):=𝒰𝒰_{i_1}\mathrm{}𝒰_{i_n},\end{array}$$ (16) and $$\begin{array}{c}a_𝒰^{}^{}(𝒰_{i_1}\mathrm{}𝒰_{i_n}):=g_𝒰(𝒰_{j_1}^{}𝒰_{i_1})\mathrm{}𝒰_{i_n}\end{array}$$ (17) where $$g_𝒰(𝒰_{j_1}^{}𝒰_{i_1})\{\begin{array}{cc}g_𝒰& \text{for}𝒰_{j_i}^{}𝒰^{},𝒰_{i_1}𝒰\hfill \\ 0& \text{otherwise}\hfill \end{array},$$ (18) for $`𝒰_{i_1},\mathrm{},𝒰_{i_n},𝒰`$ and $`𝒰_{j_1}^{},𝒰^{}\text{}`$. Here $`𝒰^{}`$ represents a quasihole and $`𝒰`$ – a quasiparticle. Two different objects $`𝒰,𝒱\text{}`$ represent (quasi-) particle states of two different sorts. There are no identical particles. It is easy to see that we have the following set of relations $$a_𝒰^{}^{}a_𝒰^+=g_𝒰\mathrm{𝟏},$$ (19) where $`𝒰\text{}`$. These relations are in fact the commutation relations for the system equipped with the infinite statistics . One can use the relation $$\widehat{a}_𝒰\widehat{a}_𝒰=\widehat{a}_{𝒰𝒱},$$ where $`\widehat{a}_𝒰`$ stands for $`a_𝒰^+`$ or $`a_𝒰^{}`$, for the extension of commutation relations corresponding for monoidal products of generating objects. These relations seems to be simple, but they lead to well–defined operator algebras . Observe that we have here elementary quantum processes of two sorts, namely creation and annihilation. Example 5. We assume that $``$ contains the cross symmetry $`\mathrm{\Psi }_{cross}`$ in addition to the previous example. The corresponding category is denoted by $`_{cross}:=(\text{},I,,,g,\mathrm{\Psi }_{cross})`$. We need here a collection $`\mathrm{\Psi }(\text{}):=\{\mathrm{\Psi }_{𝒰^{},𝒱}:𝒰^{},𝒱\text{}\}`$ as the initial data for the description of exchange processes of (quasi-) particles and (quasi-) holes, see the Appendix. We have here the following relations $$\begin{array}{c}b_𝒰^{}^{}b_𝒰^+b_𝒰^+b_𝒰^{}^{}\mathrm{\Psi }_{𝒰^{},𝒰}:=g_𝒰\mathrm{𝟏},\end{array}$$ (20) where $$\begin{array}{c}b_𝒰^+:=a_𝒰^+,\hfill \\ b_𝒰^{}^{}(𝒱_1𝒱_2):=\hfill \\ \left[(a^{}id_{𝒱_2})(id_{𝒱_1}a^{})(\mathrm{\Psi }_{𝒰^{},𝒱}id_{𝒱_2})\right](U^{}𝒱_1𝒱_2),\hfill \end{array}$$ (21) and $$a^{}(𝒰^{}𝒱):=a_𝒰^{}^{}(𝒱).$$ (22) In this way we have here a collection of elementary processes $``$ which contains creation, annihilation and exchange processes. Example 6. We replace the cross symmetry by the braid one. In this case we obtain the following relations $$\begin{array}{c}c_𝒰^{}^{}c_𝒰^+c_𝒰^+c_𝒰^{}^{}\mathrm{\Psi }_{𝒰^{},𝒰}:=g_𝒰\mathrm{𝟏},\end{array}$$ (23) and in addition $$\begin{array}{c}c_𝒰^{}^{}c_𝒱^+c_𝒱^+c_𝒰^{}^{}\mathrm{\Psi }_{𝒰^{},𝒱}=0,\hfill \\ c_𝒰^+c_𝒱^+c_𝒱^+c_𝒰^+\mathrm{\Psi }_{𝒰,𝒱}=0,\hfill \\ c_𝒰^{}^{}c_𝒱^{}^{}c_𝒱^{}^{}c_𝒰^{}^{}\mathrm{\Psi }_{𝒰^{},𝒱^{}}=0.\hfill \end{array}$$ (24) Note that for the braid symmetry there are additional elementary quantum processes, namely the exchange processes of identical particles on lattice in two dimensional case. Example 7. Let $`_{cross}:=(\text{},I,,,g,\mathrm{\Psi }_{cross})`$ be a category with commutation relation1s like in the Example 5. we denote by $`𝒩_{cross}:=𝒩(\text{},I,\underset{¯}{},,g^{},\mathrm{\Psi }_{cross}^{})`$ a second category with a new cross $`\mathrm{\Psi }_{cross}^{}`$ and pairing $`g^{}`$. One can define the following two sets $`c^+:=\{c_𝒰^+:\text{}(n)\text{}(n+1)\}`$ and $`c^{}:=\{c_𝒰^{}^{}:\text{}(n)\text{}(n1)\}`$ of operators in it. For a cross symmetric generalized transmutation $`:𝒩`$, we have the relation for these operators $$\begin{array}{c}c_{(𝒰^{})}^{}c_{(𝒰)}^+c_{(𝒰)}^+c_{(𝒰^{})}^{}\mathrm{\Psi }_{(𝒰^{}),(𝒰)}^{}:=g_{(𝒰)}^{}\mathrm{𝟏},\end{array}$$ (25) Note that the category $`𝒩`$ can be braided or symmetric. In these cases we obtain additional relations such as (25). It is obvious that the concept of category symmetries is related to the systems with generalized statistics . Note that the braid commutation relations, consistency conditions and corresponding Fock space representation with well–defined scalar product has been considered previously, see for instance. Some interesting examples of related formalism has been studied previously by Fiore . Observe that the above concept of category symmetries can be futher developed in a few respects. One can consider the corresponding noncommmutative calculi, It should be interesting to study Hamiltonians in terms of described here creation and annihilation operators and study the concrete physical models. ## Appendix Let us briefly recall the fundamental concept of the category theory for the fixing of notation. For more details see the textbook of Mac Lane . A category $``$ contains a collection $`𝒪b()`$ of objects and a collection $`hom()`$ of arrows (morphisms). The collection $`hom()`$ is the union of mutually disjoint sets $`hom(𝒰,𝒱)`$ of arrows $`f:𝒰𝒱`$ from $`𝒰`$ to $`𝒱`$ defined for every pair of objects $`𝒰,𝒱𝒪b()`$. It may happen that for a pair $`𝒰,𝒱𝒪b()`$ the set $`hom(𝒰,𝒱)`$ is empty. The associative composition of morphisms is also defined. A functor $`:𝒩`$ of the category $``$ into the category $`𝒩`$ is a map which sends objects of $``$ into objects of $`𝒩`$ and morphisms of $``$ into morphisms of $`𝒩`$ such that $`(fg)=(f)(g)`$ for every morphisms $`f:𝒱𝒲`$ and $`g:𝒰𝒱`$ of $``$. The generalization to multifunctors is obvious. One can consider an arbitrary object of a category as an example of constant functor. For instance an $`n`$–ary functor $`:^{\times n}𝒩`$ sends an $`n`$–tuple of objects of $``$ into an object of $`𝒩`$. The corresponding condition for morphisms is evident. In this paper we restrict our attention for a description how functors act on objects, we omit the action on morphisms for simplicity. The reader can complete our description. Now we recall the concept of natural transformations. Let $``$ and $`𝒢`$ be two functors of the category $``$ into the category $`𝒩`$. A natural transformation $`s:𝒢`$ of $``$ into $`𝒢`$ is a collection of morphisms $`s=\{s_𝒰:(𝒰)𝒢(𝒰),𝒰𝒪b()\}`$ such that $$s_𝒱(f)=𝒢(f)s_𝒰$$ (26) for every morphism $`f:𝒰𝒱`$ of $``$. The set of all natural transformations of $``$ into $`𝒢`$ is denoted by $`𝒩at(,𝒢)`$. It is easy to see that the composition $`tf`$ of natural transformation $`s`$ of $``$ into $`𝒢`$ and $`t`$ of $`𝒢`$ into $``$ is a natural transformation of $``$ into $``$. If $`𝒢`$, then we say that the natural transformation $`s:𝒢`$ is a natural transformation of $``$ into itself. Now let us briefly explain the notions of monoidal categories adopted for our goal. A monoidal category $`(,I)`$ is in fact a category $``$ equipped with a monoidal operation (a bifunctor) $`:\times `$, a unit object $`I`$, and collections of natural isomorphisms: (i) an associativity constraint $`\psi =\{\psi _{𝒰,𝒱,𝒲}:(𝒰𝒱)𝒲𝒰(𝒱𝒲)`$}, (ii) a left unity constraint $`l=\{l_𝒰:I𝒰𝒰\}`$ (iii) and a right unity constraint $`r=\{r_𝒰:𝒰k𝒰\}`$ such that the following diagrams $$\begin{array}{ccc}& (𝒰𝒱)(𝒲𝒳)& \\ \psi _{𝒰𝒱,𝒲,𝒳}& & \psi _{𝒰,𝒱,𝒲𝒳}\\ & & \\ ((𝒰𝒱)𝒲)𝒳& & 𝒰(𝒱(𝒲𝒳))\\ & & \\ \psi _{𝒰,𝒱,𝒲}id& & id\psi _{𝒱,𝒲,𝒳}\\ & & \\ (𝒰(𝒱𝒲))𝒳& & 𝒰((𝒱𝒲)𝒳)\\ & \psi _{𝒰,𝒱𝒲,𝒳}& \end{array}$$ (27) $$\begin{array}{ccc}& \psi _{𝒱,I,𝒲}& \\ (𝒱I)𝒲& & 𝒱(I𝒲)\\ & & \\ r_𝒱id& & idl_𝒲\\ & & \\ & 𝒱𝒲& \end{array}$$ (28) commute. It is interesting that in a monoidal category any diagram built from the constraints $`\psi ,l,r`$, and the identities by composing and tensoring, commutes. This is just the famous Mac Lane’s coherence theorem. A monoidal category $``$ is said to be strict, if all natural isomorphisms $`\psi _{𝒰,𝒱,𝒲},l_𝒰,r_𝒰`$ are identity. It is also interesting that every monoidal category is equivalent to certain strict one. This means that we can restrict our attention to strict monoidal categories. A (left) $``$-operation in a monoidal category $``$ is a transformation $`()^{}`$ of functor $``$ into the opposite functor $`^{op}`$ such that $$\begin{array}{cc}()^{}=id_{},& ()^{}=^{op}()^{}\end{array}$$ (29) where $`𝒰`$ and $`𝒱`$ are arbitrary objects of the category $``$. A (left) pairing $`g`$ in the category $``$ is a transformation of the functor $`()^{}`$ into $`I`$, where $`I`$ is a field satisfying some compatibility axioms, see . This means that $`g`$ is a set $`g\{g_𝒰\}`$ of $`I`$–valued mappings $$\begin{array}{c}g\{g_𝒰:𝒰^{}𝒰I,𝒰𝒪b()\}\end{array}$$ (30) Let $``$ be a monoidal category equipped with a (left) $``$-operation $`()^{}`$ and a (left) pairing $`g`$, then such category is said to be a category with (left) duality and it is denoted by $`=_{left}(,,I,,g)`$. One can introduce a (right) duality structure in the category $``$ in a similar way. Note that both dualities in $``$ the right and the left one are in general two independent structures. But it is possible to introduce an additional structure which making these two structures equivalent. Such equivalence can be established by the following set of natural isomorphisms $$\mathrm{\Psi }\{\mathrm{\Psi }_{𝒰^{},𝒱}:𝒰^{}𝒱𝒱𝒰^{}\}.$$ (31) where $$\begin{array}{c}\mathrm{\Psi }_{𝒰^{}𝒱^{},𝒲}=(\mathrm{\Psi }_{𝒰^{},𝒲}id_𝒱)(id_𝒰\mathrm{\Psi }_{𝒱^{},𝒲}),\hfill \\ \mathrm{\Psi }_{𝒰^{},𝒱𝒲}=(id_𝒱\mathrm{\Psi }_{𝒰^{},𝒲})(\mathrm{\Psi }_{𝒰^{},𝒱}id_𝒲),\hfill \end{array}$$ (32) for every objects $`𝒰,𝒱,𝒲`$ in $``$. These transformations are called a generalized cross symmetry, . We can identify the right and left duality in the category equipped with such generalized cross symmetry. The monoidal category equipped with such symmetry is denoted by $`=(,,I,,g,\mathrm{\Psi }_{cross})`$. Note that the generalized cross symmetry is not a braid symmetry in general. For the braid symmetry in the category with duality we need additional transformations like $$\mathrm{\Psi }\{\mathrm{\Psi }_{𝒰,𝒱}:𝒰𝒱𝒱𝒰\}$$ (33) for arbitrary objects $`𝒰,𝒱`$ in $``$ and $$\mathrm{\Psi }\{\mathrm{\Psi }_{𝒰^{},𝒱^{}}:𝒰^{}𝒱^{}𝒱^{}𝒰^{}\}$$ (34) for objects $`𝒰^{},𝒱^{}`$ in $``$. We need also some new commutative diagrams for all these transformations and pairings. In fact a family of natural isomorphisms $$\begin{array}{c}\mathrm{\Psi }\{\mathrm{\Psi }_{𝒰𝒲}:𝒰𝒲𝒲𝒰\}\end{array}$$ (35) such that we have the following relations $$\begin{array}{c}\mathrm{\Psi }_{𝒰𝒱,𝒲}=(\mathrm{\Psi }_{𝒰,𝒲}id_𝒱)(id_𝒰\mathrm{\Psi }_{𝒱,𝒲}),\hfill \\ \mathrm{\Psi }_{𝒰,𝒱𝒲}=(id_𝒱\mathrm{\Psi }_{𝒰,𝒲})(\mathrm{\Psi }_{𝒰,𝒱}id_𝒲),\hfill \end{array}$$ (36) is said to be a braiding or a braid symmetry on $``$. The monoidal category with unique duality and braid symmetry is said to be rigid . This category is denoted by $`=(,,I,,g.\mathrm{\Psi }_{braid})`$. If in addition we have the relation $$\begin{array}{c}\mathrm{\Psi }_{𝒰,𝒱}^2=id_{𝒰𝒱},\end{array}$$ (37) for every objects $`𝒰,𝒱`$, then the set $`S:=\{\mathrm{\Psi }_{𝒰,𝒱}\}`$ is said to be a (vector) symmetry or tensor symmetry and the corresponding category $``$ is called a symmetric monoidal or tensor category, see . Let us consider some examples of monoidal categories which can be useful for the study of category symmetry. The most simple example of a monoidal category is provided by the category $`𝒱ect(k)`$ of vector spaces over a field $`k`$. The monoidal operation in this category is defined by the usual tensor product of vector spaces. Another example is given by the category $`𝒱ect_G(k)`$ of $`G`$–graded vector spaces, where $`G`$ is a grading group. In the supersymmetry the grading group is the group of integer $`Z_2`$. For anyons we have $`GZ_n`$, where $`n>2`$, .There is a category $`{}_{}{}^{}`$ of all left $``$-modules, where $``$ is an unital and associative algebra. Observe that the usual tensor product $`𝒰𝒱`$ of two left $``$–modules $`𝒰`$ and $`𝒱`$ is not a left $``$-module but a left $``$-module! Hence this category is not a monoidal category. But it is easy to see that in the particular case when $``$ is a bialgebra, i.e. we have a comultiplication $`\mathrm{}:`$ in $``$, the category $`{}_{}{}^{}`$ is monoidal. For instance there is the category $`_G`$ of finite dimensional representations of compact matrix quantum group $`G`$, . There is also a category of Hopf modules or crossed modules . Observe that there is also a category $`^{}`$ of right $``$–comodules, where $``$ is a Hopf algebra. The monoidal operation in $`^{}`$ is given as the following tensor product of $``$–comodules $$\rho _{𝒰𝒱}=(idm_{})(id\tau id)(\rho _𝒰\rho _𝒱),$$ (38) where $`\tau :𝒰𝒰`$ is the twist, $`m_{}:`$ is the multiplication in $`H`$. As an example for a category with duality we can give a category $`{}_{}{}^{}`$ of left $``$–modules. In this case the monoidal operation corresponds to the tensor product of representations, the $``$-operation corresponds to the contragradient representation and the generalized cross symmetry corresponds to the intertwiner between an arbitrary representation and its contragradient. Hence it is also called a statistics operator. Note that if there is a bialgebra $``$ such that the category $``$ is equivalent to the category of left modules (representations) over $``$, then the bialgebra is said to be a generalized symmetry or (bi-)algebra symmetry for the given physical system. One can describe states in the quantum Hall effect as a result of symmetry in such generalized sense . The symmetry algebra for Klein–Gordon equation on quantum Minkowski space is considered in . Note that the category of representations of the so–called weak Hopf algebra is rigid . Also the category of quantum compact matrix groups of Woronowicz is rigid . For a coquasitriangular Hopf algebra $`H`$ with a coquasitriangular structure $`,:HHI`$ we obtain the category $`^H`$ of right $`H`$-comodules which is also braided monoidal . The braid symmetry $`\mathrm{\Psi }\{\mathrm{\Psi }_{𝒰,𝒱}:𝒰𝒱𝒱𝒰;𝒰,𝒱Ob\}`$ in $``$ is defined by the equation $$\begin{array}{c}\mathrm{\Psi }_{𝒰,𝒱}(uv)=\mathrm{\Sigma }v_1,u_1v_0u_0,\end{array}$$ (39) where $`\rho (u)=\mathrm{\Sigma }u_0u_1𝒰H`$, and $`\rho (v)=\mathrm{\Sigma }v_0v_1𝒱H`$ for every $`u𝒰,v𝒱`$. Let $`G`$ be an arbitrary group, then the group algebra $`H:=\text{}G`$ is a Hopf algebra for which the comultiplication, the counit, and the antypode are given by the formulae $$\begin{array}{cccc}\mathrm{}(g):=gg,& \eta (g):=1,& S(g):=g^1& \text{for}gG.\end{array}$$ respectively. If $`H\text{}G`$, where $`G`$ is an Abelian group, then the coquasitriangular structure on $`H`$ is given as a bicharacter on $`G`$ . Note that for Abelian groups we use the additive notation. A mapping $`ϵ:G\times G\text{}\{0\}`$ is said to be a bicharacter on $`G`$ if and only if we have the following relations $$ϵ(\alpha ,\beta +\gamma )=ϵ(\alpha ,\beta )ϵ(\alpha ,\gamma ),ϵ(\alpha +\beta ,\gamma )=ϵ(\alpha ,\gamma )ϵ(\beta ,\gamma )$$ (40) for $`\alpha ,\beta ,\gamma G`$. If in addition $$ϵ(\alpha ,\beta )ϵ(\beta ,\alpha )=1,$$ (41) for $`\alpha ,\beta G`$, then $`ϵ`$ is said to be a normalized bicharacter or a commutation factor on $`G`$ . The category $`^H`$ of right comodules, where $`H:=\text{}G`$ for certain Abelian group $`G`$ and $`,ϵ(,)`$ is a bicharacter like above is denoted by $`(G,ϵ)`$. Note that if $`𝒰`$ is a $`H`$-comodule, where $`H=\text{}G`$, then $`𝒰`$ is a $`G`$-graded vector space, i.e $`𝒰=\underset{\alpha G}{}𝒰_\alpha `$. This means that a coaction of $`H:=\text{}G`$ on $`𝒰`$ is equivalent to $`G`$-gradation of $`𝒰`$. Acknowledgments The author would like to thank to Prof, W. Rühl for kind invitation to Universität Kaiserslautern and discussion, to M. Greulach for any other help. The work is partially sponsored by DAAD, Bonn, Germany, and by Polish Committee for Scientific Research (KBN) under Grant 2P03B130.12.
warning/0002/gr-qc0002049.html	ar5iv	text	# Proof of a generalized Geroch conjecture for the hyperbolic Ernst equation ## 1 A generalized Geroch conjecture In terms of Weyl canonical coordinates $`(z,\rho )`$, the Ernst equation of general relativity can be expressed in the form $$(\mathrm{R}e)\left\{\frac{^2}{z^2}\pm \frac{1}{\rho }\frac{}{\rho }\left(\rho \frac{}{\rho }\right)\right\}=\left(\frac{}{z}\right)^2\pm \left(\frac{}{\rho }\right)^2,$$ (1.1) where the upper signs correspond to the elliptic equation associated with stationary axisymmetric (spinning body) gravitational fields and the lower signs correspond to the hyperbolic equation associated with colliding gravitational plane wave pairs and cylindrical gravitational waves.<sup>1</sup><sup>1</sup>1In the latter case, one of the Weyl coordinates has the character of a time coordinate. In practice a notation more appropriate for the physical problem being treated would be in order. In 1972 R. Geroch asserted a conjecture<sup>2</sup><sup>2</sup>2R. Geroch, J. Math. Phys. 13, 394-404 (1972). concerning the solution manifold of the elliptic Ernst equation that was eventually proved<sup>3</sup><sup>3</sup>3I. Hauser and F. J. Ernst, A new proof of an old conjecture, in Gravitation and Geometry, Eds. Rindler and Trautman, Bibliopolis, Naples (1987). by the present authors, who used their own homogeneous Hilbert problem version of the Kinnersley–Chitre realization of the Geroch group. In 1986, at the suggestion of S. Chandrasekhar, we turned our attention from stationary axisymmetric fields to colliding gravitational plane wave pairs. While the Kinnersley–Chitre transformations could still be used to generate scores of exact analytic solutions of the hyperbolic Ernst equation, we were aware of the fact that there might exist a significantly larger group, for, whereas any $`𝐂^3`$ solution of the axis-accessible elliptic Ernst equation can be shown to be automatically an analytic solution, a solution of the hyperbolic Ernst equation can be even $`𝐂^{\mathrm{}}`$ without being analytic.<sup>4</sup><sup>4</sup>4Even the elliptic equation admits a larger group if solutions are considered that are everywhere axis-inaccessible. Clearly, one should not expect a non-analytic solution of the hyperbolic Ernst equation to be related to Minkowski space by a K–C transformation, for these transformations preserve analyticity. ### A. Linear systems for the Ernst equation Any discussion of the Geroch group or its extensions requires a knowledge of at least one linear system<sup>5</sup><sup>5</sup>5Such linear systems have been found by many authors, including Chinea, Harrison, Kinnersley and Chitre, Maison, Neugebauer and Papanicolaou. A more complicated type of linear system was found by Belinskii and Zakharov. $$dF(𝐱,\tau )=\mathrm{\Gamma }(𝐱,\tau )F(𝐱,\tau )$$ (1A.1) for the Ernst equation. Here $`𝐱`$ is shorthand for the nonignorable spacetime coordinates (e.g., $`z`$ and $`\rho `$), $`\tau `$ is a spacetime-independent complex-valued parameter, and the $`1`$-form $`2\times 2`$ matrix $`\mathrm{\Gamma }(𝐱,\tau )`$ satisfies the integrability condition $$d\mathrm{\Gamma }(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )=0$$ (1A.2) if and only if the Ernst equation is satisfied. The symbol $`\mathrm{\Gamma }(𝐱,\tau )`$ was chosen because of the resemblance of the last equation to a zero-curvature condition for a connection $`1`$-form. If there exists one such $`\mathrm{\Gamma }(𝐱,\tau )`$ for the Ernst equation, then there are infinitely many, for if $$\mathrm{\Gamma }^{}(𝐱,\tau ):=p(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )p(𝐱,\tau )^1+dp(𝐱,\tau )p(𝐱,\tau )^1,$$ (1A.3) where $`p(𝐱,\tau )`$ is an invertible matrix, then $$d\mathrm{\Gamma }^{}(𝐱,\tau )\mathrm{\Gamma }^{}(𝐱,\tau )\mathrm{\Gamma }^{}(𝐱,\tau )=p(𝐱,\tau )\left\{d\mathrm{\Gamma }(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )\right\}p(𝐱,\tau )^1.$$ (1A.4) This transformation is nothing but a gauge transformation, the analog of the effect that a mere change of basis has upon a connection $`1`$-form. Under such a gauge transformation, the matrix $`F(𝐱,\tau )`$ transforms into the matrix $$F^{}(𝐱,\tau ):=p(𝐱,\tau )F(𝐱,\tau ).$$ (1A.5) While, in one sense, the various possible representations of the linear system may be regarded as equivalent, in another sense they may be quite different, with the matrices $`F(𝐱,\tau )`$ and $`F^{}(𝐱,\tau )`$ possibly having very different domains in the space $`R^2\times C`$, as well as different continuity and/or differentiability properties. Often one representation is more useful for one part of the analysis, while another representation is more useful for another part. Different formalisms may also differ with respect to the number of columns that the matrix $`F`$ has. Here we shall follow an approach that we described long ago that effectively sidesteps the question of number of columns by introducing an auxilliary $`2\times 2`$ matrix potential $`(𝐱,\tau )`$ such that $$F(𝐱,\tau )=(𝐱,\tau )F(𝐱_0,\tau ),$$ (1A.6) $$d(𝐱,\tau )=\mathrm{\Gamma }(𝐱,\tau )(𝐱,\tau )$$ (1A.7) and $$(𝐱_0,\tau )=I,$$ (1A.8) where $`I`$ is a unit matrix, and $`𝐱_0`$ is a selected spacetime point within the domain of $`(𝐱)`$. Clearly, under a gauge transformation (1A.5), $`(𝐱,\tau )`$ transforms into $$^{}(𝐱,\tau ):=p(𝐱,\tau )(𝐱,\tau )p(𝐱_0,\tau )^1.$$ (1A.9) One of the simplest formulations of the linear system is that of G. Neugebauer,<sup>6</sup><sup>6</sup>6G. Neugebauer, Bäcklund transformations of axially symmetric stationary gravitational fields, Phys. Lett. A 12, L67 (1979). in which $`\mathrm{\Gamma }(𝐱,\tau )=\mathrm{\Gamma }_N(𝐱,\tau )`$, where $$\mathrm{\Gamma }_N(𝐱,\tau ):=\left(\frac{\tau z\pm \rho }{\tau z\rho }\right)^{\frac{1}{2}}\left(\begin{array}{cc}0& \frac{d(𝐱)}{2f(𝐱)}\\ \frac{d^{}(𝐱)}{2f(𝐱)}& 0\end{array}\right)+\left(\begin{array}{cc}\frac{d(𝐱)}{2f(𝐱)}& 0\\ 0& \frac{d^{}(𝐱)}{2f(𝐱)}\end{array}\right),$$ (1A.10) where $``$ is a $`2`$-dimensional duality operator such that $$d\rho =\pm dz,dz=d\rho ,$$ (1A.11) the upper signs applying in the stationary axisymmetric (elliptic) case, and the lower signs applying in the gravitational wave (hyperbolic) case. Here $`\mathrm{\Gamma }(𝐱,\tau )`$ is expressed directly in terms of the Ernst potential $`(𝐱)`$ and its complex conjugate, with $`f(𝐱):=\mathrm{R}e(𝐱)`$. Using these notations, the Ernst equation (1.1) can be expressed as $$(\mathrm{R}e)d(\rho d)=\rho dd.$$ (1A.12) A slightly different linear system that is due to the authors and is more suited to our purpose employs $`\mathrm{\Gamma }=\mathrm{\Gamma }_{HE}`$, where $$\mathrm{\Gamma }_{HE}(𝐱,\tau ):=\left(\frac{\tau z\pm \rho }{\tau z\rho }\right)^{\frac{1}{2}}\left(\frac{Idf(𝐱)Jd\chi (𝐱)}{2f(𝐱)}\right)\sigma _3J\frac{d\chi (𝐱)}{2f(𝐱)}$$ (1A.13) and $$\chi :=\mathrm{I}m,J:=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _3:=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$ (1A.14) The 1-form $`\mathrm{\Gamma }_{HE}`$ can be obtained from $`\mathrm{\Gamma }_N`$ by the gauge transformation (1A.3) corresponding to $`p=p_{NHE}`$, where $$p_{NHE}(𝐱,\tau )=\frac{1}{2\sqrt{\|f(𝐱)\|}}\left(\begin{array}{cc}1\pm i& 1i\\ 1i& 1i\end{array}\right).$$ (1A.15) On the other hand, the Kinnersley–Chitre formulation of the linear system<sup>7</sup><sup>7</sup>7W. Kinnersley and D. M. Chitre, Symmetries of the stationary Einstein-Maxwell field equations, III, J. Math. Phys. 19, 1926–1931 (1978). corresponds to the choice $`\mathrm{\Gamma }(𝐱,\tau )=\mathrm{\Gamma }_{KC}(𝐱,\tau )`$, where $$\mathrm{\Gamma }_{KC}(𝐱,\tau ):=\frac{1}{2}\mathrm{\Lambda }(𝐱,\tau )^1dH(𝐱)\mathrm{\Omega },\mathrm{\Omega }:=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),$$ (1A.16) with $$\mathrm{\Lambda }(𝐱,\tau ):=\tau (z\pm \rho )$$ (1A.17) and $`H(𝐱)`$ a $`2\times 2`$ matrix generalization of the Ernst potential $`(𝐱)`$ that can be introduced in the following manner. It is well-known that any vacuum spacetime possessing two commuting Killing vector fields can be described in terms of a $`2\times 2`$ real symmetric matrix $`h(𝐱)`$ (a $`2\times 2`$ block of the metric tensor) that depends exclusively on the nonignorable coordinates, and that this matrix satisfies the equation $$d[\rho dhh^1]=d[\rho h^1dh]=0,$$ (1A.18) where $$\rho :=\sqrt{\|deth\|}.$$ (1A.19) Equation (1A.18) can be used to justify the introduction of a complex $`H`$-potential that satisfies the equations $$\rho d(\mathrm{I}mH)=ih\mathrm{\Omega }dh\text{ and }\mathrm{R}eH=h,$$ (1A.20) or, equivalently, $$2(z\pm \rho )dH=(H+H^{})\mathrm{\Omega }dH,$$ (1A.21) where $$HH^T=2z\mathrm{\Omega }\text{ and }\mathrm{R}eH=h.$$ (1A.22) Then it is not difficult to establish that $`\mathrm{\Gamma }(𝐱,\tau )`$ as given by Eq. (1A.16) satisfies the zero-curvature condition (1A.2) if and only if $`:=H_{22}`$ satisfies the Ernst equation. The reader can verify that the K–C connection (1A.16) is related to the H–E connection (1A.13) by $$\mathrm{\Gamma }_{KC}(𝐱,\tau ):=p(𝐱,\tau )\mathrm{\Gamma }_{HE}(𝐱,\tau )p(𝐱,\tau )^1+dp(𝐱,\tau )p(𝐱,\tau )^1,$$ (1A.23) and $`_{KC}(𝐱,\tau )`$ is related to $`_{HE}(𝐱,\tau )`$ by $$_{KC}(𝐱,\tau )=p(𝐱,\tau )_{HE}(𝐱,\tau )p(𝐱_0,\tau )^1,$$ (1A.24) where $$p(𝐱,\tau )=\frac{1}{\sqrt{\|h_{22}(𝐱)\|}}\left(\begin{array}{cc}1& h_{12}(𝐱)\\ 0& \|h_{22}(𝐱)\|\end{array}\right)P^M(𝐱,\tau ),$$ (1A.25) $$P^M(𝐱,\tau ):=\left(\begin{array}{cc}1& \pm i(\tau z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& \mu (𝐱,\tau )^1\end{array}\right)\frac{1}{\sqrt{2}}(\sigma _3\sigma _2)$$ (1A.26) and $$\mu (𝐱,\tau ):=\sqrt{(\tau z)^2\pm \rho ^2},\underset{\tau \mathrm{}}{lim}\frac{\mu (𝐱,\tau )}{\tau }:=1.$$ (1A.27) Note that, for fixed $`𝐱`$, $`\mu (𝐱,\tau )`$ is a holomorphic function of $`\tau `$ throughout a cut complex plane. It has branch points of index $`1/2`$ at the zeroes of $`\mu (𝐱,\tau )`$, which are at the end points of the branch cut, and a simple pole at $`\tau =\mathrm{}`$. We shall assume that this brief review of the three formulations of the linear system for the Ernst equation and the relationships among these formulations will suffice. In the rest of this paper we shall suppress the subscript $`KC`$ on $`\mathrm{\Gamma }_{KC}(𝐱,\tau )`$ and $`_{KC}(𝐱,\tau )`$ as we proceed to discuss how a group $`𝒦`$ such as the Geroch group can be described in terms of its action upon the potentials $`_{KC}(𝐱,\tau )`$ associated with the spacetimes in question. ### B. The set $`𝒮_{}`$ of Kinnersley–Chitre $``$-potentials In order to discuss in a meaningful way the action of the group $`𝒦`$ upon the potentials $`(𝐱,\tau )`$, we must first identify the set $`𝒮_{}`$ of $``$-potentials being considered, and this requires, in particular, the specification of the domain of $`(𝐱,\tau )`$. This can best be done by first specifying the domain of $`H(𝐱)`$ \[and $`(𝐱)`$\], and then choosing the gauge of $`(𝐱,\tau )`$ so as to minimize its singularities in the complex $`\tau `$-plane. Throughout the rest of this paper we shall be concerned exclusively with the hyperbolic Ernst equation, where we find it convenient to introduce null coordinates $`r:=z\rho `$ and $`s:=z+\rho `$ and to adopt the $``$-potential domain (see Fig. 1) $$D:=\mathrm{d}om:=\{(r,s):r_1<r<r_2,s_2<s<s_1,r<s\}.$$ (1B.1) It is to be understood that $`r_1`$ may be $`\mathrm{}`$ and/or $`s_1`$ may be $`+\mathrm{}`$. Moreover, we restrict attention to domains $`D`$ such that $`r_1<s_2`$ and $`r_2<s_1`$; i.e., $`\rho >0`$ at both the lower left vertex $`(r_1,s_2)`$ and the upper right vertex $`(r_2,s_1)`$, while $`\rho `$ may be greater than, less than or equal to zero at the lower right vertex $`𝐱_2:=(r_2,s_2)`$. Finally, we select one point $`𝐱_0:=(r_0,s_0)D`$ such that the null line segments $`\{(r,s_0):r_1<r<r_2\}`$ and $`\{(r_0,s):s_2<s<s_1\}`$ lie entirely within $`D`$; and at this point we assign the Minkowski space value $`(𝐱_0)=1`$ to the complex $``$-potential.<sup>8</sup><sup>8</sup>8We have also considered more general domains and a more general choice for $`𝐱_0`$, but to include discussion of these extensions here would unnecessarily complicate our exposition. It is our intention to solve an initial value problem in which $`(𝐱)`$ is determined throughout $`D`$ from its values specified on the two null line segments through the point $`𝐱_0`$. For a given choice of the triple $`(𝐱_0,𝐱_1,𝐱_2)`$, we shall define $`𝒮_{}`$ $`:=`$ the set of all complex-valued functions $``$ such that $`\mathrm{d}om=D`$, the derivatives $`_r(𝐱)`$, $`_s(𝐱)`$ and $`_{rs}(𝐱)`$ exist and are continuous at all $`𝐱D`$, $`f:=\mathrm{R}e>0`$ and $``$ satisfies Eq. (1A.12) thoughout $`D`$, and $`(𝐱_0)=1`$. The metric components $`h_{ab}`$ corresponding to each given $`𝒮_{}`$ are defined by $`h_{22}:=f`$, $`d\omega :=\rho f^2d\chi `$ such that $`\omega (𝐱_0):=0`$, $`h_{12}:=\omega h_{22}`$ and $`h_{11}:=[(h_{12})^2+\rho ^2]/h_{22}`$. Naturally, we shall let $`\mathrm{d}omH=D`$ and assign the value $$H(𝐱_0)=H^M(𝐱_0),$$ (1B.3) where $`H^M`$ is the Minkowski space $`H`$-potential with values $$H^M(𝐱)=\left(\begin{array}{cc}\rho ^2& 0\\ 2iz& 1\end{array}\right).$$ (1B.4) For a given choice of the triple $`(𝐱_0,𝐱_1,𝐱_2)`$, we shall define $`𝒮_H`$ $`:=`$ the set of all complex-valued $`2\times 2`$ matrix functions $`H`$ with $`\mathrm{d}omH:=D`$ such that there exists $`𝒮_{}`$ for which $`\mathrm{R}eH=h`$, $`d(\mathrm{I}mH)`$ exists and satisfies $`\rho d(\mathrm{I}mH)=ih\mathrm{\Omega }dh\text{ and the gauge condition (}\text{1B.3}\text{) holds.}`$ Let $`^{(3)}(𝐱)`$ denote the open interval with end points $`r,r_0`$ and $`^{(4)}(𝐱)`$ denote the open interval with end points $`s,s_0`$, and let $`\overline{}^{(3)}(𝐱)`$ and $`\overline{}^{(4)}(𝐱)`$ denote, respectively, the closures of these two intervals. Furthermore, let $`(𝐱)`$ $`:=`$ $`^{(3)}(𝐱)^{(4)}(𝐱),\text{ and}`$ (1B.6) $`\overline{}(𝐱)`$ $`:=`$ $`\overline{}^{(3)}(𝐱)\overline{}^{(4)}(𝐱).`$ (1B.7) Note that $`\overline{}^{(3)}(𝐱)`$ is empty if $`r=r_0`$ and $`\overline{}^{(4)}(𝐱)`$ is empty if $`s=s_0`$. When neither $`r=r_0`$ nor $`s=s_0`$, the set $`\overline{}(𝐱)`$ comprises two disjoint closed sets (for $`𝐱_0`$ chosen as indicated earlier). The gauge of the $``$-potential can be chosen so that $$\mathrm{d}om:=\{(𝐱,\tau ):𝐱D,\tau C\overline{}(𝐱)\}.$$ (1B.8) For a given choice of the triple $`(𝐱_0,𝐱_1,𝐱_2)`$, we shall define $`𝒮_{}`$ $`:=`$ the set of all complex-valued $`2\times 2`$ matrix functions $``$ with domain (1B.8) such that there exists $`H𝒮_H`$ such that, for all $`𝐱D`$ and $`\tau [C\overline{}(𝐱)]\{r_0,s_0\}`$, $`d(𝐱,\tau )`$ exists and Eq. (1A.7) holds, subject to the condition (1A.8), and, for each $`(r,s)D`$, $`((r,s_0),\tau )`$ and $`((r_0,s),\tau )`$ are continuous functions of $`\tau `$ at $`\tau =s_0`$ and at $`\tau =r_0`$, respectively. Remember that at $`𝐱=𝐱_0`$, $`(𝐱,\tau )`$ reduces to the $`2\times 2`$ unit matrix. With these definitions one can establish the properties enumerated in the following theorem, the proof of which is (except for conventions and notations and the choice of the domain $`D`$) essentially the same as that given in two earlier papers<sup>9</sup><sup>9</sup>9I. Hauser and F. J. Ernst, Initial value problem for colliding gravitational plane waves-III/IV, J. Math. Phys. 31, 871–881 (1990), 32, 198–209 (1991). In these papers we used ‘$`P`$’ in place of ‘$``$’. on the IVP (initial value problem) for colliding gravitional plane wave pairs by the present authors. The complex-valued functions $`^{(3)}`$ and $`^{(4)}`$ with respective domains $`^{(3)}:=\{r:r_1<r<r_2\}`$ and $`^{(4)}:=\{s:s_2<s<s_1\}`$ serve as initial value data for the $``$-potential on the null line segments through the point $`𝐱_0`$. ###### THEOREM 1 (Initial Value Problem) * For each $`H𝒮_H`$, the corresponding $`𝒮_{}`$ exists and is unique; and, for each $`𝐱D`$, $`(𝐱,\tau )`$ is a holomorphic function of $`\tau `$ throughout $`C\overline{}(𝐱)`$ and, in at least one neighborhood of $`\tau =\mathrm{}`$, $$(𝐱,\tau )=I+(2\tau )^1\left[H(𝐱)H(𝐱_0)\right]\mathrm{\Omega }+O(\tau ^2).$$ (1B.10) * For each $`𝒮_{}`$, there is only one $`H𝒮_H`$ for which $`d(𝐱,\tau )=\mathrm{\Gamma }(𝐱,\tau )(𝐱,\tau )`$. * With the understanding that $$\mathrm{d}om\nu :=\{(𝐱,\tau ):𝐱D\text{ and }\tau C\overline{}(𝐱)\}$$ (1B.11) and that $`\nu (𝐱,\mathrm{})=1`$, we have $$det(𝐱,\tau )=\nu (𝐱,\tau ):=\frac{\mu (𝐱_0,\tau )}{\mu (𝐱,\tau )}=\left(\frac{\tau r_0}{\tau r}\right)^{1/2}\left(\frac{\tau s_0}{\tau s}\right)^{1/2}.$$ (1B.12) * The member of $`𝒮_{}`$ that corresponds to $`^M`$ is given by $$^M(𝐱,\tau )=\left(\begin{array}{cc}1& i(\tau z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& \nu (𝐱,\tau )\end{array}\right)\left(\begin{array}{cc}1& i(\tau z_0)\\ 0& 1\end{array}\right).$$ (1B.13) * For each $`𝒮_{}`$, there is exactly one $`H𝒮_H`$ such that $`=H_{22}`$. * If, for each $`i\{3,4\}`$, $`^{(i)}`$ is $`𝐂^{n_i}`$ ($`n_i1`$), then, for all $`0k<n_3`$ and $`0mn_4`$, the partial derivatives $`^{k+m}H(𝐱)/r^ks^m`$ exist and are continuous throughout $`D`$. If, for each $`i\{3,4\}`$, $`^{(i)}`$ is analytic, then $`H`$ is analytic. * For each choice of complex valued functions $`^{(3)}`$ and $`^{(4)}`$ for which (for $`i\{3,4\}`$) $`\mathrm{d}om^{(i)}=^{(i)}`$, $`^{(i)}`$ is $`𝐂^1`$, $`f^{(i)}:=\mathrm{R}e^{(i)}<0`$ throughout $`^{(i)}`$, and $`^{(3)}(r_0)=1=^{(4)}(s_0)`$, there exists exactly one $`𝒮_{}`$ such that $$^{(3)}(r)=(r,s_0)\text{ and }^{(4)}(s)=(r_0,s)$$ for all $`r^{(3)}`$ and $`s^{(4)}`$, respectively. ### C. Homogeneous Hilbert problem The HHP that we developed for effecting K–C transformations<sup>10</sup><sup>10</sup>10I. Hauser and F. J. Ernst, A homogeneous Hilbert problem for the Kinnersley–Chitre transformations, J. Math. Phys. 21, 1126-1140 (1980). (adapted to the hyperbolic case) involved a closed contour in the complex $`\tau `$-plane surrounding the arcs that comprise $`\overline{}(𝐱)`$. This was fine as long as we were dealing with the analytic case, but now we must instead formulate an HHP on those arcs themselves, and this will involve the limiting values of $`(𝐱,\tau )`$ as $`\tau `$ approaches points on those arcs. What we discovered concerning these limiting values is contained in the following theorems, the proofs of which are based upon a classic method of reducing the solving of a total differential equation to the solving of a pair of ordinary linear differential equations along characteristic lines in $`D`$. The Picard method of successive approximations and certain well known theorems of infinite sequences of functions are used to demonstrate existence, continuity and differentiability properties of the solution.<sup>11</sup><sup>11</sup>11For each $`\sigma R^1`$ and for fixed $`𝐱D`$, the limits of $`_{HE}(𝐱,\sigma \pm \zeta )`$ as $`\zeta 0(\mathrm{I}m\zeta >0)`$ exist. Moreover, $`_{HE}(𝐱,\tau ^{})=_{HE}(𝐱,\tau )^{}`$ and $`det(𝐱,\tau )=1`$. For these reasons, we found it convenient to use the H–E representation of the linear system in developing this proof, translating the results into corresponding results for the K–C representation. ###### THEOREM 2 (Limits of $``$) * For each $`𝐱D`$ and and $`\sigma (𝐱)`$ the limits $`^\pm (𝐱,\sigma ):=lim_{\zeta 0}(𝐱,\sigma \pm \zeta )(\mathrm{I}m\zeta >0)`$ exist. * Further, let $`\alpha `$ and $`\beta `$ be points of $`\overline{}(𝐱)`$ such that $`\alpha \{r_0,s_0\}`$ and $`\beta \{r,s\}`$, and let $`\tau C\overline{}(𝐱)`$. Then the following limits all exist and are equal as indicated: $`\underset{\sigma \alpha }{lim}^\pm (𝐱,\sigma )`$ $`=`$ $`\underset{\tau \alpha }{lim}(𝐱,\tau ),`$ (1C.1) $`\underset{\sigma \beta }{lim}[^\pm (𝐱,\sigma )^1]`$ $`=`$ $`\underset{\tau \beta }{lim}[(𝐱,\tau )^1].`$ (1C.2) We shall employ $`\mathrm{}`$ as a generic superscript that stands for $`n`$, $`n+`$, $`\mathrm{}`$ or ‘an’ (analytic). The symbols $`𝐂^n`$ and $`𝐂^{\mathrm{}}`$ are self explanatory. We shall say that $`f`$ is $`𝐂^{n+}`$ if its $`n`$th derivative $`D^nf`$ exists throughout $`\mathrm{d}omf`$ and $`D^nf`$ obeys a Hölder condition of arbitrary index on each closed subinterval of $`\mathrm{d}omf`$.<sup>12</sup><sup>12</sup>12The index may be different for different closed subintervals of $`\mathrm{d}omf`$. If $`f`$ is a real- or complex-valued function, the domain of which is a union of disjoint intervals of $`R^1`$, and $`[a,b]`$ is a given closed subinterval of $`\mathrm{d}omf`$, then $`f`$ is said to obey a Hölder condition of index $`0<\gamma 1`$ on $`[a,b]`$; i.e., to be $`H(\gamma )`$ on $`[a,b]`$, if there exists $`M(a,b,\gamma )>0`$ such that $`\|f(x^{})f(x)\|M(a,b,\gamma )\|x^{}x\|^\gamma `$ for all $`x,x^{}[a,b]`$. The same terminology is used if $`f(x)`$ is a matrix with real or complex elements, and $`\|f(x)\|`$ is its norm. Dfn. of the groups $`K^{\mathrm{}}`$ and $`K`$ * In order to describe our extensions $`𝒦^{\mathrm{}}`$ of the Geroch group, we shall introduce groups $`K^{\mathrm{}}`$ of $`2\times 2`$ matrix pairs; namely, the multiplicative groups of all ordered pairs $`𝐯=(v^{(3)},v^{(4)})`$ of $`2\times 2`$ matrix functions such that, for both $`i=3`$ and $`i=4`$, $$\mathrm{d}omv^{(i)}=^{(i)},detv^{(i)}=1,v^{(i)}\text{ is }𝐂^{\mathrm{}}$$ (1C.3) and the condition $$v^{(i)}(\sigma )^{}𝒜^M(𝐱_0,\sigma )v^{(i)}(\sigma )=𝒜^M(𝐱_0,\sigma )\text{ for all }\sigma ^{(i)}$$ (1C.4) holds, where $$𝒜^M(𝐱_0,\sigma ):=(\sigma z_0)\mathrm{\Omega }+\mathrm{\Omega }h^M(𝐱_0)\mathrm{\Omega },h^M(𝐱_0):=\left(\begin{array}{cc}\rho _0^2& 0\\ 0& 1\end{array}\right).$$ (1C.5) Moreover, the symbol $`K`$ will denote the multiplicative group of all ordered pairs $`𝐯=(v^{(3)},v^{(4)})`$ of $`2\times 2`$ matrix functions such that, for both $`i=3`$ and $`i=4`$, $$\mathrm{d}omv^{(i)}=^{(i)},detv^{(i)}=1,v^{(i)}\text{ is }H(1/2)\text{ on each closed subinterval of }^{(i)}$$ (1C.6) and the condition (1C.4) holds. End of Dfn. Dfn. of the HHP corresponding to $`(𝐯,_0)`$ * For each $`𝐯K^{\mathrm{}}`$ and $`_0𝒮_{}`$, the HHP corresponding to $`(𝐯,_0)`$ will mean the set of all functions $``$ \[which are not presumed to be members of $`𝒮_{}`$\] such that $`\mathrm{d}om=\{(𝐱,\tau ):𝐱D,\tau C\overline{}(𝐱)\}`$ and such that, for each $`𝐱D`$, the functions $`(𝐱)`$ whose domains are $`C\overline{}(𝐱)`$ and whose values are $`(𝐱,\tau )`$ is a solution of the HHP corresponding to $`(𝐯,_0,𝐱)`$, i.e., a member of the set of all $`2\times 2`$ matrix functions $`(x)`$ such that + $`(𝐱)`$ is holomorphic throughout $`\mathrm{d}om(𝐱):=C\overline{}(𝐱)`$, + $`(𝐱,\mathrm{})=I`$, + $`^\pm (𝐱)`$ exist, and $`Y^{(i)}(𝐱,\sigma )`$ $`:=`$ $`^+(𝐱,\sigma )v^{(i)}(\sigma )[_0^+(𝐱,\sigma )]^1`$ $`=`$ $`^{}(𝐱,\sigma )v^{(i)}(\sigma )[_0^{}(𝐱,\sigma )]^1`$ for each $`i\{3,4\}`$ and $`\sigma ^{(i)}(𝐱)`$, + $`(𝐱)`$ is bounded at $`𝐱_0`$ and $`\nu (𝐱)^1(𝐱)`$ is bounded at $`𝐱`$, and the function $`Y(𝐱)`$ whose domain is $`(𝐱)`$ and whose values are given by $`Y(𝐱,\sigma ):=Y^{(i)}(𝐱,\sigma )`$ for each $`\sigma ^{(i)}(𝐱)`$ is bounded at $`𝐱_0`$ and at $`𝐱`$. The members of the HHP corresponding to $`(𝐯,_0)`$ will be called its solutions. End of Dfn. Notes: * $`^+(𝐱)`$ and $`^{}(𝐱)`$ denote the functions that have the common domain $`(𝐱)`$ and the values ($`\mathrm{I}m\zeta >0`$) $$^\pm (𝐱,\sigma ):=\underset{\zeta 0}{lim}(𝐱,\sigma \pm \zeta ).$$ (1C.8) It is understood that $`^+(𝐱)`$ and $`^{}(𝐱)`$ exist if and only if the above limits exist for every $`\sigma (𝐱)`$. $`\nu ^+(𝐱)`$ and $`\nu ^{}(𝐱)`$ are similarly defined. * $`\nu (𝐱)`$ denotes the function whose domain is $`C\overline{}(𝐱)`$ and whose values $`\nu (𝐱,\tau )`$ are defined in Eq. (1B.12). * It is to be understood that $`(𝐱)`$, with domain $`C\overline{}(𝐱)`$, is bounded at $`𝐱_0`$ if there exists a neighborhood $`\mathrm{n}bd(𝐱_0)`$ of the set $`\{r_0,s_0\}`$ in the space $`C`$ such that $$\{(𝐱,\tau ):\tau \mathrm{n}bd(𝐱_0)\overline{}(𝐱)\}$$ (1C.9) is bounded. Likewise, $`(𝐱)`$ is said to be bounded at $`𝐱`$ if there exists a neighborhood $`\mathrm{n}bd(𝐱)`$ of the set $`\{r,s\}`$ in the space $`C`$ such that $$\{(𝐱,\tau ):\tau \mathrm{n}bd(𝐱)\overline{}(𝐱)\}$$ (1C.10) is bounded. * We say that $`Y(𝐱)`$, with domain $`(𝐱)`$, is bounded at $`𝐱_0`$ if there exists a neighborhood $`\mathrm{n}bd(𝐱_0)`$ of the set $`\{r_0,s_0\}`$ in the space $`R^1`$ such that $$\{Y(𝐱,\sigma ):\sigma \mathrm{n}bd(𝐱_0)(𝐱)\}$$ (1C.11) is bounded. Likewise, $`Y(𝐱)`$ is bounded at $`𝐱`$ if there exists a neighborhood $`\mathrm{n}bd(𝐱)`$ of the set $`\{r,s\}`$ in the space $`R^1`$ such that $$\{Y(𝐱,\sigma ):\sigma \mathrm{n}bd(𝐱)(𝐱)\}$$ (1C.12) is bounded. ###### THEOREM 3 (Properties of HHP solution) Suppose that $`𝐯K^{\mathrm{}}`$, $`_0𝒮_{}`$ and $`𝐱D`$ exist such that a solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,_0,𝐱)`$ exists. Then * $`^+(𝐱)`$, $`^{}(𝐱)`$ and $`Y(𝐱)`$ are continuous throughout $`(𝐱)`$, * $`^\pm (𝐱)`$ are bounded at $`𝐱_0`$, and $`[\nu ^\pm (𝐱)]^1^\pm (𝐱)`$ are bounded at $`𝐱`$, * $`det(𝐱)=\nu (𝐱),detY(𝐱)=1`$, * the solution $`(𝐱)`$ is unique, and * the solution of the HHP corresponding to $`(𝐯,_0,𝐱_0)`$ is given by $$(𝐱_0,\tau )=I$$ (1C.13) for all $`\tau C`$. Proofs: * The statement that $`^+(𝐱)`$ and $`^{}(𝐱)`$ are continuous is a direct consequence of a theorem by P. Painlevé which is stated and proved by N. I. Muskhelishvili.<sup>13</sup><sup>13</sup>13N. I. Muskhelishvili, Singular Integral Equations, Ch. 2, Sec. 14, pp. 33-34 (Dover, 1992). The continuity of $`Y(𝐱)`$ then follows from its definition by Eq. (3), the fact that $`v^{(i)}`$ is continuous and the fact that $`_0^+(𝐱)`$ and $`_0^{}(𝐱)`$ are continuous. End of proof. * From Eq. (3), $$^\pm (𝐱)=Y^{(i)}(𝐱)_0^\pm (𝐱)[v^{(i)}]^1$$ (1C.14) for each $`i\{3,4\}`$. The function $`Y(𝐱)`$ is bounded at $`𝐱`$ and at $`𝐱_0`$ according to condition (4) in the definition of the HHP, and $`v^{(i)}`$ and its inverse are continuous throughout $`^{(i)}`$. Finally, $`_0^\pm (𝐱)`$ is bounded at $`𝐱_0`$ and $`[\nu ^\pm (𝐱)]^1_0^\pm (𝐱)`$ is bounded at $`𝐱`$, so, from Eq. (1C.14), $`^\pm (𝐱)`$ is bounded at $`𝐱_0`$, and $`[\nu (𝐱)]^1^\pm (𝐱)`$ is bounded at $`𝐱`$. End of proof. * Conditions (1), (2), (3) and (4) of the definition of the HHP imply that $$Z_1(𝐱):=det(𝐱)/\nu (𝐱)\text{ is holomorphic throughout }C\overline{}(𝐱),$$ (1C.15) $$Z_1(𝐱,\mathrm{})=1,$$ (1C.16) $$\begin{array}{c}\text{the limits }Z_1^\pm (𝐱)\text{ exist and }\hfill \\ detY(𝐱,\sigma )=Z_1^+(𝐱,\sigma )=Z_1^{}(𝐱,\sigma )\text{ for all }\sigma (𝐱),\hfill \end{array}$$ (1C.17) $$\begin{array}{c}\nu (𝐱)Z_1(𝐱)\text{ is bounded at }𝐱_0\text{ and }\hfill \\ \nu (𝐱)^1Z_1(𝐱)\text{ is bounded at }𝐱,\hfill \end{array}$$ (1C.18) and $$detY(𝐱)=Z_1^\pm (𝐱)\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{.}$$ (1C.19) From the above statements (1C.15) and (1C.17) together with the theorem of Riemann<sup>14</sup><sup>14</sup>14See Sec. 24, Ch. 1, of A Course of Higher Mathematics, Vol. III, Part Two, by V. I. Smirnov (Addison-Wesley, 1964). on analytic continuation across an arc, $`Z_1(𝐱)`$ has a holomorphic extension to the domain $`C\{r,s,r_0,s_0\}`$; and, from the statements (1C.18) and (1C.19), together with the theorem of Riemann<sup>15</sup><sup>15</sup>15See Sec. 133 of Theory of Functions of a Complex Variable, Vol. 1, by C. Caratheodory, 2nd English edition (Chelsea Publishing Company, 1983). on isolated singularities of holomorphic functions, $`Z_1(𝐱)`$ has a further holomorphic extension $`Z_1^{ex}(𝐱)`$ to $`C`$. Finally, the theorem of Liouville<sup>16</sup><sup>16</sup>16See Secs. 167-168 of the text by Caratheodory cited above. on entire functions that do not have an essential singularity at $`\tau =\mathrm{}`$, together with Eq. (1C.16), then yields $$Z_1^{ex}(𝐱,\tau )=1\text{ for all }C.$$ (1C.20) Thus, we have shown that $`det(𝐱)=\nu (𝐱)`$, whereupon Eq. (1C.17) yields $`detY(𝐱)=1`$. End of proof. * Suppose that $`^{}(𝐱)`$ is also a solution of the HHP corresponding to $`(𝐯,_0,𝐱)`$. Since $`det(𝐱)=\nu (𝐱)`$, $`(𝐱)`$ is invertible. Conditions (1), (2), (3) and (4) in the definition of the HHP imply that $$Z_2(𝐱):=^{}(𝐱)(𝐱)^1\text{ is holomorphic throughout }C\overline{}(𝐱),$$ (1C.21) $$Z_2(𝐱,\mathrm{})=I,$$ (1C.22) $$\begin{array}{c}\text{the limits }Z_2^\pm (𝐱)\text{ exist and }\hfill \\ Y^{}(𝐱)Y(𝐱)^1=Z_2^+(𝐱)=Z_2^{}(𝐱)\text{ throughout }(𝐱),\hfill \end{array}$$ (1C.23) $$Z_2(𝐱)\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{,}$$ (1C.24) and $$Y^{}(𝐱)Y(𝐱)^1=Z_2^\pm (𝐱)\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{.}$$ (1C.25) The same kind of reasoning that was used in the proof of part (iii) of the theorem nets $`Z(𝐱)=I`$. So $`^{}(𝐱)=(𝐱)`$. End of proof. * When $`𝐱=𝐱_0`$, $`(x)`$ and its closure $`\overline{}(𝐱)`$ are empty. So, condition (1) of the HHP definition implies that $`(𝐱_0)`$ is holomorphic throughout $`C`$, whereupon condition (2) tells us that $`(𝐱_0)`$ has the value $`I`$ throughout $`C`$. \[$`^\pm (𝐱)`$ are empty sets when $`𝐱=𝐱_0`$; and conditions (3) and (4) hold trivially when $`𝐱=𝐱_0`$.\] End of proof. ### D. The generalized Geroch conjecture At this point we shall conjecture that for each $`\mathrm{}`$, where $`\mathrm{}`$ may be $`n`$ or $`n+`$, where $`n3`$, $`\mathrm{}`$ or ‘an’ (analytic), the following theorems hold: * There exists a subset $`𝒮_{}^{\mathrm{}}`$ of $`𝒮_{}`$ such that, for each $`_0𝒮_{}^{\mathrm{}}`$ and each $`𝐯K^{\mathrm{}}`$, there exists exactly one solution $`𝒮_{}^{\mathrm{}}`$ of the HHP corresponding to $`(𝐯,_0)`$, enabling us to define a mapping $$[𝐯]:𝒮_{}^{\mathrm{}}𝒮_{}^{\mathrm{}}$$ (1D.1) such that, for each $`_0𝒮_{}^{\mathrm{}}`$, $$[𝐯](_0)=$$ (1D.2) is that unique solution of the HHP corresponding to $`(𝐯,_0)`$. We then define our extension $`𝒦^{\mathrm{}}`$ of the K–C group by $$𝒦^{\mathrm{}}:=\{[𝐯]:𝐯K^{\mathrm{}}\}.$$ (1D.3) * The mapping $`[𝐯]`$ is the identity map on $`𝒮_{}^{\mathrm{}}`$ iff $`𝐯Z^{(3)}\times Z^{(4)}`$, where $$Z^{(i)}:=\{\delta ^{(i)},\delta ^{(i)}\}$$ (1D.4) and $$\delta ^{(i)}(\sigma )=I\text{ for all }\sigma ^{(i)}\text{.}$$ (1D.5) * The set $`𝒦^{\mathrm{}}`$ is a group of permutations of $`𝒮_{}^{\mathrm{}}`$ such that the mapping $`𝐯[𝐯]`$ is a homomorphism of $`K^{\mathrm{}}`$ onto $`𝒦^{\mathrm{}}`$; and the mapping $`\{\mathrm{𝐯𝐰}:𝐰Z^{(3)}\times Z^{(4)}\}[𝐯]`$ is an isomorphism \[viz, the isomorphism of $`K^{\mathrm{}}/(Z^{(3)}\times Z^{(4)})`$ onto $`𝒦^{\mathrm{}}`$\]. * The group $`𝒦^{\mathrm{}}`$ is transitive \[i.e., for each $`_0,𝒮_{}^{\mathrm{}}`$ there exists at least one element of $`𝒦^{\mathrm{}}`$ that transforms $`_0`$ into $``$\]. It will later be seen when we come to Thm. 35 that to prove the first part of the above generalized Geroch conjecture it is sufficient to prove that, for each $`𝐯𝒮_𝒴^{\mathrm{}}`$ with $`\mathrm{}=n`$, $`n+(n3)`$, $`\mathrm{}`$ or ‘an’, the solution $``$ of the HHP corresponding to $`(𝐯,^M)`$ exists, and $`𝒮_𝒴^{\mathrm{}}`$. For this reason, we shall now focus on the HHP corresponding to $`(𝐯,^M)`$. We shall begin with a study of an Alekseev-type singular integral equation and a Fredholm integral equation of the second kind that are, under suitable circumstances, equivalent to the HHP corresponding to $`(𝐯,^M)`$. Ultimately we shall have to return to the identification of the sets $`𝒮_{}^{\mathrm{}}`$ for $`\mathrm{}=n`$, $`n+`$, $`\mathrm{}`$ and ‘an’ (analytic), which will require us to introduce the concept of generalized Abel transforms of the initial data functions $`^{(3)}`$ and $`^{(4)}`$. ## 2 An Alekseev-type singular integral equation that is equivalent to the HHP corresponding to $`(𝐯,^M)`$ when $`𝐯K^{1+}`$ Using an ingenious argument G. A. Alekseev<sup>17</sup><sup>17</sup>17G. A. Alekseev, The method of the inverse scattering problem and singular integral equation for interacting massless fields, Dokl. Akad. Nauk SSSR 283, 577–582 (1985) \[Sov. Phys. Dokl. (USA) 30, 565 (1985)\], Exact solutions in the general theory of relativity, Trudy Matem. Inst. Steklova 176, 215–262 (1987). derived a singular integral equation, supposing that $`(\tau )`$ was analytic in a neighborhood of $`\{r,s\}`$ except for branch points of index $`1/2`$ at $`\tau =r`$ and $`\tau =s`$. We shall now show that the same type integral equation arises in connection with solutions of our new HHP that need not be analytic. ### A. A preliminary theorem Henceforth, whenever there is no danger of ambiguity, the arguments ‘$`𝐱`$’ and ‘$`𝐱_0`$’ will be suppressed. For example, ‘$`(\tau )`$’ and ‘$`^\pm (\sigma )`$’ will generally be used as abbreviations for ‘$`(𝐱,\tau )`$’ and ‘$`^\pm (𝐱,\sigma )`$’, respectively; and ‘$`\nu (\tau )`$’, ‘$`\nu ^\pm (\sigma )`$’ and ‘$`\overline{}`$’ will generally stand for ‘$`\nu (𝐱,\tau )`$’, ‘$`\nu ^\pm (𝐱,\sigma )`$’ and ‘$`\overline{}(𝐱)`$’, respectively. ###### THEOREM 4 (Alekseev preliminaries) * Suppose that the solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,_0,𝐱)`$ exists. Then, for each $`\tau C\overline{}(𝐱)`$, $`[\nu ^+(\sigma ^{})]^1{\displaystyle \frac{^+(\sigma ^{})+^{}(\sigma ^{})}{\sigma ^{}\tau }}\text{ is summable over }\sigma ^{}\overline{}(𝐱),`$ $`\text{with assigned orientation in the direction of increasing }\sigma ^{},`$ (2A.1) and $$\left[\nu (\tau )\right]^1(\tau )=I+\frac{1}{2\pi i}_\overline{}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1\frac{^+(\sigma ^{})+^{}(\sigma ^{})}{\sigma ^{}\tau },$$ (2A.2) where the meaning we attribute to the symbol $`_\overline{}`$ should be obvious. * Moreover, for each $`\sigma (𝐱)`$, $`[\nu ^+(\sigma ^{})]^1{\displaystyle \frac{^+(\sigma ^{})+^{}(\sigma ^{})}{\sigma ^{}\sigma }}\text{ is summable over }\sigma ^{}\overline{}(𝐱)`$ in the principal value (PV) sense, (2A.3) and $$\frac{1}{2}[\nu ^+(\sigma )]^1\left\{^+(\sigma )^{}(\sigma )\right\}=I+\frac{1}{2\pi i}_\overline{}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1\frac{^+(\sigma ^{})+^{}(\sigma ^{})}{\sigma ^{}\sigma }.$$ (2A.4) Proofs: * From Thms. 3(i) and (ii), the function of $`\sigma ^{}`$ given by $`\nu ^\pm (\sigma ^{})^1^\pm (\sigma ^{})(\sigma ^{}\tau )^1`$ is continuous throughout $``$ and is bounded at $`𝐱`$, while $`^\pm (\sigma ^{})(\sigma ^{}\tau )^1`$ is bounded at $`𝐱_0`$. Moreover, it is clear that $`\nu ^\pm (\sigma ^{})`$ and $`\nu ^\pm (\sigma ^{})^1`$ are summable on $`\overline{}`$, and $`\nu ^{}(\sigma ^{})=\nu ^+(\sigma ^{})`$ throughout $``$. Statement (2A.1) can now be obtained by employing the well-known theorem<sup>18</sup><sup>18</sup>18See Integration, by Edward J. McShane (Princeton University Press, 1944). that the product of any complex-valued function which is summable on $`[a,b]R^1`$ by a function which is continuous and bounded on $`[a,b]`$(any given finite set) is also summable on $`[a,b]`$. To obtain the conclusion (2A.2), one employs Cauchy’s integral formula and the HHP condition $`(\mathrm{})=I`$ to infer that $$\nu (\tau )^1(\tau )=I\frac{1}{2\pi i}_\mathrm{\Lambda }𝑑\tau ^{}\frac{[\nu (\tau ^{})]^1(\tau ^{})}{\tau ^{}\tau },$$ (2A.5) where $`\mathrm{\Lambda }`$ is a closed positively oriented contour enclosing $`\overline{}`$ but not the point $`\tau `$, which we may assume to be rectangular. This equation can be expressed in the form $$\nu (\tau )^1(\tau )=I\frac{1}{2\pi i}_{\mathrm{\Lambda }^+}𝑑\tau ^{}\frac{[\nu ^+(\tau ^{})]^1^+(\tau ^{})}{\tau ^{}\tau }\frac{1}{2\pi i}_\mathrm{\Lambda }^{}𝑑\tau ^{}\frac{[\nu ^{}(\tau ^{})]^1^{}(\tau ^{})}{\tau ^{}\tau },$$ (2A.6) where $`\mathrm{\Lambda }^\pm :=\mathrm{\Lambda }\overline{C}^\pm `$ denote the parts of the contour $`\mathrm{\Lambda }`$ that lie respectively in the upper and lower half planes, $`\overline{C}^\pm `$. To evaluate each of the integrals, one applies a well known generalization<sup>19</sup><sup>19</sup>19See Remark 2 in Sec. 2, Ch. II, of Analytic Functions by M. A. Evgrafov (Dover Publications, 1978). of Cauchy’s integral theorem which asserts that the integral of a function about a simple piecewise smooth contour $`𝒦`$ is zero if the given function is holomorphic throughout $`𝒦_{int}`$ and is continuous throughout $`𝒦𝒦_{int}`$. In the case of the first integral, we select the contour as in Fig. 2. The other integral is evaluated in a similar way, using a contour in $`\overline{C}^{}`$. Here $`a^i`$ and $`b^i`$ are the left and right endpoints, respectively, of the arc $`\overline{}^{(i)}`$. The radius of each semicircular arc is $`\alpha `$ and each of the vertical segments of the closed contours has length $`\sqrt{2}\alpha `$. One ultimately takes the limit as $`\alpha 0`$. From a well known theorem<sup>20</sup><sup>20</sup>20See Cor. 27.7 in Ref. 18. on Lebesgue integrals, $$_{a^i+\alpha }^{b^i\alpha }𝑑\sigma ^{}\frac{[\nu ^\pm (\sigma ^{})]^1^\pm (\sigma ^{})}{\sigma ^{}\tau }_{a^i}^{b^i}𝑑\sigma ^{}\frac{[\nu ^\pm (\sigma ^{})]^1^\pm (\sigma ^{})}{\sigma ^{}\tau }\text{ as }\alpha 0.$$ (2A.7) Upon applying the above statement (2A.7) and the easily proved statement that the integral on each semicircular arc $`0`$ as $`\alpha 0`$, and using the fact that $`\nu ^{}(\sigma ^{})=\nu ^+(\sigma ^{})`$ for all $`\sigma ^{}(𝐱)`$, one obtains the conclusion (2A.2). End of proof. * To obtain statement (2A.3) and Eq. (2A.4) when $`\sigma ^{(3)}(𝐱)`$, we again employ the Cauchy integral formula and the generalized Cauchy integral theorem, this time using (for the integral over $`\mathrm{\Lambda }^+`$) the positively oriented closed contours depicted in Fig. 3. The case $`\sigma ^{(4)}(𝐱)`$ is treated similarly. Here the radius of the semicircular arc about $`\sigma `$ is $`\beta `$ and each of the vertical segments of the left closed contour has length $`\sqrt{2}\beta `$. The radius of each of the other semicircular arcs is $`\alpha `$, and each of the vertical segments of the right closed contour has length $`\sqrt{2}\alpha `$. One ultimately takes the limit as $`\alpha 0`$ followed by the limit as $`\beta 0`$. It is clear that the integral on the semicircular arc with center $`\sigma `$ has the limit $`\frac{1}{2}\nu ^+(\sigma )^+(\sigma )`$ as $`\beta 0`$. End of proof. ### B. Derivation of an Alekseev-type singular integral equation Proceeding from equations (2A.2) and (2A.4), one can construct a singular integral equation of the Alekseev type and, if $`𝐯K^{1+}`$, a Fredholm equation of the second kind. We begin by observing that Eq. (1B.13) implies that, for each $`\sigma (𝐱)\{r_0,s_0\}`$, $`{\displaystyle \frac{1}{2}}\left\{^{M+}(\sigma )+^M(\sigma )\right\}`$ $`=`$ $`\left(\begin{array}{cc}1& i(\sigma z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\left(\begin{array}{cc}1& i(\sigma z_0)\\ 0& 1\end{array}\right),`$ (2B.7) and $`{\displaystyle \frac{1}{2}}[\nu ^+(\sigma )]^1\left\{^{M+}(\sigma )^M(\sigma )\right\}=`$ (2B.14) $`\left(\begin{array}{cc}1& i(\sigma z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& i(\sigma z_0)\\ 0& 1\end{array}\right).`$ If $``$ is a solution of the HHP corresponding to $`(𝐯,^M)`$, Eq. (3) tells us that, for any $`\sigma (𝐱)`$, $$^\pm (\sigma )v^{(i)}(\sigma )=Y^{(i)}(\sigma )^{M\pm }(\sigma ),$$ (2B.15) and, therefore, $`{\displaystyle \frac{1}{2}}\left\{^+(\sigma )+^{}(\sigma )\right\}v^{(i)}(\sigma )=`$ (2B.22) $`Y^{(i)}(\sigma )\left(\begin{array}{cc}1& i(\sigma z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right)\left(\begin{array}{cc}1& i(\sigma z_0)\\ 0& 1\end{array}\right),`$ and $`{\displaystyle \frac{1}{2}}[\nu ^+(\sigma )]^1\left\{^+(\sigma )^{}(\sigma )\right\}v^{(i)}(\sigma )=`$ (2B.29) $`Y^{(i)}(\sigma )\left(\begin{array}{cc}1& i(\sigma z)\\ 0& 1\end{array}\right)\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right)\left(\begin{array}{cc}1& i(\sigma z_0)\\ 0& 1\end{array}\right).`$ This motivates the introduction of two new $`2\times 2`$ matrices. Dfn. of functions $`W^{(i)}(𝐱)`$ and $`𝒴^{(i)}(𝐱)`$ * For each $`𝐯K`$, we let $`W^{(i)}(𝐱)`$ denote the function whose domain is $`^{(i)}`$ and whose value for each $`\sigma ^{(i)}`$ is $`W^{(i)}(𝐱,\sigma )`$ $`:=`$ $`W^{(i)}(𝐱)(\sigma ):=v^{(i)}(\sigma )\left(\begin{array}{cc}1& i(\sigma z_0)\\ 0& 1\end{array}\right),`$ (2B.32) and, for each solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,^M,𝐱)`$, we let $`𝒴^{(i)}(𝐱)`$ denote the function whose domain is $`^{(i)}(𝐱)`$ and whose value for each $`\sigma ^{(i)}(𝐱)`$ is $`𝒴^{(i)}(𝐱,\sigma )`$ $`:=`$ $`𝒴^{(i)}(𝐱)(\sigma ):=Y^{(i)}(𝐱,\sigma )\left(\begin{array}{cc}1& i(\sigma z)\\ 0& 1\end{array}\right).`$ (2B.35) End of Dfn. In terms of these matrices we may write \[suppressing ‘$`𝐱`$’\] $$^\pm (\sigma )W^{(i)}(\sigma )=𝒴^{(i)}(\sigma )\left(\begin{array}{cc}1& 0\\ 0& \nu ^\pm (\sigma )\end{array}\right)$$ (2B.36) as well as $`{\displaystyle \frac{1}{2}}\left\{^+(\sigma )+^{}(\sigma )\right\}W^{(i)}(\sigma )`$ $`=`$ $`𝒴^{(i)}(\sigma )\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),`$ (2B.39) and $`{\displaystyle \frac{1}{2}}[\nu ^+(\sigma )]^1\left\{^+(\sigma )^{}(\sigma )\right\}W^{(i)}(\sigma )`$ $`=`$ $`𝒴^{(i)}(\sigma )\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right).`$ (2B.42) Dfns. of $`W(𝐱)`$, $`𝒴(𝐱)`$, $`W_a(𝐱)`$ and $`𝒴_a(𝐱)`$ * Let $`W(𝐱)`$ and $`𝒴(𝐱)`$ denote the functions<sup>21</sup><sup>21</sup>21We shall frequently suppress $`𝐱`$. with domain $`(𝐱)`$ and values $$\begin{array}{c}W(𝐱,\sigma ):=W(𝐱)(\sigma ):=W^{(i)}(𝐱,\sigma )\text{ and }\hfill \\ 𝒴(𝐱,\sigma ):=𝒴(𝐱)(\sigma ):=𝒴^{(i)}(𝐱,\sigma )\hfill \\ \text{for each }i\{3,4\}\text{ and }\sigma ^{(i)}(𝐱)\text{.}\hfill \end{array}$$ (2B.43) Moreover, let $$\begin{array}{c}W_a(𝐱,\sigma ):=a^{th}\text{ column of }W(𝐱,\sigma )\text{ and }\hfill \\ 𝒴_a(𝐱,\sigma ):=a^{th}\text{ column of }𝒴(𝐱,\sigma ),\text{ where }a\{1,2\}.\hfill \end{array}$$ (2B.44) End of Dfn. ###### THEOREM 5 (Alekseev-type equation) For each $`𝐯K`$, $`𝐱D`$, solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,^M,𝐱)`$, $`\tau C\overline{}(𝐱)`$ and $`\sigma (𝐱)`$, the following statement holds: $$[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{})W_2^T(\sigma ^{})(\sigma ^{}\tau )^1\text{ is summable over }\sigma ^{}\overline{}(𝐱),$$ (2B.45) $$\nu (\tau )^1(\tau )=I+\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{})\frac{W_2^T(\sigma ^{})J}{\sigma ^{}\tau },$$ (2B.46) $$\begin{array}{c}\hfill [\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{})W_2^T(\sigma ^{})(\sigma ^{}\sigma )^1\text{ is summable over }\sigma ^{}\overline{}(𝐱)\\ \hfill \text{in the PV sense,}\end{array}$$ (2B.47) $`𝒴_2(\sigma )`$ $`=`$ $`W_2(\sigma ){\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{}){\displaystyle \frac{W_2^T(\sigma ^{})JW_2(\sigma )}{\sigma ^{}\sigma }},`$ (2B.48) and $`0`$ $`=`$ $`W_1(\sigma )+{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{}){\displaystyle \frac{W_2^T(\sigma ^{})JW_1(\sigma )}{\sigma ^{}\sigma }}.`$ (2B.49) Here we have employed the symbol $`J:=i\mathrm{\Omega }=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)`$. Proof: The statements (2B.45) to (2B.47) are obtained by using Eqs. (2B.39) and (2B.42) together with the relation $$W(\sigma )^1=JW^T(\sigma )J$$ to replace $`^+(\sigma ^{})^{}(\sigma ^{})`$ and $`[\nu ^+(\sigma ^{})]^1[^+(\sigma ^{})+^{}(\sigma )]`$ in statements (2A.1) to (2A.3) of Thm. 4. The same replacements are to be made in the integrands on the right side of Eq. (2A.4) in Thm. 4. Equation (2B.48) is obtained by multiplying both sides of Eq. (2A.4) by $`W_2(\sigma )`$ and replacing the product on the left side with the second column of (2B.42) multiplied by $`[\nu ^+(\sigma )]^1`$. Equation (2B.49) is obtained by multiplying both sides of Eq. (2A.4) by $`W_1(\sigma )`$ and replacing the product on the left side with the first column of (2B.42). End of proof. Equation (2B.49) has the form of the singular integral equation which Alekseev obtained in the analytic case. ### C. Extension of the function $`𝒴(𝐱)`$ from $`(𝐱)`$ to $`\overline{}(𝐱)`$ Since $`C^TJC=0`$ (the zero matrix) for any $`2\times 1`$ matrix $`C`$, Eq. (2B.48) is expressible in the following form for each $`i\{3,4\}`$: $`𝒴_2^{(i)}(\sigma )=W_2^{(i)}(\sigma )`$ (2C.1) $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _{a^i}^{b^i}}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1^{(i)}(\sigma ^{})W_2^{(i)}(\sigma ^{})^TJ\left[{\displaystyle \frac{W_2^{(i)}(\sigma )W_2^{(i)}(\sigma ^{})}{\sigma ^{}\sigma }}\right]`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _{a^{7i}}^{b^{7i}}}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1^{(7i)}(\sigma ^{})W_2^{(7i)}(\sigma ^{})^TJ\left[{\displaystyle \frac{W_2^{(i)}(\sigma )}{\sigma ^{}\sigma }}\right],`$ for all $`\sigma ^{(i)}(𝐱)`$, where recall that $`a^i:=inf\{x^i,x_0^i\}`$ and $`b^i:=sup\{x^i,x_0^i\}`$. Without indicating the parallel proof, we simply remark that one can also show that $`𝒴_1^{(i)}(\sigma )=W_1^{(i)}(\sigma )`$ (2C.2) $`+{\displaystyle \frac{1}{\pi i}}{\displaystyle _{a^i}^{b^i}}𝑑\sigma ^{}\nu ^+(\sigma ^{})𝒴_2^{(i)}(\sigma ^{})W_1^{(i)}(\sigma ^{})^TJ\left[{\displaystyle \frac{W_1^{(i)}(\sigma )W_1^{(i)}(\sigma ^{})}{\sigma ^{}\sigma }}\right]`$ $`+{\displaystyle \frac{1}{\pi i}}{\displaystyle _{a^{7i}}^{b^{7i}}}𝑑\sigma ^{}\nu ^+(\sigma ^{})𝒴_2^{(7i)}(\sigma ^{})W_1^{(7i)}(\sigma ^{})^TJ\left[{\displaystyle \frac{W_1^{(i)}(\sigma )}{\sigma ^{}\sigma }}\right],`$ Now, from Thms. 3(i) and (ii), Eq. (2B.39) and Eq. (2B.42), $$\begin{array}{c}\nu ^+(\sigma ^{})𝒴_2(\sigma ^{})W_1(\sigma ^{})^TJ\hfill \\ \text{and }[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{})W_2(\sigma ^{})^TJ\hfill \\ \text{are summable over }\sigma ^{}\overline{}(𝐱).\hfill \end{array}$$ (2C.3) From the definition of $`W^{(i)}`$ by Eq. (2B.32) and the definition of $`K^{\mathrm{}}`$, the following statement holds for each $`𝐱D`$ and $`i\{3,4\}`$: $$\begin{array}{c}\text{If }𝐯K^1\text{, then }W^{(i)}\text{ is }𝐂^1\text{ throughout }^{(i)}\text{,}\hfill \\ [W^{(i)}(\sigma ^{})W^{(i)}(\sigma )](\sigma ^{}\sigma )^1\text{ is a continuous }\hfill \\ \text{function of }(\sigma ^{},\sigma )\text{ throughout }^{(i)}\times ^{(i)},\text{ and}\hfill \\ W^{(i)}(\sigma )(\sigma ^{}\sigma )^1\text{ is a }𝐂^1\text{ function of }(\sigma ^{},\sigma )\hfill \\ \text{throughout }\overline{}^{(7i)}(𝐱)\times \stackrel{ˇ}{}^{(i)}(x^{7i}),\hfill \end{array}$$ (2C.4) where $$\begin{array}{ccc}\hfill \stackrel{ˇ}{}^{(3)}(s)& :=& \{\sigma ^{(3)}:\sigma <s\},\text{ and }\hfill \\ \hfill \stackrel{ˇ}{}^{(4)}(r)& :=& \{\sigma ^{(4)}:r<\sigma \}.\hfill \end{array}$$ (2C.5) Note that (See Fig. 4) $$^{(i)}(𝐱)\stackrel{ˇ}{}^{(i)}(x^{7i})^{(i)}.$$ (2C.6) From the above statements (2C.3) and (2C.4), and from the theorem that asserts the summability over a finite interval of the product of a summable function by a continuous function, the extension of $`𝒴^{(i)}(𝐱)`$ that we shall define below exists. Note that $`\mathrm{}`$ is $`n1`$, $`n+`$ (with $`n1`$), $`\mathrm{}`$ or ‘an’. Dfn. of an extension of $`𝒴^{(i)}(𝐱)`$ when $`𝐯K^{\mathrm{}}`$ * For each $`𝐯K^{\mathrm{}}`$, $`𝐱D`$, solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,^M,𝐱)`$ and $`i\{3,4\}`$, let $`𝒴^{(i)}(𝐱)`$ denote the function whose extended domain is $`\stackrel{ˇ}{}(x^{7i})`$ and whose value for each $`\sigma \stackrel{ˇ}{}^{(i)}(𝐱^{7i})`$ is given by \[suppressing ‘$`𝐱`$’\] $`𝒴_1^{(i)}(\sigma )`$ $`:=`$ right side of Eq. (2C.2), (2C.7) $`𝒴_2^{(i)}(\sigma )`$ $`:=`$ right side of Eq. (2C.1). (2C.8) End of Dfn. ###### LEMMA 6 (Continuity and differentiability of $`W^{(i)}`$) * If $`𝐯K^{\mathrm{}}`$, then $`W^{(i)}`$ is $`𝐂^{\mathrm{}}`$ throughout its domain $`^{(i)}`$, and the function whose domain is $`\overline{}^{(7i)}(𝐱)\times \stackrel{ˇ}{}^{(i)}(x^{7i})`$ and whose values for each $`(\sigma ^{},\sigma )`$ in this domain is $`W^{(i)}(\sigma )(\sigma ^{}\sigma )^1`$ is also $`𝐂^{\mathrm{}}`$. * If $`𝐯K^{\mathrm{}}`$, then the function of $`(\sigma ^{},\sigma )`$ whose domain is $`^{(i)}\times ^{(i)}`$ and whose value for each $`(\sigma ^{},\sigma )`$ in this domain is $`[W^{(i)}(\sigma )W^{(i)}(\sigma ^{})]/(\sigma ^{}\sigma )`$ is $`𝐂^{n1}`$ if $`\mathrm{}`$ is $`n1`$, is $`𝐂^{(n1)+}`$ if $`\mathrm{}`$ is $`n^+`$ ($`n1`$), is $`𝐂^{\mathrm{}}`$ if $`\mathrm{}`$ is $`\mathrm{}`$, and is $`𝐂^{an}`$ if $`\mathrm{}`$ is ‘an’. Proofs: * The conclusion follows by using the definition of $`W^{(i)}`$ by Eq. (2B.32) together with the definition of $`K^{\mathrm{}}`$. End of proof. * The conclusions when $`\mathrm{}`$ is $`n`$, $`\mathrm{}`$ or ‘an’ are well known. As regards the case when $`\mathrm{}`$ is $`n^+`$ ($`n1`$), one can construct a simple proof (which we shall not reproduce here) using the relation $$\frac{W^{(i)}(\sigma )W^{(i)}(\sigma ^{})}{\sigma \sigma ^{}}=_0^1𝑑t(DW^{(i)})(t\sigma +(1t)\sigma ^{}),$$ (2C.9) where $`D^pW^{(i)}`$ ($`1pn`$) denotes the function whose domain is $`^{(i)}`$ and whose value for each $`\sigma ^{(i)}`$ is $$(D^pW^{(i)})(\sigma ):=\frac{^pW^{(i)}(\sigma )}{\sigma ^p};$$ (2C.10) and $`DW^{(i)}:=D^1W^{(i)}`$. End of proof. We shall leave the proof of the following basic lemma to the reader. ###### LEMMA 7 (Integral of product) Suppose $`[a,b]R^1`$, $`S`$ is a connected open subset of $`R^m`$ ($`m1`$), $`f`$ is a real-valued function defined almost everywhere on and summable over $`[a,b]`$, and $`g`$ is a real-valued function whose domain is $`[a,b]\times S`$ and which is continuous. Let $`\sigma :=(\sigma ^1,\mathrm{},\sigma ^m)`$ denote any member of $`S`$, and let $`F`$ denote the function whose domain is $`S`$ and whose value at each $`\sigma S`$ is $$F(\sigma ):=_a^b𝑑\sigma ^{}f(\sigma ^{})g(\sigma ^{},\sigma ).$$ (2C.11) Then the following statements hold: * $`F`$ is continuous. * If $`g(\sigma ^{},\sigma )/\sigma ^k`$ exists for each $`k\{1,\mathrm{},m\}`$ and is a continuous function of $`(\sigma ^{},\sigma )`$ throughout $`[a,b]\times S`$, then $`F`$ is $`𝐂^1`$ and $$\frac{F(\sigma )}{\sigma ^k}=_a^b𝑑\sigma ^{}f(\sigma ^{})\frac{g(\sigma ^{},\sigma )}{\sigma ^k}$$ (2C.12) for all $`k\{1,\mathrm{},m\}`$ and $`\sigma S`$. In particular, $`F`$ is $`𝐂^1`$ if $`g`$ is $`𝐂^1`$. * Assume (for simplicity) that $`m=1`$. If there exists a positive integer $`p`$ such that $`^pg(\sigma ^{},\sigma )/(\sigma )^p`$ exists and is a continuous function of $`(\sigma ^{},\sigma )`$ throughout $`[a,b]\times S`$, then $`F`$ is $`bC^p`$ and $$\frac{^pF(\sigma )}{\sigma ^p}=_a^b𝑑\sigma ^{}f(\sigma ^{})\frac{^pg(\sigma ^{},\sigma )}{\sigma ^p}$$ (2C.13) for all $`\sigma S`$. In particular, $`F`$ is $`𝐂^p`$ if $`g`$ is $`𝐂^p`$. * Assume (for simplicity) that $`m=1`$. If $`[c,d]S`$ and $`g`$ obeys a Hölder condition on $`[a,b]\times [c,d]`$, then $`F`$ obeys a Hölder condition on $`[c,d]`$. * Assume (for simplicity) that $`m=1`$. If $`g`$ is analytic \[i.e., if $`g`$ has an analytic extension to an open subset $`[aϵ,b+ϵ]\times S`$ of $`R^2`$, then $`F`$ is analytic. ###### LEMMA 8 (Generalization of Lem. 7) All of the conclusions of the preceding lemma remain valid when the only alteration in the premises is to replace the statement that $`f`$ and $`g`$ are real valued by the statement that they are complex valued or are finite matrices (such that the product $`fg`$ exists) with complex-valued elements. Proof: Use the definition $$_a^b𝑑\sigma ^{}h(\sigma ^{}):=_a^b𝑑\sigma ^{}\mathrm{R}eh(\sigma ^{})+i_a^b𝑑\sigma ^{}\mathrm{I}mh(\sigma ^{})$$ for any complex-valued function $`h`$ whose real and imaginary parts are summable over $`[a,b]`$. The rest is obvious. End of proof. ###### THEOREM 9 (Continuity and differentiability of extended $`𝒴^{(i)}(𝐱)`$) For each $`𝐯K^{\mathrm{}}`$, $`𝐱D`$, solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,^M,𝐱)`$ and $`i\{3,4\}`$, $`𝒴^{(i)}(𝐱)`$ \[see Eqs. (2C.7) and (2C.8)\] is $`𝐂^{n1}`$ if $`\mathrm{}`$ is $`n`$, is $`𝐂^{(n1)+}`$ if $`\mathrm{}`$ is $`n+`$, is $`𝐂^{\mathrm{}}`$ if $`\mathrm{}`$ is $`\mathrm{}`$ and is $`𝐂^{an}`$ if $`\mathrm{}`$ is ‘an’. Proof: Apply Lemmas 6, 7 and 8 to the definitions (2C.7) and (2C.8) of $`𝒴_1^{(i)}(𝐱)`$ and $`𝒴_2^{(i)}(𝐱)`$. It is then easily shown that the second term on the right side of each of the Eqs. (2C.2) and (2C.1) \[with $`\sigma \stackrel{ˇ}{}^{(i)}(x^{7i})`$\] is $`𝐂^{n1}`$ if $`\mathrm{}`$ is $`n`$, is $`𝐂^{(n1)+}`$ if $`\mathrm{}`$ is $`n+`$, is $`𝐂^{\mathrm{}}`$ if $`\mathrm{}`$ is $`\mathrm{}`$ and is $`𝐂^{an}`$ if $`\mathrm{}`$ is ‘an’. The first and third terms on the right sides of each of the Eqs. (2C.2) and (2C.1) are, on the other hand, both $`𝐂^{\mathrm{}}`$ even when $`\mathrm{}`$ is $`n`$ or is $`n+`$. However, a $`𝐂^n`$ function is also a $`𝐂^{n1}`$ function; and a $`𝐂^{n+}`$ function is also a $`𝐂^{(n1)+}`$ function. End of proof. Dfns. of $`𝒴^{(i)}`$, $`𝒴`$ and the partial derivatives of $`𝒴`$ * Henceforth, $`𝒴^{(i)}(i\{3,4\})`$ will denote the function whose domain is $$\mathrm{d}om𝒴^{(i)}:=\{(𝐱,\sigma ):𝐱D,\sigma \stackrel{ˇ}{}(x^{7i})\}$$ (2C.14) and whose values are given by $$𝒴^{(i)}(𝐱,\sigma ):=𝒴^{(i)}(𝐱)(\sigma ),$$ (2C.15) where $`𝒴^{(i)}(𝐱)`$ is the extension of the original $`𝒴^{(i)}(𝐱)`$ that is defined by Eqs. (2C.7) and (2C.8). We shall let $`𝒴`$ denote the function whose domain is $$\mathrm{d}om𝒴:=\{(𝐱,\sigma ):𝐱D,\sigma \overline{}(𝐱)\}$$ and whose values are given by $$𝒴(𝐱,\sigma ):=𝒴^{(i)}(𝐱,\sigma )\text{ whenever }\sigma \overline{}^{(i)}(𝐱).$$ \[Thus, $`𝒴(𝐱,\sigma )=𝒴(𝐱)(\sigma )`$.\] Also, for each $`𝐱D`$, $`i\{3,4\}`$ and $`\sigma \overline{}^{(i)}(𝐱)`$, we shall let $$\frac{^{l+m+n}𝒴(𝐱,\sigma )}{r^ls^m\sigma ^n}:=\frac{^{l+m+n}𝒴^{(i)}(𝐱,\sigma )}{r^ls^m\sigma ^n},$$ if the above partial derivative of $`𝒴^{(i)}`$ exists. End of Dfn. The domain of $`𝒴^{(i)}`$, as defined above, is an open subset of $`R^3`$; and (though the domain of $`𝒴`$ is not an open subset of $`R^3`$) the partial derivatives of $`𝒴`$ are defined in terms of partial derivatives of $`𝒴^{(i)}`$ and, therefore, employ sequences of points in $`R^3`$ which may converge to a given point along any direction in $`R^3`$. This has formal advantages when one employs the derivatives of $`𝒴`$ at the boundary of its domain. ###### COROLLARY 10 (The extension $`𝒴(𝐱)`$ when $`𝐯K^{\mathrm{}}`$) * Suppose $`𝐯K^1`$, $`𝐱D`$ and the solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,^M,𝐱)`$ exists. Then $`𝒴(𝐱)`$ has a unique continuous extension to $`\overline{}(𝐱)`$. * If $`𝐯K^{\mathrm{}}`$, then the extension $`𝒴(𝐱)`$ is $`𝐂^{n1}`$ if $`\mathrm{}`$ is $`n`$, is $`𝐂^{(n1)+}`$ if $`\mathrm{}`$ is $`n+`$, is $`𝐂^{\mathrm{}}`$ if $`\mathrm{}`$ is $`\mathrm{}`$ and is $`𝐂^{an}`$ if $`\mathrm{}`$ is ‘an’. Proof: Statement (ii) of this corollary follows from Thm. 9. The uniqueness follows, of course, from the fact that a function defined and continuous on an open subset of $`R^1`$ has no more than one continuous extension to the closure of that subset. End of proof. ### D. Equivalence of the HHP to an Alekseev-type equation when $`𝐯K^{1+}`$ ###### THEOREM 11 (HHP-Alekseev equivalence theorem) Suppose $`𝐯K^{1+}`$ and $`𝐱D`$, and suppose that $`(𝐱)`$ and $`𝒴_1(𝐱)`$ are $`2\times 2`$ and $`2\times 1`$ matrix functions, respectively, such that $$\mathrm{d}om(𝐱)=C\overline{}(𝐱),\mathrm{d}om𝒴_1(𝐱)=\overline{}(𝐱)\text{ and }𝒴_1(𝐱)\text{ is }𝐂^{0+}.$$ (2D.1) Then the following two statements are equivalent to one another: * The function $`(𝐱)`$ is a solution of the HHP corresponding to $`(𝐯,^M,𝐱)`$, and $`𝒴_1(𝐱)`$ is the function whose restriction to $`(𝐱)`$ is defined in terms of $`^+(𝐱)+^{}(𝐱)`$ by Eq. (2B.39) \[where $`𝐱`$ is suppressed\] and whose existence and uniqueness \[for the given $`(𝐱)`$\] is asserted by Cor. 10 when $`\mathrm{}`$ is $`1+`$. * The restriction of $`𝒴_1(𝐱)`$ to $`(𝐱)`$ is a solution of the singular integral equation (2B.49), and $`(𝐱)`$ is defined in terms of $`𝒴_1(𝐱)`$ by Eq. (2B.46) \[where $`𝐱`$ is suppressed\]. Proof: That (i) implies (ii) has already been proved. \[See Thm. 5 and Cor. 10.\] The proof that (ii) implies (i) will be given in four parts: * Assume that statement (ii) is true. From the definition of $`(𝐱)`$ by Eq. (2B.46), $$(𝐱)\text{ is holomorphic;}$$ (2D.2) and, from two theorems of Plemelj<sup>22</sup><sup>22</sup>22See Secs. 16 and 17 of Ch. II of Ref. 13 (pp. 37-43). $$^+(𝐱)\text{ and }^{}(𝐱)\text{ exist }$$ (2D.3) and, since $`\nu ^{}(\sigma )=\nu ^+(\sigma )`$ for all $`\sigma (𝐱)`$, $$\frac{1}{2}\left[^+(\sigma )+^{}(\sigma )\right]=𝒴_1(\sigma )W_2^T(\sigma )J$$ (2D.4) and $$\frac{1}{2}[\nu ^+(\sigma )]^1\left[^+(\sigma )^{}(\sigma )\right]=I\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{})\frac{W_2^T(\sigma ^{})J}{\sigma ^{}\sigma }$$ (2D.5) for all $`\sigma (𝐱)`$. Upon multiplying Eqs. (2D.4) and (2D.5) through by $`W_2(\sigma )`$ on the right, one obtains, for all $`\sigma (𝐱)`$, $`{\displaystyle \frac{1}{2}}\left[^+(\sigma )+^{}(\sigma )\right]W_2(\sigma )`$ $`=`$ $`\left(\begin{array}{c}0\\ 0\end{array}\right),`$ (2D.8) and $`{\displaystyle \frac{1}{2}}[\nu ^+(\sigma )]^1\left[^+(\sigma )^{}(\sigma )\right]W_2(\sigma )`$ $`=`$ $`𝒴_2(\sigma ),`$ (2D.9) where $`𝒴_2(𝐱)`$ has the domain $`\overline{}(𝐱)`$ and the values $`𝒴_2(\sigma ):=W_2(\sigma )`$ (2D.10) $`{\displaystyle \frac{1}{\pi I}}{\displaystyle _\overline{}}𝑑\sigma ^{}[\nu ^+(\sigma ^{})]^1𝒴_1(\sigma ^{}){\displaystyle \frac{W_2^T(\sigma ^{})J[W_2(\sigma )W_2(\sigma ^{})]}{\sigma ^{}\sigma }}`$ $`\text{for all }\sigma \overline{}(𝐱).`$ From Lemmas 6(ii), 7(iv) and 8, $$𝒴_2(𝐱)\text{ is }𝐂^{0+}.$$ (2D.11) Upon multiplying Eqs. (2D.4) and (2D.5) through by $`W_1(\sigma )`$ on the right, upon using the fact that $`detW(\sigma )=1`$ is equivalent to the equation $$W_2^T(\sigma )JW_1(\sigma )=(1),$$ (2D.12) and, upon using Eq. (2B.49), one obtains, for all $`\sigma (𝐱)`$, $`{\displaystyle \frac{1}{2}}\left[^+(\sigma )+^{}(\sigma )\right]W_1(\sigma )`$ $`=`$ $`𝒴_1(\sigma )`$ (2D.13) and $`{\displaystyle \frac{1}{2}}\left[^+(\sigma )^{}(\sigma )\right]W_1(\sigma )`$ $`=`$ $`\left(\begin{array}{c}0\\ 0\end{array}\right).`$ (2D.16) * We next note that the four equations (2D.8), (2D.9), (2D.13) and (2D.16) are collectively equivalent to the single equation $$^\pm (\sigma )W(\sigma )=𝒴(\sigma )\left(\begin{array}{cc}1& 0\\ 0& \nu ^\pm (\sigma )\end{array}\right)\text{ for all }\sigma (𝐱),$$ (2D.17) where $`𝒴(\sigma )`$ is defined to be the $`2\times 2`$ matrix whose first and second columns are $`𝒴_1(\sigma )`$ and $`𝒴_2(\sigma )`$, respectively. From the definition of $`W(\sigma )`$ by Eqs. (2B.32) and (2B.43), and from the expression for $`^M(\tau )`$ that is given by Eq. (1B.13), Eq. (2D.17) is equivalent to the statement $`^+(\sigma )v^{(i)}(\sigma )[^{M+}(\sigma )]^1`$ $`=`$ $`^{}(\sigma )v^{(i)}(\sigma )[^M(\sigma )]^1`$ (2D.18) $`=`$ $`Y(\sigma )\text{ for all }\sigma ^{(i)}(𝐱),`$ where $$Y(\sigma ):=𝒴(\sigma )\left(\begin{array}{cc}1& i(\sigma z)\\ 0& 1\end{array}\right)\text{ for all }\sigma \overline{}(𝐱).$$ (2D.19) From the above Eq. (2D.19) and from statements (2D.1) and (2D.11), $$\begin{array}{c}\text{the function }Y(𝐱)\text{ whose domain is }\overline{}(𝐱)\text{ and whose value for each }\hfill \\ \sigma \overline{}(𝐱)\text{ is }Y(𝐱)(\sigma ):=Y(𝐱,\sigma )\text{ is }𝐂^{0+}\text{ and is, therefore, continuous.}\hfill \end{array}$$ (2D.20) So, $$Y(𝐱)\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{.}$$ (2D.21) * We now return to the definition of $`(𝐱)`$ in terms of $`𝒴_1(𝐱)`$ by Eq. (2B.46). From Lemma 6(i) when $`\mathrm{}`$ is $`1`$, and from statement (2D.1) concerning $`𝒴_1(𝐱)`$ being $`𝐂^{0+}`$ on its domain $`\overline{}(𝐱)`$, note that the factors in the numerator of the integrand in Eq. (2B.46) have the following properties: $$\begin{array}{c}𝒴_1(\sigma ^{})[W_2(\sigma ^{})]^TJ\text{ is defined for all }\sigma ^{}\overline{}(𝐱)\hfill \\ \text{and obeys a Hölder condition on }\overline{}(𝐱);\hfill \end{array}$$ (2D.22) and $$\begin{array}{c}[\nu ^\pm (\sigma ^{})]^1\text{ is }H(1/2)\text{ on each closed subinterval of }(𝐱)\hfill \\ \text{and converges to zero as }\sigma ^{}r\text{ and as }\sigma ^{}s.\hfill \end{array}$$ (2D.23) Also, recall that $$\begin{array}{c}\nu (\tau )\text{ is that branch of }(\tau r_0)^{1/2}(\tau s_0)^{1/2}(\tau r)^{1/2}(\tau s)^{1/2}\hfill \\ \text{which has the cut }\overline{}(𝐱)\text{ and the value }1\text{ at }\tau =\mathrm{}\text{.}\hfill \end{array}$$ (2D.24) Several theorems on Cauchy intgrals near the end points of the lines of inegration are given in Ref. 13, Sec. 29, Ch. 4. In particular, by applying Muskelishvili’s Eq. (29.4) to our Eq. (2B.46), one obtains the following conclusion from the above statements (2D.22) and (2D.23): $$\nu (\tau )^1(\tau )\text{ converges as }\tau r\text{ and as }\tau s.$$ (2D.25) Moreover, by applying Muskelishvili’s Eqs. (29.5) and (29.6) to our Eq. (2B.46), one obtains the following conclusion from the above statements (2D.22) to (2D.24): $$(\tau )\text{ converges as }\tau r_0\text{ and as }\tau s_0.$$ (2D.26) * From the above statements (2D.2), (2D.3), (2D.18), (2D.21), (2D.25) and (2D.26), all of the defining conditions for a solution of the HHP corresponding to $`(𝐯,^M,𝐱)`$ are satisfied by $`(𝐱)`$ as defined in terms of $`𝒴_1(𝐱)`$ by Eq. (2B.46). End of proof. We already know from Thm. 3(iv) that there is not more than one solution of the HHP corresponding to $`(𝐯,^M,𝐱)`$. ###### COROLLARY 12 (Uniqueness of $`𝒴_1(𝐱)`$) For each $`𝐯K^{1+}`$ and $`𝐱D`$, there is not more than one $`2\times 1`$ matrix function $`𝒴_1(𝐱)`$ such that $`\mathrm{d}om𝒴_1(𝐱)`$ $`=`$ $`\overline{}(𝐱),`$ (2D.27) $`𝒴_1(𝐱)`$ is $`𝐂^{0+}`$ (2D.28) and $`𝒴_1(𝐱,\sigma ):=𝒴_1(𝐱)(\sigma )`$ satisfies the singular integral equation (2B.49) for all $`\sigma (𝐱)`$. Proof: Suppose that $`𝒴_1(𝐱)`$ and $`𝒴_1^{}(𝐱)`$ are $`2\times 1`$ matrix functions, both of which have domain $`\overline{}(𝐱)`$, are $`𝐂^{0+}`$ and satisfy Eq. (2B.49) for all $`\sigma (𝐱)`$ and for the same given $`𝐯K^{1+}`$. Let $`(𝐱)`$ and $`^{}(𝐱)`$ be the $`2\times 2`$ matrix functions with domain $`C\overline{}(𝐱)`$ that are defined in terms of $`𝒴_1(𝐱)`$ and $`𝒴_1^{}(𝐱)`$, respectively, by Eq. (2B.46). Then, from the preceding Thm. 11, $`(𝐱)`$ and $`^{}(𝐱)`$ are both solutions of the HHP corresponding to $`(𝐯,^M,𝐱)`$; and, therefore, from Thm. 3(iv), $$(𝐱)=^{}(𝐱);$$ (2D.29) and, from Eq. (2D.13) in the proof of Thm. 11 and from statements (2D.27) and (2D.28), $$𝒴_1(𝐱)=𝒴_1^{}(𝐱).$$ (2D.30) End of proof. ## 3 A Fredholm integral equation of the second kind that is equivalent to the Alekseev-type singular integral equation when $`𝐯K^{2+}`$ If $`𝐯K^{1+}`$ and the particular solution $`𝒴_1(𝐱)`$ of Eq. (2B.49) that has a $`𝐂^{0+}`$ extension to $`\overline{}(𝐱)`$ exists, then it can be shown that $`𝒴_1(𝐱)`$ is also a solution of a Fredholm integral equation of the second kind. ### A. Derivation of Fredholm equation from Alekseev-type equation We shall employ a variant of the Poincaré-Bertrand commutator theorem. Suppose that $`L`$ is a smooth oriented line or contour in $`C\{\mathrm{}\}`$ and $`\varphi `$ is a complex-valued function whose domain is $`L\times L`$ and which obeys a Hölder condition on $`L\times L`$. Then the conventional Poincaré-Bertrand theorem asserts $`[{\displaystyle \frac{1}{\pi i}}{\displaystyle _L}𝑑\tau ^{\prime \prime },{\displaystyle \frac{1}{\pi i}}{\displaystyle _L}𝑑\tau ^{}]{\displaystyle \frac{\varphi (\tau ^{},\tau ^{\prime \prime })}{(\tau ^{\prime \prime }\tau )(\tau ^{}\tau ^{\prime \prime })}}`$ $`=`$ $`\varphi (\tau ,\tau )\text{ for all }\tau L`$ (3A.1) minus its end points, where the above bracketed expression is the commutator of the path integral operators. We are, of course, concerned here only with the case $`L=\overline{}(𝐱)`$; and our variant asserts that, for any function $`\varphi `$ which is $`𝐂^{0+}`$ on $`\overline{}(𝐱)^2`$, $`[{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime }),{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(\sigma ^{})^1]{\displaystyle \frac{\varphi (\sigma ^{},\sigma ^{\prime \prime })}{(\sigma ^{\prime \prime }\sigma )(\sigma ^{}\sigma ^{\prime \prime })}}`$ $`=`$ $`\varphi (\sigma ,\sigma )\text{ for all}`$ (3A.2) $`\sigma (𝐱),`$ or, alternatively, $`[{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime })^1,{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(\sigma ^{})]{\displaystyle \frac{\varphi (\sigma ^{},\sigma ^{\prime \prime })}{(\sigma ^{\prime \prime }\sigma )(\sigma ^{}\sigma ^{\prime \prime })}}`$ $`=`$ $`\varphi (\sigma ,\sigma )\text{ for all}`$ (3A.3) $`\sigma (𝐱),`$ We shall not supply the proof here, as an elegant and thorough proof of the Poincaré-Bertrand theorem (3A.1) is given by Sec. 23 of Muskhelishvili’s treatise, and what we have done is to construct proofs of (3A.2) and (3A.3) that parallel his proof step by step. We shall now apply Eq. (3A.2) to the Alekseev-type equation (2B.49), which we express in the form $$\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(\sigma ^{})^1\frac{𝒴_1(\sigma ^{})d_{21}(\sigma ^{},\sigma ^{\prime \prime })}{\sigma ^{}\sigma ^{\prime \prime }}=W_1(\sigma ^{\prime \prime })\text{ for all }\sigma ^{\prime \prime }(𝐱),$$ (3A.4) where, for all $`\sigma \overline{}(𝐱)`$ and $`\sigma ^{}\overline{}(𝐱)`$, $$d_{21}(\sigma ^{},\sigma ):=W_{22}(\sigma ^{})W_{11}(\sigma )W_{12}(\sigma ^{})W_{21}(\sigma ).$$ (3A.5) We suppose that, for a given $`𝐯K^{1+}`$ and $`𝐱D`$, a solution $`𝒴_1(𝐱)`$ of the Alekseev-type equation (2B.49) exists and is $`𝐂^{0+}`$ on $`\overline{(𝐱)}`$. Then the product $`𝒴_1(𝐱)d_{21}`$ is $`𝐂^{0+}`$ on $`\overline{}(𝐱)^2`$. Also, $`detW(\sigma )=d_{21}(\sigma ,\sigma )=1`$. Therefore, upon multiplying both sides of Eq. (3A.4) by $`(\sigma ^{\prime \prime }\sigma )^1`$ and then applying the PV integral operator $$\frac{1}{\pi i}_\overline{}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime }),$$ Eq. (3A.2) gives us $$𝒴_1(\sigma )\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(\sigma ^{})^1𝒴_1(\sigma ^{})K_{21}(\sigma ^{},\sigma )=U_1(\sigma ),$$ (3A.6) where, for each $`\sigma \overline{}(𝐱)\{r,s\}`$, $$U_1(\sigma ):=\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(\sigma ^{})\frac{W_1(\sigma ^{})}{\sigma ^{}\sigma };$$ (3A.7) and, for each $`(\sigma ^{},\sigma )\overline{}(𝐱)\times [\overline{}(𝐱)\{r,s\}]`$, $$K_{21}(\sigma ^{},\sigma ):=\frac{1}{\pi i}_\overline{}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime })\frac{d_{21}(\sigma ^{},\sigma ^{\prime \prime })}{(\sigma ^{\prime \prime }\sigma )(\sigma ^{}\sigma ^{\prime \prime })}.$$ (3A.8) So far, we have only established that Eq. (3A.6) holds for all $`\sigma (𝐱)`$. However, using the expressions $`U_1(\sigma )`$ $`=`$ $`W_1(\sigma ){\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(\sigma ^{}){\displaystyle \frac{W_1(\sigma ^{})W_1(\sigma )}{\sigma ^{}\sigma }},`$ (3A.9) $`K_{21}(\sigma ^{},\sigma )`$ $`=`$ $`k_{21}(\sigma ^{},\sigma ){\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime })\left[{\displaystyle \frac{k_{21}(\sigma ^{},\sigma ^{\prime \prime })k_{21}(\sigma ^{},\sigma )}{\sigma ^{\prime \prime }\sigma }}\right],`$ (3A.10) where $$k_{21}(\sigma ^{},\sigma ):=\frac{d_{21}(\sigma ^{},\sigma )1}{\sigma ^{}\sigma },$$ (3A.11) it is not difficult to prove the following lemma. ###### LEMMA 13 (Properties of $`U_1(𝐱)`$ and $`K_{21}(𝐱)`$) For each $`𝐱D`$ and $`𝐯K^{\mathrm{}}`$, $`U_1(𝐱)`$ is $`𝐂^{n1}`$ if $`\mathrm{}`$ is $`n`$ and $`n2`$, is $`𝐂^{(n1)+}`$ if $`\mathrm{}`$ is $`n+`$ and $`n2`$, are $`𝐂^{\mathrm{}}`$ if $`\mathrm{}`$ is $`\mathrm{}`$ and is $`𝐂^{an}`$ if $`\mathrm{}`$ is ‘an’; and $`K_{21}(𝐱)`$ is $`𝐂^{n2}`$ if $`\mathrm{}`$ is $`n`$ and $`n2`$, is $`𝐂^{(n2)+}`$ if $`\mathrm{}`$ is $`n+`$ and $`n2`$, is $`𝐂^{\mathrm{}}`$ if $`\mathrm{}`$ is $`\mathrm{}`$, and is $`𝐂^{an}`$ if $`\mathrm{}`$ is ‘an’. If $`𝐯K^{1+}`$, then $`U_1(𝐱)`$ is $`𝐂^{0+}`$ and $`K_{21}(𝐱)`$ is also $`𝐂^{0+}`$ {but, as we recall, its domain is only $`\overline{}(𝐱)\times [\overline{}(𝐱)\{r,s\}]`$}. From this it follows that $$U_1(𝐱)\text{ is continuous on }\overline{}(𝐱)$$ (3A.12) and $$K_{21}(𝐱)\text{ is continuous on }\mathrm{d}omK_{21}(𝐱).$$ (3A.13) Moreover, Lem. 7 remains valid if $`S`$ is a closed or a semi-closed subinterval of $`R^1`$. Therefore, from (3A.13) and Lem. 7(i) \[with a closed or a semi-closed $`SR^1`$\], the integral in Eq. (3A.6) is a continuous function of $`\sigma `$ throughout $`\overline{}(𝐱)`$ if $`𝐯K^2`$, and throughout $`\overline{}(𝐱)\{r,s\}`$ if $`𝐯K^2`$; and it then follows from the fact that $$𝒴_1(𝐱)\text{ is continuous on }\overline{}(𝐱)$$ (3A.14) and from (3A.12) that Eq. (3A.6) holds for all $`\sigma \overline{}(𝐱)`$ if $`𝐯K^2`$, and for all $`\sigma \overline{}(𝐱)\{r,s\}`$ if $`𝐯K^2`$. Thus, we have the following theorem. ###### THEOREM 14 (Fredholm equation) Suppose that, for a given $`𝐯K^{1+}`$ and $`𝐱D`$, a solution $`𝒴_1(𝐱)`$ of the Alekseev-type equation (2B.49) exists and is $`𝐂^{0+}`$ on $`\overline{}(𝐱)`$. Then the Fredholm equation (3A.6) holds for all $`\sigma \overline{}(𝐱)`$ if $`𝐯K^2`$ and for all $`\sigma \overline{I}(𝐱)\{r,s\}`$ if $`𝐯K^2`$. ### B. Equivalence of Alekseev-type equation and Fredholm equation when $`𝐯K^{2+}`$ The Fredholm equation (3A.6) generally has a singular kernel and is generally not equivalent to the Alekseev-type equation (2B.49). In this section we shall restrict our attention to the case $`𝐯K^{2+}`$. ###### THEOREM 15 (Alekseev-Fredholm equivalence theorem) Suppose $`𝐯K^{2+}`$, $`𝐱D`$ and $`𝒴_1(𝐱)`$ is a $`2\times 1`$ column matrix function with domain $`\overline{}(𝐱)`$. Then $`U_1(𝐱)`$ is $`𝐂^{1+}`$ and $`K_{21}(𝐱)`$ is $`𝐂^{0+}`$. Also, the following two statements are equivalent to one another: * $`𝒴_1(𝐱)`$ is $`𝐂^{0+}`$ and is the solution of Eq. (2B.49) for all $`\sigma (𝐱)`$. * $`𝒴_1(𝐱)`$ is summable over $`\overline{}(𝐱)`$ and is a solution of Eq. (3A.6) for all $`\sigma \overline{}(𝐱)`$. Proof: From Lem. 13, $`U_1(𝐱)`$ is $`𝐂^{1+}`$ and $`K_{21}(𝐱)`$ is $`𝐂^{0+}`$; and Thm. 14 already asserts that statement (i) implies statement (ii). It remains only to prove that statement (ii) implies statement (i). Grant statement (ii). Since $`U_1(𝐱)`$ is $`𝐂^{1+}`$ and $`K_{21}(𝐱)`$ is $`𝐂^{0+}`$ on $`\overline{}(𝐱)`$ and since $`𝒴_1(𝐱)`$ is summable over $`\overline{}(𝐱)`$, Eq. (3A.6) and Lem. 7(iv) yield $$𝒴_1(𝐱)\text{ is }𝐂^{0+}\text{ on }\overline{}(𝐱).$$ (3B.1) Next, using the Poincaré-Beltrami variant, one deduces the following equivalent of the Fredholm equation (3A.6): $$𝒴_1(\sigma )+\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(\sigma ^{})\frac{\psi (\sigma ^{})+W_1(\sigma ^{})}{\sigma ^{}\sigma }=0,$$ (3B.2) where $$\psi (\sigma ):=\frac{1}{\pi i}𝑑\sigma ^{}\nu ^+(\sigma ^{})^1𝒴_1(\sigma ^{})k_{21}(\sigma ^{},\sigma ).$$ (3B.3) From Lem. 7(iv) and (3B.1), $$\psi \text{ is }𝐂^{0+}\text{ on }\overline{}(𝐱).$$ (3B.4) Next, after replacing ‘$`\sigma `$’ by ‘$`\sigma ^{\prime \prime }`$’ in Eq. (3B.2) and then applying the operator $$\frac{1}{\pi i}_\overline{}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime })^1\frac{1}{\sigma ^{\prime \prime }\sigma },$$ one finds that $$\frac{1}{\pi i}_\overline{}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime })^1\frac{𝒴_1(\sigma ^{\prime \prime })}{\sigma ^{\prime \prime }\sigma }+\psi (\sigma )+W_1(\sigma )=0,$$ (3B.5) from which equation one can derive the Alekseev-type equation (2B.49). End of proof. Let us summarize the results given by Thm. 11 and Thm. 15 when $`𝐯K^{2+}`$. ###### THEOREM 16 (Summary) Suppose $`𝐯K^{2+}`$, $`𝐱D`$, and $`(𝐱)`$ and $`𝒴_1(𝐱)`$ are $`2\times 2`$ and $`2\times 1`$ matrix functions, respectively, such that $$\mathrm{d}om(𝐱)=C\overline{}(𝐱)\text{ and }\mathrm{d}om𝒴_1(𝐱)=\overline{}(𝐱).$$ (3B.6) Then the following three statements are equivalent to one another: * The function $`(𝐱)`$ is the solution of the HHP corresponding to $`(𝐯,^M,𝐱)`$, and $`𝒴_1(𝐱)`$ is the function whose restriction to $`(𝐱)`$ is defined by Eq. (2B.39) and whose extension to $`\overline{}(𝐱)`$ is then defined by Eqs. (2C.1) and (2C.2). \[The existence and uniqueness of this extension is asserted by Cor. 10.\] * The function $`𝒴_1(𝐱)`$ is $`𝐂^{0+}`$ and its restriction to $`(𝐱)`$ is a solution of the Alekseev-type equation (2B.49); and $`(𝐱)`$ is defined in terms of $`𝒴_1(𝐱)`$ by Eq. (2B.46). * The function $`𝒴_1(𝐱)`$ is summable over $`\overline{}(𝐱)`$ and is a solution of the Fredholm equation (3A.6) for all $`\sigma \overline{}(𝐱)`$. Proof: Directly from Thm. 11 and Thm. 15. End of proof. ###### COROLLARY 17 (Uniqueness of solutions) When $`𝐯K^{2+}`$, each of the solutions defined in (i), (ii) and (iii) of the preceding theorem is unique if it exists. Proof: This follows from the preceding theorem and the uniqueness theorem \[Thm. 3(iv)\] for the HHP. End of proof. ## 4 Existence and properties of the HHP solution $``$ when $`𝐯K^{2+}`$ ### A. Homogeneous equations, theorems, etc. By considering a homogeneous version of the Fredholm equation (3A.6), we found it possible to employ the Fredholm alternative theorem to establish the existence of the solution of the HHP corresponding to $`(𝐯,^M)`$ when $`𝐯K^{2+}`$. Dfn. of HHP<sub>0</sub> * The HHP that is defined as in Sec. 1 except that the condition (2) is replaced by the condition $$(𝐱,\mathrm{})=0\text{ (HHP}\text{0}\text{ condition)}$$ (4A.1) will be called the HHP<sub>0</sub> corresponding to $`(𝐯,_0,𝐱)`$. End of Dfn. Clearly, the $`2\times 2`$ matrix function $`(𝐱)`$ with the domain $`C\overline{}(𝐱)`$ and the value $`(𝐱,\tau )=0`$ for all $`\tau `$ in this domain is a solution of the HHP<sub>0</sub> corresponding to $`(𝐯,_0,𝐯)`$. It will be called the zero solution. Dfn. of equation number with attached subscript ‘$`0`$ * To each linear integral equation that occurs in these notes from Thm. 4 to Thm. 16, inclusive, and that has a term that is an integral whose integrand involves ‘$``$’, ‘$`^\pm `$’, ‘$`𝒴`$’ or ‘$`𝒴^{(i)}`$’ (or one of their columns), there corresponds a homogeneous integral equation that will be designated by the symbol that results when the subscript ‘$`0`$’ is attached to the equation number for the inhomogeneous integral equation. End of Dfn. Dfn. of theorem label (etc.) with attached subscript ‘$`0`$ * When a new valid assertion results from subjecting a labelled assertion to the following substitutions, that new valid assertion will bear the same label with an attached subscript ‘$`0`$’. + ‘HHP’ $``$ ‘HHP<sub>0</sub> + $`(𝐱,\mathrm{})=I`$$``$$`(𝐱,\mathrm{})=0`$’ in condition (2) of the HHP + each integral equation $``$ the corresponding homogeneous integral equation + each equation number for an integral equation $``$ the same equation number with an attached subscript ‘$`0`$’. End of Dfn. ### B. Only a zero solution of homogeneous equation For our immediate purpose, we shall need the following explicit version of Thm. 16<sub>0</sub>: ###### THEOREM 18 (Theorem 16<sub>0</sub>) Suppose $`𝐯K^{2+}`$, $`𝐱D`$, and $`(𝐱)`$ and $`𝒴_1(𝐱)`$ are $`2\times 2`$ and $`2\times 1`$ matrix functions, respectively, such that $$\mathrm{d}om(𝐱)=C\overline{}(𝐱)\text{ and }\mathrm{d}om𝒴_1(𝐱)=\overline{}(𝐱).$$ (4B.1) Then the following three statements are equivalent to one another: * The function $`(𝐱)`$ is a solution of the HHP<sub>0</sub> corresponding to $`(𝐯,^M,𝐱)`$; and $`𝒴_1(𝐱)`$ is the continuous function whose restriction to $`(𝐱)`$ is defined in terms of $`^\pm (𝐱)`$ by Eq. (2B.39), and whose existence and uniqueness are asserted by Cor. 10<sub>0</sub>. * The function $`𝒴_1(𝐱)`$ is $`𝐂^{0+}`$ and its restriction to $`(𝐱)`$ is a solution of Eq. (2B.49)<sub>0</sub>; and $`(𝐱)`$ is defined in terms of $`𝒴_1(𝐱)`$ by Eq. (2B.46)<sub>0</sub>. * The function $`𝒴_1(𝐱)`$ is summable over $`\overline{}(𝐱)`$ and is a solution of the homogeneous Fredholm integral equation (3A.6)<sub>0</sub> for all $`\sigma \overline{}(𝐱)`$. Proof: This theorem summarizes Thms. 11<sub>0</sub> and 15<sub>0</sub> for the case $`𝐯K^{2+}`$. End of proof. ###### THEOREM 19 (Only a zero solution of HHP<sub>0</sub>) For each $`𝐯K`$, $`_0𝒮_{}`$ and $`𝐱D`$, the only solution of the HHP<sub>0</sub> corresponding to $`(𝐯,_0,𝐱)`$ is its zero solution. Proof: The proof will be given in four parts: * From the hypothesis $`_0𝒮_{}`$, $$\left[_0(𝐱,\tau ^{})\right]^{}𝒜_0(𝐱,\tau )_0(𝐱,\tau )=𝒜_0(𝐱_0,\tau )\text{ for all }\tau C\overline{}(𝐱),$$ (4B.2) where $$𝒜_0(𝐱,\tau ):=(\tau z)\mathrm{\Omega }+\mathrm{\Omega }h_0(𝐱)\mathrm{\Omega }$$ (4B.3) and $`h_0(𝐱)`$ is computed from $`_0𝒮_{}`$ in the usual way.<sup>23</sup><sup>23</sup>23To prove Eq. (4B.2), one first shows that Eq. (1A.21) is equivalent to $`𝒜_0\mathrm{\Gamma }_0=\frac{1}{2}\mathrm{\Omega }dH_0\mathrm{\Omega }`$ and then uses (1A.7) to show that the differential of the left side of Eq. (4B.2) vanishes. The rest is obvious. Since $$h_0(𝐱_0):=h^M(𝐱_0)=\left(\begin{array}{cc}\rho _0^2& 0\\ 0& 1\end{array}\right)$$ (4B.4) in our gauge, $$𝒜_0(𝐱_0,\tau )=𝒜^M(𝐱_0,\tau ).$$ (4B.5) Equation (4B.2) is clearly expressible in the alternative form $`_0(𝐱,\tau )\left[𝒜^M(𝐱_0,\tau )\right]^1\left[_0(𝐱,\tau ^{})\right]^{}=\left[𝒜_0(𝐱,\tau )\right]^1`$ (4B.6) $`\text{for all }\tau C\overline{}(𝐱),`$ since $`[_0(𝐱,\tau )]^1`$ exists for all $`\tau C\overline{}(𝐱)`$, and $$\left[𝒜_0(𝐱,\tau )\right]^1=\frac{_0(𝐱,\tau )}{\rho ^2(\tau z)^2},$$ (4B.7) where $$_0(𝐱,\tau ):=h_0(𝐱)(\tau z)\mathrm{\Omega },$$ (4B.8) exists for all $`\tau C\{r,s\}`$. * Next, condition (3) in the definition of the HHP (and the HHP<sub>0</sub>) that is given in Sec. 1 asserts that $`^\pm (𝐱)`$ exist, and Eq. (3) is expressible in the form $$^\pm (𝐱,\sigma )=Y^{(i)}(\sigma )_0^\pm (𝐱,\sigma )[v^{(i)}(\sigma )]^1\text{ for each }i\{3,4\}\text{ and }\sigma (𝐱).$$ (4B.9) From the definition of the group $`K`$, $$[v^{(i)}(\sigma )]^1[𝒜^M(𝐱_0,\sigma )]^1[v^{(i)}(\sigma )^{}]^1=𝒜^M(𝐱_0,\sigma )^1\text{ for all }\sigma ^{(i)}\{r,s\}.$$ (4B.10) Therefore, from Eqs. (4B.9), (4B.10) and (4B.6), $`^\pm (𝐱,\sigma )[𝒜^M(𝐱_0,\sigma )]^1[^{}(𝐱,\sigma )]^{}`$ $`=`$ $`Y(𝐱,\sigma )[𝒜_0(𝐱,\sigma )]^1Y(𝐱,\sigma )^{}`$ (4B.11) $`\text{for all }\sigma (𝐱);`$ or, equivalently, with the aid of Eqs. (4B.7), (4B.8) and (4B.4), $`\left[{\displaystyle \frac{\rho ^2(\sigma z)^2}{\rho _0^2(\sigma z_0)^2}}\right]^\pm (𝐱,\sigma )^M(\sigma )[^{}(𝐱,\sigma )]^{}=`$ (4B.12) $`Y(𝐱,\sigma )_0(𝐱,\sigma )Y(𝐱,\sigma )^{}\text{ for all }\sigma (𝐱),`$ where $$^M(\tau ):=\left(\begin{array}{cc}\rho _0^2& i(\tau z_0)\\ i(\tau z_0)& 1\end{array}\right).$$ (4B.13) * Next, let $`Z(𝐱)`$ denote the function with the (tentative) domain $`C\overline{}(𝐱)`$ and the values $`Z(𝐱,\tau )`$ $`:=`$ $`Z(𝐱)(\tau )`$ (4B.14) $`:=`$ $`\nu (𝐱,\tau )^1(𝐱,\tau )^M(\tau )[\nu (𝐱,\tau ^{})^1(𝐱,\tau ^{})]^{}`$ $`\text{for all }\tau C\overline{}(𝐱),`$ where note that $$\nu (𝐱,\tau )^2=\frac{(\tau r)(\tau s)}{(\tau r_0)(\tau s_0)}=\frac{(\tau z)^2\rho ^2}{(\tau z_0)^2\rho _0^2}.$$ (4B.15) We again appeal to the trilogy of elementary theorems due to Riemann and Liouville.<sup>24</sup><sup>24</sup>24See Refs. 14, 15 and 16. Using these, we shall define an extension of $`Z(𝐱)`$, and we shall let $`Z(𝐱)`$ denote this extension as well. From condition (1) in the definition of the HHP (and the HHP<sub>0</sub>), and from Eqs. (4B.14), (4B.13) and (4A.1), $$Z(𝐱,\tau )\text{ is a holomorphic function of }\tau \text{ throughout }C\overline{}(𝐱),$$ (4B.16) and $$Z(𝐱,\mathrm{})=0.$$ (4B.17) Let ($`\mathrm{I}m\zeta >0`$) $$Z^\pm (𝐱,\sigma ):=\underset{\zeta 0}{lim}Z(𝐱,\sigma \pm \zeta )\text{ for all }\sigma (𝐱),$$ (4B.18) which exist according to condition (3) in the definition of the HHP (and the HHP<sub>0</sub>). Then, from Eqs. (4B.14), (4B.15) and (4B.12), $$Z^+(𝐱,\sigma )=Z^{}(𝐱,\sigma )=Y(𝐱,\sigma )_0(𝐱,\sigma )Y(𝐱,\sigma )^{}\text{ for all }\sigma (𝐱).$$ (4B.19) The above equation permits us to define a single valued extension of $`Z(𝐱)`$ to the domain $`C\{r,s,r_0,s_0\}`$ by letting $$Z(𝐱,\sigma ):=Z^\pm (𝐱,\sigma )=Y(𝐱,\sigma )_0(𝐱,\sigma )Y(𝐱,\sigma )^{}\text{ for all }\sigma (𝐱),$$ (4B.20) whereupon, from (4B.16), (4B.20) and the theorem on analytic continuation across an arc, $$Z(𝐱,\tau )\text{ is a holomorphic function of }\tau \text{ throughout }C\{r,s,r_0,s_0\}.$$ (4B.21) We next apply condition (4) in the definition of the HHP (and HHP<sub>0</sub>). Since, according to condition (4), $`\nu (𝐱)^1(𝐱)`$ and $`Y(𝐱)`$ are both bounded at $`𝐱`$, Eqs. (4B.14) and (4B.20) yield $$\begin{array}{c}\text{There exists a positive real number }M_1(𝐱)\text{ such that }\hfill \\ Z(𝐱,\tau )<M_1(𝐱)\text{ as }\tau r\text{ and as }\tau s\hfill \\ \text{through any sequence of points in }C\{r,s,r_0,s_0\}.\hfill \end{array}$$ (4B.22) Since $`(𝐱)`$ and $`Y(𝐱)`$ are both bounded at $`𝐱_0`$, Eqs. (4B.14), (4B.15) and (4B.20) yield $$\begin{array}{c}\text{There exists a positive real number }M_2(𝐱)\text{ such that }\hfill \\ (\tau r_0)(\tau s_0)Z(𝐱,\tau )<M_2(𝐱)\text{as }\tau r_0\text{ and as }\tau s_0\hfill \\ \text{through any sequence of points in }C\{r,s,r_0,s_0\}.\hfill \end{array}$$ (4B.23) However, since $`Y(𝐱)`$ is bounded at $`𝐱_0`$, Eq. (4B.20) yields $$\begin{array}{c}\text{There exists a positive real number }M_3(𝐱)\text{ such that }\hfill \\ Z(𝐱,\sigma )<M_3(𝐱)\text{ as }\sigma r_0\text{ and as }\sigma s_0\hfill \\ \text{through any sequence of points in }(𝐱).\hfill \end{array}$$ (4B.24) The theorem on isolated singularities, together with statements (4B.21) to (4B.24), now informs us that $$Z(𝐱)\text{ has a holomorphic extension [which we also denote by }Z(𝐱)\text{] to C,}$$ (4B.25) whereupon Eq. (4B.17) and the (generalized) theorem of Liouville yield $$Z(𝐱,\tau )=0\text{ for all }\tau C.$$ (4B.26) * Putting (4B.14) and (4B.26) together, one obtains $$(𝐱,\sigma )^M(\sigma )(𝐱,\sigma )^{}=0\text{ for all }\sigma C\overline{}(𝐱).$$ (4B.27) Note from Eq. (4B.13), $`^M(\sigma )`$ is hermitian, $$\begin{array}{ccc}\hfill \mathrm{t}r^M(\sigma )& =& 1+\rho _0^2\text{ and }\hfill \\ \hfill det^M(\sigma )& =& (s_0\sigma )(\sigma r_0).\hfill \end{array}$$ (4B.28) Recall that $`\|r,r_0\|<\|s,s_0\|`$ for any type A triple $`(x_0,x_1,x_2)`$; and it is clear that $$^M(\sigma )\text{ is hermitian and positive definite for all }\|r,r_0\|<\sigma <\|s,s_0\|.$$ (4B.29) Therefore, Eq. (4B.27) implies $$(𝐱,\sigma )=0\text{ for all }\sigma \text{ such that }\|r,r_0\|<\sigma <\|s,s_0\|.$$ (4B.30) However, $`(𝐱,\tau )`$ is a holomorphic function of $`\tau `$ throughout $`C\overline{}(𝐱)`$, and this domain contains the open interval between $`\|r,r_0\|`$ and $`\|s,s_0\|`$. So, $$(𝐱,\tau )=0\text{ for all }\tau C\overline{}(𝐱).$$ (4B.31) End of proof. ###### THEOREM 20 (Only a zero solution of (3A.6)<sub>0</sub>) The only solution of the homogeneous Fredholm integral equation of the second kind Eq. (3A.6)<sub>0</sub> is its zero solution. Proof: Let $`𝒴_1(𝐱)`$, with domain $`\overline{}(𝐱)`$, denote a solution of Eq. (3A.6)<sub>0</sub>; and let $`(𝐱)`$, with domain $`C\overline{}(𝐱)`$, be defined in terms of $`𝒴_1(𝐱)`$ by Eq. (2B.46)<sub>0</sub>. Using Thm. 18, one obtains $$(𝐱)\text{ is a solution of the HHP}\text{0}\text{ corresponding to }(𝐯,^M,𝐱),$$ (4B.32) whereupon Thm. 19 delivers $$(𝐱,\tau )=0\text{ for all }\tau C\overline{}(𝐱).$$ (4B.33) It follows that $$^\pm (𝐱,\sigma )=0\text{ for all }\sigma (𝐱),$$ (4B.34) whereupon, from Thm. 19(i), Eq. (2B.39) and the continuity of $`𝒴_1(𝐱)`$, $$𝒴_1(𝐱,\sigma )=0\text{ for all }\sigma \overline{}(𝐱).$$ (4B.35) End of proof. ### C. Existence and uniqueness of HHP solution At this point, we note that Eq. (3A.6) is a regular Fredholm equation in disguise when $`𝐯K^{2+}`$. In integrals such as those in Thm. 5, it is sometimes useful to introduce a new variable of integration for the purpose of getting rid of the singularities of the integrands at $`\sigma ^{}\{r,s,r_0,s_0\}`$. This is especially important when one has to consider derivatives of the integrals with respect to $`r`$ and $`s`$. Dfns. of $`\mathrm{\Theta }`$, $`\theta (𝐱)`$ and $`\sigma (𝐱)`$ * Let $`\mathrm{\Theta }`$ denote that union of arcs $$\mathrm{\Theta }:=[0,\frac{\pi }{2}]+[\pi ,\frac{3\pi }{2}]$$ (4C.1) whose assigned orientations are in the direction of increasing $`\theta [0,\pi /2]`$ and $`\theta [\pi ,3\pi /2]`$. For each $`𝐱D`$, let $$𝜽(𝐱):\overline{}(𝐱)\mathrm{\Theta }$$ (4C.2) be a mapping such that $$𝜽(𝐱)(\sigma ):=𝜽(𝐱,\sigma ),$$ (4C.3) where $$\begin{array}{c}\hfill 0𝜽(𝐱,\sigma )\frac{\pi }{2}\text{ and }\mathrm{cos}[2𝜽(𝐱,\sigma )]:=\frac{2\sigma (r_0+r)}{r_0r}\\ \hfill \text{when }\sigma \overline{}^{(3)}(𝐱)\end{array}$$ (4C.4) and $$\begin{array}{c}\hfill \pi 𝜽(𝐱,\sigma )\frac{3\pi }{2}\text{ and }\mathrm{cos}[2𝜽(𝐱,\sigma )]:=\frac{2\sigma (s_0+s)}{s_0s}\\ \hfill \text{when }\sigma \overline{}^{(4)}(𝐱).\end{array}$$ (4C.5) Also let $$𝝈(𝐱):\mathrm{\Theta }\overline{}(𝐱)$$ (4C.6) be a mapping such that $$𝝈(𝐱)(\theta ):=𝝈(𝐱,\theta ),$$ (4C.7) where $`𝝈(𝐱,\theta )`$ $`:=`$ $`r_0\mathrm{cos}^2\theta +r\mathrm{sin}^2\theta \text{ when }\theta [0,{\displaystyle \frac{\pi }{2}}]`$ (4C.8) and $`𝝈(𝐱,\theta )`$ $`:=`$ $`s_0\mathrm{cos}^2\theta +s\mathrm{sin}^2\theta \text{ when }\theta [\pi ,{\displaystyle \frac{3\pi }{2}}].`$ (4C.9) End of Dfn. The mapping $`𝜽(𝐱)`$ is monotonic and is a continuous bijection (one-to-one and onto) of $`(𝐱)`$ onto $`\mathrm{\Theta }`$, and $`𝝈(𝐱)`$ is its inverse mapping. Moreover, $`𝝈(𝐱)`$ is analytic \[which means that it has an analytic extension to an open subset of $`R^1`$\]. Note, in particular, that $`\sqrt{{\displaystyle \frac{s𝝈(𝐱,\theta ^{})}{s_0𝝈(𝐱,\theta ^{})}}}\text{ is a real positive-valued analytic}`$ $`\text{function of }(𝐱,\theta ^{})\text{ on }D\times [0,{\displaystyle \frac{\pi }{2}}]`$ (4C.10) and $`\sqrt{{\displaystyle \frac{𝝈(𝐱,\theta ^{})r}{𝝈(𝐱,\theta ^{})r_0}}}\text{ is a real positive-valued analytic}`$ $`\text{function of }(𝐱,\theta ^{})\text{ on }D\times [\pi ,{\displaystyle \frac{3\pi }{2}}],`$ (4C.11) since the left and right cuts are assumed not to overlap. The following equation is equivalent to Eq. (3A.6) and has a $`𝐂^{0+}`$ kernel and a $`𝐂^{1+}`$ inhomogeneous term: $`y_1(𝐱,\theta ){\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}y_1(𝐱,\theta ^{})\kappa _{21}(𝐱,\theta ^{},\theta )`$ (4C.12) $`=`$ $`u_1(𝐱,\theta )\text{ for all }\theta \mathrm{\Theta }:=[0,\pi /2][\pi ,3\pi /2],`$ where $`y_1(𝐱,\theta )`$ $`:=`$ $`𝒴_1(𝐱,\sigma (𝐱,\theta )),`$ (4C.13) $`u_1(𝐱,\theta )`$ $`:=`$ $`U_1(𝐱,\sigma (𝐱,\theta )),`$ (4C.14) $`\kappa _{21}(𝐱,\theta ^{},\theta )`$ $`:=`$ $`q(𝐱,\theta ^{})K_{21}(𝐱,\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta ))`$ (4C.15) and $$q(𝐱,\theta ):=\{\begin{array}{c}(r_0r)\mathrm{cos}^2\theta \sqrt{\frac{s\sigma (𝐱,\theta )}{s_0\sigma (𝐱,\theta )}}\text{ when }\theta [0,\pi /2],\hfill \\ (s_0s)\mathrm{cos}^2\theta \sqrt{\frac{\sigma (𝐱,\theta )r}{\sigma (𝐱,\theta )r_0}}\text{ when }\theta [\pi ,3\pi /2].\hfill \end{array}$$ (4C.16) Equations (3A.9) and (3A.10) are expressible in the following forms, in which $`𝐱`$ and $`𝐱_0`$ are no longer suppressed: $`U_1(𝐱,\sigma )`$ $`=`$ $`W_1(\sigma ){\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}p(𝐱,\theta ^{}){\displaystyle \frac{W_1(\sigma (𝐱,\theta ^{}))W_1(\sigma )}{\sigma (𝐱,\theta ^{})\sigma }},`$ (4C.17) $`K_{21}(𝐱,\sigma ^{},\sigma )`$ $`=`$ $`k_{21}(\sigma ^{},\sigma ){\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{\prime \prime }p(𝐱,\theta ^{\prime \prime }){\displaystyle \frac{k_{21}(\sigma ^{},\sigma (𝐱,\theta ^{\prime \prime }))k_{21}(\sigma ^{},\sigma )}{\sigma (𝐱,\theta ^{\prime \prime })\sigma }},`$ (4C.18) where $$p(𝐱,\theta ):=\{\begin{array}{c}(r_0r)\mathrm{sin}^2\theta \sqrt{\frac{s_0\sigma (𝐱,\theta )}{s\sigma (𝐱,\theta )}}\text{ when }\theta [0,\pi /2],\hfill \\ (s_0s)\mathrm{sin}^2\theta \sqrt{\frac{\sigma (𝐱,\theta )r_0}{\sigma (𝐱,\theta )r}}\text{ when }\theta [\pi ,3\pi /2].\hfill \end{array}$$ (4C.19) ###### THEOREM 21 (Fredholm determinant not zero) The Fredholm determinant corresponding to the kernel $`\kappa _{21}(𝐱)`$ is not zero. Therefore, there exists exactly one solution of Eq. (4C.12) for each given $`𝐯K^{2+}`$ and $`𝐱D`$; or, equivalently, there exists exactly one solution of Eq. (3A.6) for each given $`𝐯K^{2+}`$ and $`𝐱D`$. Proof: This follows from Thm. 20 and the Fredholm alternative. End of proof. Thus, in summation, we have the following theorem: ###### THEOREM 22 (Existence and uniqueness of HHP solution) * If $`𝐯K^{2+}`$, then the HHP<sub>0</sub> corresponding to $`(𝐯,^M,𝐱)`$ is equivalent to the homogeneous Fredholm equation of the second kind that is obtained from Eq. (3A.6) by deleting the term $`U_1(\sigma )`$, provided that the term $`W_1(\sigma )`$ is also deleted from the expression (2B.48) for $`𝒴_2(\sigma )`$. * For any given $`𝐱D`$, and $`𝐯K`$, the HHP<sub>0</sub> corresponding to $`(v,^M,𝐱)`$ has the unique solution $`(𝐱,\tau )=0`$ for all $`\tau C(𝐱)`$. * Therefore, if $`𝐯K^{2+}`$, the only solution of the homogeneous Fredholm equation is the zero solution. Hence from the Fredholm alternative theorem, the inhomogeneous Fredholm equation (3A.6) has exactly one solution. We conclude that there exists one and only one solution of the HHP corresponding to $`(𝐯,^M)`$ when $`𝐯K^{2+}`$. Proof: Directly from Thms. 18, 19, 20 and 21. End of proof. ### D. The $`2\times 2`$ matrix $`H(𝐱)`$ associated with each solution of the HHP corresponding to $`(𝐯,_0,𝐱)`$ when $`𝐯K`$ ###### THEOREM 23 (Properties of $`H(𝐱)`$ and $`h(𝐱)`$) For each $`𝐯K`$, $`_0𝒮_{}`$, $`𝐱D`$ and solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,_0,𝐱)`$, there exists exactly one $`2\times 2`$ matrix $`H(𝐱)`$ such that $`(𝐱,\tau )`$ $`=`$ $`I+(2\tau )^1\left[H(𝐱)H^M(𝐱_0)\right]\mathrm{\Omega }+O(\tau ^2)`$ (4D.1) $`\text{in at least one neighborhood of }\tau =\mathrm{}.`$ Moreover, $`H(𝐱_0)`$ $`=`$ $`H^M(𝐱_0),`$ (4D.2) $`H(𝐱)H(𝐱)^T`$ $`=`$ $`2z\mathrm{\Omega },`$ (4D.3) $`h(𝐱)`$ $`:=`$ $`\mathrm{R}eH(𝐱)\text{ is symmetric,}`$ (4D.4) and $$h(𝐱_0)=\left(\begin{array}{cc}\rho _0^2& 0\\ 0& 1\end{array}\right).$$ (4D.5) Proof: From conditions (1) and (2) in the definition of the HHP, there exists exactly one $`2\times 2`$ matrix $`B(𝐱)`$ such that $`(𝐱,\tau )`$ $`=`$ $`I+(2\tau )^1B(𝐱)+O(\tau ^2)`$ $`\text{in at least one neighborhood of }\tau =\mathrm{}.`$ Let $$H(𝐱):=H^M(𝐱_0)+B(𝐱)\mathrm{\Omega },$$ whereupon statement (4D.1) follows. From Thm. 3(v) \[Eq. (1C.13)\], $`B(𝐱_0)=0`$, whereupon Eq. (4D.2) follows. Next, from Thm. 3(iii), $`det(𝐱,\tau )`$ $`=`$ $`\nu (𝐱,\tau )`$ (4D.6) $`=`$ $`1+(2\tau )^1(r+sr_0s_0)+O(\tau ^2)`$ $`=`$ $`1+\tau ^1(zz_0)+O(\tau ^2)`$ $`\text{in at least one neighborhood of }\tau =\mathrm{}.`$ Moreover, from Eq. (1B.4), $$H^M(𝐱_0)[H^M(𝐱_0)]^T=2z_0\mathrm{\Omega }.$$ (4D.7) For any $`2\times 2`$ matrix $`M`$, $`M\mathrm{\Omega }M^T=\mathrm{\Omega }detM`$. In particular, $$(𝐱,\tau )\mathrm{\Omega }(𝐱,\tau )^T=\mathrm{\Omega }\nu (𝐱,\tau ).$$ (4D.8) The next step is to consider Eq. (4D.8) in at least one neighborhood of $`\tau =\mathrm{}`$ for which the expansions given by Eqs. (4D.1) and (4D.6) hold. The reader can then easily deduce Eq. (4D.3) by using Eq. (4D.7) and the relations $`\mathrm{\Omega }^T=\mathrm{\Omega }`$ and $`\mathrm{\Omega }^2=I`$. The statement (4D.4) follows from Eq. (4D.3) and the relation $`\mathrm{\Omega }^{}=\mathrm{\Omega }`$. Equation (4D.5) is derived from Eqs. (4D.2) and (1B.4). End of proof. ###### THEOREM 24 (Quadratic relation) For each $`𝐯K`$, $`_0𝒮_{}`$, $`𝐱D`$ and solution $`(𝐱)`$ of the HHP corresponding to $`(𝐯,_0,𝐱)`$, let $`h(𝐱)`$ be defined as in the preceding theorem, and let $$𝒜(𝐱,\tau )=(\tau z)\mathrm{\Omega }+\mathrm{\Omega }h(𝐱)\mathrm{\Omega }.$$ (4D.9) Then $$^{}(𝐱,\tau )𝒜(𝐱,\tau )(𝐱,\tau )=𝒜(𝐱_0,\tau )\text{ for all }\tau [C\overline{}(𝐱)]\{\mathrm{}\},$$ (4D.10) where $$^{}(𝐱,\tau ):=[(𝐱,\tau ^{})]^{}\text{ for all }\tau C\overline{}(𝐱).$$ (4D.11) Proof: Note that parts (1) and (2) in the proof of Thm. 19 remain valid here. For the sake of convenience, we repeat below Eq. (4B.12) from part (2) of that proof. $`[\nu (𝐱,\sigma )]^2^\pm (𝐱,\sigma )^M(\sigma )[^{}(𝐱,\sigma )]^{}`$ $`=`$ $`Y(𝐱,\sigma )_0(𝐱,\sigma )Y(𝐱,\sigma )^{}`$ (4D.12) $`\text{for all }\sigma (𝐱),`$ where $`^M(\tau )`$ $`:=`$ $`\left(\begin{array}{cc}\rho _0^2& i(\tau z_0)\\ i(\tau z_0)& 1\end{array}\right),`$ (4D.15) $`_0(𝐱,\tau )`$ $`:=`$ $`h_0(𝐱)(\tau z)\mathrm{\Omega }`$ (4D.16) $`=`$ $`[\rho ^2(\tau z)^2]𝒜_0(𝐱,\tau )^1,`$ $`\nu (𝐱,\tau )^2`$ $`=`$ $`{\displaystyle \frac{(\tau r)(\tau s)}{(\tau r_0)(\tau s_0)}}`$ (4D.17) $`=`$ $`{\displaystyle \frac{(\tau z)^2\rho ^2}{(\tau z_0)^2\rho _0^2}}.`$ Next, let $`Z(𝐱)`$ denote the function with the (tentative) domain $`[C\overline{}(𝐱)]\{\mathrm{}\}`$ and the values $$Z(𝐱,\tau ):=\nu (𝐱,\tau )^1(𝐱,\tau )^M(\tau )\left[\nu (𝐱,\tau ^{})^1(𝐱,\tau ^{})\right]^{}.$$ (4D.18) From conditions (1) and (2) in the definition of the HHP, and from Eqs. (4D.15) and (4D.18), $$\begin{array}{c}Z(𝐱,\tau )\text{ is a holomorphic function of }\tau \hfill \\ \text{throughout }[C\overline{}(𝐱)]\{\mathrm{}\}\hfill \\ \text{and has a simple pole at }\tau =\mathrm{}.\hfill \end{array}$$ (4D.19) Note that Eq. (4D.5) enables us to express (4D.15) in the form $$^M(\tau )=h(𝐱_0)(\tau z_0)\mathrm{\Omega }.$$ (4D.20) Also, note that Eqs. (4D.3) and (4D.4) imply that $$H(𝐱)+H(𝐱)^{}=2h(𝐱)+2z\mathrm{\Omega }$$ (4D.21) and that Eq. (4D.17) yields $`\nu (𝐱,\tau )^2`$ $`=`$ $`1+2\tau ^1(z_0z)+O(\tau ^2)`$ (4D.22) $`\text{in at least one neighborhood of }\tau =\mathrm{}.`$ Upon using the relation $`\nu (𝐱,\tau ^{})^{}=\nu (𝐱,\tau )`$ and upon inserting (4D.1), (4D.20) and (4D.22) into the right side of Eq. (4D.18), one obtains the following result with the aid of Eqs. (4D.2) and (4D.17): $`Z(𝐱,\tau )`$ $`=`$ $`(\tau z)\mathrm{\Omega }+h(𝐱)+O(\tau ^1)`$ (4D.23) $`\text{in at least one neighborhood of }\tau =\mathrm{}.`$ We again appeal to the trilogy of elementary theorems due to Riemann and Liouville.<sup>25</sup><sup>25</sup>25See Refs. 14, 15 and 16. We let $`Z^\pm (𝐱,\sigma )`$ be defined for all $`\sigma (𝐱)`$ by Eq. (4B.18), whereupon Eqs. (4D.12) and (4D.18) yield $$Z^+(𝐱,\sigma )=Z^{}(𝐱,\sigma )=Y(𝐱,\sigma )_0(𝐱,\sigma )Y(𝐱,\sigma )^{}\text{ for all }\sigma (𝐱).$$ (4D.24) The above equation permits us to define a single valued extension of $`Z(𝐱)`$ to the domain $`C\{r,s,r_0,s_0,\mathrm{}\}`$ by letting $$Z(𝐱,\sigma ):=Z^\pm (𝐱,\sigma )=Y(𝐱,\sigma )_0(𝐱,\sigma )Y(𝐱,\sigma )^{}\text{ for all }\sigma (𝐱),$$ (4D.25) whereupon (4D.19), (4D.25) and the theorem on analytic continuation across an arc tell us that $$\begin{array}{c}Z(𝐱,\tau )\text{ is a holomorphic function of }\tau \text{throughout }\hfill \\ C\{r,s,r_0,s_0,\mathrm{}\}\text{ and has a simple pole at }\tau =\mathrm{}.\hfill \end{array}$$ (4D.26) We next use condition (4) in the definition of the HHP, and we obtain the statements (4B.22), (4B.23) and (4B.24) exactly as we did in the proof of Thm. 19. The theorem on isolated singularities, together with the statements (4D.26), (4B.22), (4B.23) and (4B.24) now inform us that $$\begin{array}{c}Z(𝐱)\text{ has a holomorphic extension [which we also denote }\hfill \\ \text{by }Z(𝐱)\text{] to }C\{\mathrm{}\}\text{ and has a simple pole at }\tau =\mathrm{},\hfill \end{array}$$ (4D.27) whereupon Eq. (4D.23) and the theorem on entire functions that do not have an essential singularity at $`\tau =\mathrm{}`$ yield $$Z(𝐱,\tau )=(\tau z)\mathrm{\Omega }+h(𝐱)\text{ for all }\tau C\{\mathrm{}\}.$$ (4D.28) We are now close to completing our proof. From Thm. 3(iii), Eqs. (4D.18), (4D.15) and (4D.17), $$detZ(𝐱,\tau )=\rho ^2(\tau z)^2.$$ (4D.29) Therefore, from Eqs. (4D.9) and (4D.28), the matrix $`(\tau z)\mathrm{\Omega }+h(𝐱)`$ is invertible when $`\tau \{r,s,\mathrm{}\}`$, and $$\left[(\tau z)\mathrm{\Omega }+h(𝐱)\right]^1=\frac{𝒜(𝐱,\tau )}{\rho ^2(\tau z)^2}.$$ (4D.30) \[Above, we have used the fact that $`M^1=\mathrm{\Omega }M^T\mathrm{\Omega }/detM`$ for any invertible $`2\times 2`$ matrix $`M`$.\] One then obtains from Eqs. (4D.18), (4D.20), (4D.28) and (4D.30), $`(𝐱,\tau )[𝒜(𝐱_0,\tau )]^1^{}(𝐱,\tau )=𝒜(𝐱,\tau )^1`$ $`\text{for all }\tau [C\overline{}(𝐱)]\{\mathrm{}\},`$ whereupon the conclusion (4D.10) follows. End of proof. ###### THEOREM 25 (More properties of $`h(𝐱)`$) Grant the same premises as in the preceding two theorems, and let $`h(𝐱)`$ be defined as before. Then $$deth(𝐱)=\rho ^2$$ (4D.31) and $$h(𝐱)\text{ is positive definite }$$ (4D.32) as well as real and symetric. Proof: Since $`h(𝐱)`$ is symmetric $$det[h(𝐱)(\tau z)\mathrm{\Omega }]=deth(𝐱)(\tau z)^2.$$ Therefore, Eq. (4D.28) and (4D.29) imply that $`deth(𝐱)=\rho ^2`$. From Eqs. (4D.18), (4D.28) and (4D.17), $`Z(𝐱,\sigma )`$ $`=`$ $`{\displaystyle \frac{(\sigma r)(s\sigma )}{(\sigma r_0)(s_0\sigma )}}(𝐱,\sigma )^M(\sigma )(𝐱,\sigma )^{}`$ (4D.33) $`=`$ $`(\sigma z)\mathrm{\Omega }+h(𝐱)\text{ for all }\|r,r_0\|<\sigma <\|s,s_0\|.`$ Equation (4D.15) provides us with $$det^M(\sigma )=(s_0\sigma )(\sigma r_0)\text{ and }\mathrm{t}r^M(\sigma )=1+\rho _0^2.$$ (4D.34) Therefore, $$\frac{(\sigma r)(s\sigma )}{(\sigma r_0)(s_0\sigma )}^M(\sigma )$$ is a positive definite hermitian matrix when $`\|r,r_0\|<\sigma <\|s,s_0\|`$. Therefore, the left side of Eq. (4D.33) is a positive definite hermitian matrix when $`\|r,r_0\|<\sigma <\|s,s_0\|`$ and must, therefore, have a real positive trace when $`\|r,r_0\|<\sigma <\|s,s_0\|`$. So, $$\mathrm{t}r[(\sigma z)\mathrm{\Omega }+h(𝐱)]=\mathrm{t}rh(𝐱)>0;$$ (4D.35) and, since the determinant of $`h(𝐱)`$ is also positive, $`h(𝐱)`$ is positive definite. End of proof. We caution the reader that the HHP solution $``$ whose existence has been proved in this section when $`𝐯K^{2+}`$ is not necessarily a member of $`𝒮_{}`$; and $`H`$ as defined by Eq. (4D.1) is not necessarily a member of $`𝒮_H`$. However, as we shall prove in Sec. 6, $`𝒮_{}`$ and $`H𝒮_H`$ when $`𝐯K^3`$. To prepare for this proof, we shall now investigate the differentiability of $``$ and $`H`$ when $`𝐯K^3`$. ## 5 Derivatives of $``$ and $`H`$ when $`𝐯K^3`$ ### A. Fredholm equation solution $`𝒴_1`$ corresponding to $`𝐯K^3`$ We again refer the reader to the mappings $`𝜽(𝐱):\overline{}(𝐱)\mathrm{\Theta }`$ and $`𝝈(𝐱):\mathrm{\Theta }\overline{}(𝐱)`$, for we shall first be discussing the solution $`y_1`$ of the Fredholm equation (4C.12) with kernel $`\kappa _{21}`$ and inhomogeneous term $`u_1`$ rather than the solution $`𝒴_1`$ of the Fredholm equation (3A.6) with kernel $`K_{21}`$ and inhomogeneous term $`U_1`$. When $`𝐯K^{2+}`$, the solution $`y_1`$ need not be differentiable. However, when $`𝐯K^3`$, the kernel $`K_{21}(𝐱,\sigma ^{},\sigma )`$ and the inhomogeneous term $`U_1(𝐱,\sigma )`$ in the Fredholm equation (3A.6) are $`𝐂^1`$ and $`𝐂^2`$ functions of $`(𝐱,\sigma ^{},\sigma )`$ and $`(𝐱,\sigma )`$, respectively; and the result is a differentiable $`y_1`$ as we shall see in Thm. 28. The following lemma is required for the proof of Thm. 28. ###### LEMMA 26 (Differentiability properties of $`u_1`$ and $`\kappa _{21}`$ when $`𝐯K^3`$) When $`𝐯K^3`$, $`u_1`$ is $`𝐂^2`$ and $`\kappa _{21}`$ is $`𝐂^1`$. Moreover, $`^2\kappa _{21}(𝐱,\theta ^{},\theta )/rs`$ exists and is a continuous function of $`(𝐱,\theta ^{},\theta )`$ throughout $`D\times \mathrm{\Theta }^2`$ \[whereupon, from a theorem of the calculus, $`^2\kappa _{21}/sr`$ also exists and is equal to $`^2\kappa _{21}/rs`$\]. Proof: The proof will be given in three parts: * From Eqs. (4C.8) and (4C.9), $`\sigma (𝐱,\theta )`$ is a real analytic function of $`\theta `$ throughout $`D\times \mathrm{\Theta }`$, $`\sigma (𝐱,\theta )`$ $``$ $`\|r,r_0\|\text{ when }\theta [0,\pi /2],`$ (5A.1) $`\sigma (𝐱,\theta )`$ $``$ $`\|s,s_0\|\text{ when }\theta [\pi ,3\pi /2],`$ (5A.2) Therefore, $$W(\sigma (𝐱,\theta ))\text{ is a }𝐂^3\text{ function of }(𝐱,\theta )\text{ throughout }D\times \mathrm{\Theta }$$ (5A.3) and $$\begin{array}{c}\hfill \lambda _{21}(\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta ^{\prime \prime }),\sigma (𝐱,\theta ))\text{ is a }𝐂^1\text{ function of }\\ \hfill (𝐱,\theta ^{},\theta ^{\prime \prime },\theta )\text{ throughout }D\times \mathrm{\Theta }^3.\end{array}$$ (5A.4) * To prove that $`{\displaystyle \frac{^2\lambda _{21}(\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta ^{\prime \prime }),\sigma (𝐱,\theta ))}{rs}}`$ exists and is a continuous function of $`(𝐱,\theta ^{},\theta ^{\prime \prime },\theta )\text{ throughout }D\times \mathrm{\Theta }^3,`$ (5A.5) we consider three distinct cases, (a), (b) and (c): + $$(\theta ^{\prime \prime },\theta )[0,\pi /2]\times [\pi ,3\pi /2]\text{ or }(\theta ^{\prime \prime },\theta )[\pi ,3\pi /2]\times [0,\pi /2].$$ (5A.6) + $$\begin{array}{c}(\theta ^{\prime \prime },\theta )[0,\pi /2]^2\text{ and }\theta ^{}[\pi ,3\pi /2],\text{ or }\hfill \\ (\theta ^{\prime \prime },\theta )[\pi ,3\pi /2]^2\text{ and }\theta ^{}[0,\pi /2].\hfill \end{array}$$ (5A.7) + $$(\theta ^{},\theta ^{\prime \prime },\theta )[0,\pi /2]^3\text{ or }(\theta ^{},\theta ^{\prime \prime },\theta )[\pi ,3\pi /2]^3.$$ (5A.8) In case (a) and case (b), it is easily seen that the denominator $`\sigma (𝐱,\theta ^{\prime \prime })\sigma (𝐱,\theta )`$ is different from zero, and hence $$\lambda _{21}(\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta ^{\prime \prime }),\sigma (𝐱,\theta ))\text{ is a }𝐂^2\text{ function of }(𝐱,\theta ^{},\theta ^{\prime \prime },\theta ),$$ (5A.9) from which the desired conclusion follows. In case (3) we employ $`{\displaystyle \frac{\sigma (𝐱,\theta )}{s}}`$ $`=`$ $`0\text{ when }\theta [0,\pi /2],`$ (5A.10) and $`{\displaystyle \frac{\sigma (𝐱,\theta )}{r}}`$ $`=`$ $`0\text{ when }\theta [\pi ,3\pi /2],`$ (5A.11) to show that the mixed second derivative of $`\lambda _{21}`$ exists and equals zero. * From Eqs. (4C.14) to (4C.18), $`u_1(𝐱,\theta )`$ $`=`$ $`W_1(\sigma (𝐱,\theta )){\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}p(𝐱,\theta ^{})L_1(\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta ))`$ (5A.12) $`\text{for all }(𝐱,\theta )D\times \mathrm{\Theta }`$ and $`\kappa _{21}(𝐱,\theta ^{},\theta )`$ $`=`$ $`q(𝐱,\theta ^{})\left[k_{21}(\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta )){\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{\prime \prime }p(𝐱,\theta ^{\prime \prime })\lambda _{21}(\sigma (𝐱,\theta ^{}),\sigma (𝐱,\theta ))\right]`$ (5A.13) $`\text{for all }(𝐱,\theta ^{},\theta )D\times \mathrm{\Theta }^2,`$ where $`p(𝐱,\theta )`$ is defined by Eq. (4C.19), and $`q(𝐱,\theta )`$ is defined by Eq. (4C.16). From statements (4C.10) and (4C.11), $$\begin{array}{c}p(𝐱,\theta )\text{ and }q(𝐱,\theta )\text{ are real analytic }\hfill \\ \text{functions of }(𝐱,\theta )\text{ throughout }D\times \mathrm{\Theta }.\hfill \end{array}$$ From statements (5A.3), (5A.4), (5A.5) and (3), it is clear that the functions $`u_1`$ and $`\kappa _{21}`$ whose values are given by Eqs. (5A.12) and (5A.13), respectively, satisfy the conclusions of our lemma. End of proof. Dfn. of a function that is $`𝐂^{N_1,\mathrm{},N_L}`$ on $`XR^L`$ * Suppose that $`X`$ is an open subset of $`R^L`$ or a closed or semi-closed subinterval of $`R^L`$, $`x=(x^1,\mathrm{},x^L)`$ denotes any point in $`X`$, $`T`$ is a topological space, $`t`$ denotes any point in $`T`$, and $`N_1,\mathrm{},N_L`$ are $`L`$ non-negative integers. Suppose, furthermore, that $`F:(X\times T)C`$ and that, for each $`L`$-tuple of integers $`(n_1,\mathrm{},n_L)`$ such that $`0n_kN_k`$ for all $`k=1,\mathrm{},L`$, $$_{1\mathrm{}L}^{n_1\mathrm{}n_L}F(x,t):=\left(\frac{}{x^1}\right)^{n_1}\mathrm{}\left(\frac{}{x^k}\right)^{n_k}\mathrm{}\left(\frac{}{x^L}\right)^{n_L}F(x,t)$$ (5A.14) exists and is a continuous function of $`(x,t)`$ throughout $`X\times T`$. \[It is understood that $`(/x^k)^0=1`$.\] Then, we shall say that $`F`$ is $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$. Also, if $`F:XC`$ and $`_{1\mathrm{}L}^{n_1\mathrm{}n_L}F(x)`$ exists and is a continuous function of $`x`$ throughout $`X`$ for each choice of $`(n_1,\mathrm{},n_L)`$ that satisfies $`0n_kN_k`$ for all $`1kL`$, then we shall say that $`F`$ is $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$. End of Dfn. Note: If $`F`$ is $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$, then a theorem of the calculus tells us that, for each $`(n_1,\mathrm{},n_L)`$ satisfying $`0n_kN_k`$ for all $`1kL`$, the existence and value of $`_{1\mathrm{}L}^{n_1\mathrm{}n_L}F`$ are unchanged when the operator factors $`/x^k`$ are subject to any permutation. The following lemma is applicable to a broad class of Fredholm integral equations and is clearly capable of further generalization in several directions. A $`2\times 2`$ matrix version of the lemma for the case $`L=2`$ was covered in a paper by the authors on the initial value problem for colliding gravitational plane wave pairs.<sup>26</sup><sup>26</sup>26I. Hauser and F. J. Ernst, J. Math. Phys. 32, 198 (1991), Sec. V. As regards the current notes, the lemma will play a key role in the proof of Thm. 28. ###### LEMMA 27 (Fredholm minor $`M`$ and determinant $`\mathrm{\Delta }`$) Let $`X`$, $`x`$ and $`N_k(k=1,\mathrm{},L)`$ be assigned the same meanings as in the preceding definition; and let $`Y`$ denote a compact, oriented, $`m`$-dimensional differentiable manifold, $`y`$ denote any point in $`Y`$, and $`dy`$ denote a volume element at point $`y`$ (the value of a distinguished non-zero $`m`$-form at $`y`$). Suppose that $`K:X\times (Y\times Y)C`$ and $`K`$ is $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$. Let us regard $`K`$ as an $`L`$-parameter family of Fredholm kernels that is employed in Fredholm integral equations of the form $$f(x,y)_Y𝑑y^{}f(x,y^{})K(x,y^{},y)=g(x,y)\text{ for all }(x,y)X\times Y,$$ (5A.15) where $`X`$ is the parameter space. Then, the corresponding Fredholm minor $`M`$ and Fredholm determinant $`\mathrm{\Delta }`$ are $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$. Proof: The Fredholm construction of $`M`$ and $`\mathrm{\Delta }`$ are given by $`M(x,y^{},y)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n!}}M^{(n)}(x,y^{},y),`$ (5A.16) $`\mathrm{\Delta }(x)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n!}}\mathrm{\Delta }^{(n)}(x),`$ (5A.17) where $`M^{(0)}(x,y^{},y)`$ $`:=`$ $`K(x,y^{},y),`$ (5A.18) $`M^{(n)}(x,y^{},y)`$ $`:=`$ $`{\displaystyle _Y}𝑑y_1\mathrm{}{\displaystyle _Y}𝑑y_nD^{(n+1)}\left(x\|\begin{array}{cc}y& y_1\mathrm{}y_n\\ y^{}& y_1\mathrm{}y_n\end{array}\right)`$ (5A.22) $`\text{for all }n>0,`$ $`\mathrm{\Delta }^{(0)}(x)`$ $`:=`$ $`1,`$ (5A.23) $`\mathrm{\Delta }^{(n+1)}(x)`$ $`:=`$ $`{\displaystyle _Y}𝑑yM^{(n)}(x,y,y)\text{ for all }n0,`$ (5A.24) and $`D^{(n)}\left(x\|\begin{array}{c}y_1\mathrm{}y_n\\ y_1^{}\mathrm{}y_n^{}\end{array}\right)`$ $`:=`$ $`\text{the determinant of that }n\times n`$ (5A.28) matrix whose element in the $`k`$th $`\text{row and }l\text{th column is }K(x,y_k^{},y_l).`$ In particular, $$D^{(0)}\left(x\|\begin{array}{c}y\\ y^{}\end{array}\right):=K(x,y^{},y).$$ (5A.29) For each bounded and closed subspace $`U`$ of $`X`$, let $`K_u`$ $`:=`$ $`sup\left\{\right\|_{1\mathrm{}L}^{n_1\mathrm{}n_L}K(x,y^{},y)\|:(x,y^{},y)U\times Y^2,`$ (5A.30) $`\text{and }0n_kN_k\text{ for all }k=1,\mathrm{},L\}.`$ Also let $$V:=_Y𝑑y.$$ (5A.31) Then, from Eqs. (5A.18) and (5A.22), and from a generalization of Hadamard’s inequality that was formulated and proved by the authors in the aforementioned paper on the initial value problem for colliding gravitational plane wave pairs \[see Thm. 7 in that paper\], $`\left\|_{1\mathrm{}L}^{n_1\mathrm{}n_L}M^n(x,y^{},y)\right\|`$ $``$ $`V^nK_U^{n+1}(n+1)^{N_1+\mathrm{}+N_L+(n+1)/2}`$ (5A.32) $`\text{for all }(x,y^{},y)U\times Y^2\text{ and all}`$ $`(n_1,\mathrm{},n_L)\text{ such that }0n_kN_k`$ $`\text{for each }k=1,\mathrm{},L.`$ It follows that, for each positive integer $`N`$, $`{\displaystyle \underset{n=0}{\overset{N}{}}}{\displaystyle \frac{1}{n!}}\left\|_{1\mathrm{}L}^{n_1\mathrm{}n_L}M^{(n)}(x,y^{},y)\right\|`$ $``$ $`{\displaystyle \underset{n=0}{\overset{N}{}}}{\displaystyle \frac{V^nK_U^{n+1}}{n!}}(n+1)^{N_1+\mathrm{}+N_L+(n+1)/2}`$ (5A.33) $`\text{for all }(x,y^{},y)U\times Y^2\text{ and all}`$ $`\text{choices (the usual) of }(n_1,\mathrm{},n_l).`$ The application of the ratio test to the series on the right side of the above inequality (5A.33) is straightforward and deomonstrates that this series converges as $`N\mathrm{}`$. Hence, from the comparison test, the series on the left side of (5A.33) converges for all $`(x,y^{},y)U\times Y^2`$ and all choices of $`(n_1,\mathrm{},n_L)`$. The theorems<sup>27</sup><sup>27</sup>27See Sec. 2, Ch. IV, of Differential and Integral Calculus by R. Courant (Interscience Publishers, Inc., 1936). of the calculus on the continuity and term-by-term differentiability of a uniformly convergent infinite series of functions then supply us with the following conclusions: $$\begin{array}{c}\text{For all choices of }(n_1,\mathrm{},n_L)\text{ for which }0n_kN_k(1kL),\hfill \\ _{1\mathrm{}L}^{n_1\mathrm{}n_L}M(x,y^{},y)\text{ exists and is a continuous function of }(x,y^{},y)\hfill \\ \text{throughout }X\times Y^2;\hfill \end{array}$$ and $`_{1\mathrm{}L}^{n_1\mathrm{}n_L}M(x,y^{},y)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n!}}\left[_{1\mathrm{}L}^{n_1\mathrm{}n_L}M^{(n)}(x,y^{},y)\right],\text{ and}`$ (5A.34) the infinite series converges absolutely and converges uniformly on each $`\text{compact subspace of }X\times Y^2.`$ Hence, $`M`$ is $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$. The proof that $`\mathrm{\Delta }`$ is also $`𝐂^{N_1,\mathrm{},N_L}`$ on $`X`$ is left for the reader. End of proof. The following theorem concerns the solution $`y_1(𝐱,\theta )`$ of the Fredholm equation (4C.12) for all $`(𝐱,\theta )D\times \mathrm{\Theta }`$. ###### THEOREM 28 (Differentiability properties of $`y_1`$ when $`𝐯K^3`$) If $`𝐯K^3`$, then $`y_1`$ is $`𝐂^{1,1}`$ on $`D`$; i.e., $`y_1(𝐱,\theta )/r`$, $`y_1(𝐱,\theta )/s`$ and $`^2y_1(𝐱,\theta )/rs`$ exist and are continuous functions of $`(𝐱,\theta )`$ throughout $`D\times \mathrm{\Theta }`$. Proof: Consider the inhomogenous Fredholm equation of the second kind (4C.12). According to Thm. 21, the Fredholm determinant for Eq. (4C.12) is not zero for all choices of $`𝐱D`$. Therefore, a unique solution of the Fredholm equation exists and is given by $$y_1(𝐱,\theta )=u_1(𝐱,\theta )+\frac{2}{\pi }_\mathrm{\Theta }𝑑\theta ^{}u_1(𝐱,\theta ^{})R(𝐱,\theta ^{},\theta )$$ (5A.35) for all $`(𝐱,\theta )D\times \mathrm{\Theta }`$, where the resolvent kernel $`R(𝐱,\theta ^{},\theta )`$ is the following ratio of the Fredholm minor and determinant: $$R(𝐱,\theta ^{},\theta )=\frac{M(𝐱,\theta ^{},\theta )}{\mathrm{\Delta }(𝐱)}.$$ (5A.36) From Lem. 26, $`\kappa _{21}`$ is $`𝐂^1`$. Moreover, $`^2\kappa _{21}(𝐱,\theta ^{},\theta )/rs`$ exists and is a continuous function of $`(𝐱,\theta ^{},\theta )`$ throughout $`D\times \mathrm{\Theta }^2`$. Therefore, $$\kappa _{21}\text{ is }𝐂^{1,1}\text{ on }D.$$ (5A.37) The preceding Lem. 27 is now applied to the present case, for which $$X=D,Y=\mathrm{\Theta },L=2,m=1,dy=2d\theta /\pi .$$ (5A.38) Thereupon, one obtains $$R\text{ is }𝐂^{1,1}\text{ on }D.$$ (5A.39) Lemma 26 also tells us that (amongst other things) $$u_1\text{ is }𝐂^{1,1}\text{ on }D.$$ (5A.40) Therefore, from Eq. (5A.35), statements (5A.39) and (5A.40), and the theorems<sup>28</sup><sup>28</sup>28See Ref. 27. of the calculus on the continuity and differentiability of an integral with respect to parameters, $$y_1\text{ is }𝐂^{1,1}\text{ on }D.$$ (5A.41) End of proof. Note that, in terms of standard notation and terminology, $`\lambda =1`$ for our particular Fredholm equation; and the statement that $`\mathrm{\Delta }(𝐱)0`$ is equivalent to the statement that $`1`$ is not a characteristic value (eigenvalue) of our kernel. ### B. Concerning the partial derivatives of $`𝒴`$, $``$, $`H`$ and $`^\pm `$ when $`𝐯K^3`$ Dfn. of $`L^{(i)}(\sigma ^{},\sigma )`$ for each $`𝐱D`$ and $`i\{3,4\}`$ * For each $`\sigma ^{}\overline{}(𝐱)`$ and $`\sigma ^{(i)}`$, let $$L^{(i)}(\sigma ^{},\sigma ):=\frac{W(\sigma ^{})W^{(i)}(\sigma )}{\sigma ^{}\sigma }.$$ End of Dfn. Employing the transformation defined by Eqs. (4C.1) to (4C.9), the definition of $`p(𝐱,\theta )`$ by Eq. (4C.19), the definition of $`q(𝐱,\theta )`$ by Eq. (4C.16) and the definition of $`L^{(i)}(\sigma ^{},\sigma )`$ that we just gave, one finds that Eqs. (2C.2), (2C.1) and (2B.46) are expressible in the forms (in which ‘$`𝐱`$’ is no longer suppressed) $`𝒴_1^{(i)}(𝐱,\sigma )`$ $`=`$ $`W_1^{(i)}(\sigma ){\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}p(𝐱,\theta ^{})y_2(𝐱,\theta ^{})W_1^T(\sigma (𝐱,\theta ^{}))JL_1^{(i)}(\sigma (𝐱,\theta ^{}),\sigma )`$ (5B.1) $`\text{for all }𝐱D\text{ and }\sigma \stackrel{ˇ}{}^{(i)}(x^{7i})`$ \[after the extension defined by (2C.7)\], $`𝒴_2^{(i)}(𝐱,\sigma )`$ $`=`$ $`W_2^{(i)}(\sigma )+{\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}q(𝐱,\theta ^{})y_1(𝐱,\theta ^{})W_2^T(\sigma (𝐱,\theta ^{}))JL_2^{(i)}(\sigma (𝐱,\theta ^{}),\sigma )`$ (5B.2) $`\text{for all }𝐱D\text{ and }\sigma \stackrel{ˇ}{}^{(i)}(x^{7i})`$ \[after the extension defined by (2C.7)\], and $`\nu (𝐱,\tau )^1(𝐱,\tau )`$ $`=`$ $`I{\displaystyle \frac{2}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}q(𝐱,\theta ^{})y_1(𝐱,\theta ^{}){\displaystyle \frac{W_2(\sigma (𝐱,\theta ^{}))J}{\sigma (𝐱,\theta ^{})\tau }}`$ (5B.3) $`\text{for all }𝐱D\text{ and }\tau C\overline{}(𝐱).`$ Furthermore, from Eqs. (4D.1), (4D.6) and (5B.3), $`H(𝐱)`$ $`=`$ $`H^M(𝐱_0)+2(zz_0)\mathrm{\Omega }{\displaystyle \frac{4i}{\pi }}{\displaystyle _\mathrm{\Theta }}𝑑\theta ^{}q(𝐱,\theta ^{})y_1(𝐱,\theta ^{})W_2^T(\sigma (𝐱,\theta ^{}))`$ (5B.4) $`\text{for all }𝐱D.`$ When proving the following theorem, one should bear in mind that $`\sigma (𝐱,\theta )`$, $`p(𝐱,\theta )`$ and $`q(𝐱,\theta )`$ are analytic functions of $`(𝐱,\theta )`$ throughout $`D\times \mathrm{\Theta }`$. ###### THEOREM 29 (Differentiability properties of $`𝒴^{(i)}`$, $``$ and $`H`$ when $`𝐯K^3`$) If $`𝐯K^3`$, then $$\begin{array}{c}𝒴^{(i)}(𝐱,\sigma )/r,𝒴^{(i)}(𝐱,\sigma )/s,^2𝒴^{(i)}(𝐱,\sigma )/rs,\hfill \\ ^2𝒴^{(i)}(𝐱,\sigma )/\sigma ^2,^2𝒴^{(i)}(𝐱,\sigma )/r\sigma ,\text{ and }^2𝒴^{(i)}(𝐱,\sigma )/s\sigma \hfill \\ \text{exist and are continuous functions of }(𝐱,\sigma )\text{ throughout }\hfill \\ \{(𝐱,\sigma ):𝐱D,\sigma \stackrel{ˇ}{}^{(i)}(x^{7i})\}.\hfill \end{array}$$ (5B.5) Also, upon letting $`\stackrel{`}{}`$ denote the restriction of $``$ to $$\mathrm{d}om\stackrel{`}{}:=\{(𝐱,\tau ):𝐱D,\tau C(𝐱)\{r,s,r_0,s_0\}\},$$ one has $$\begin{array}{c}\stackrel{`}{}(𝐱,\tau )/r,\stackrel{`}{}(𝐱,\tau )/s,\text{ and }^2\stackrel{`}{}(𝐱,\tau )/rs\hfill \\ \text{exist and are continuous functions of }(𝐱,\tau )\hfill \\ \text{throughout }\mathrm{d}om\stackrel{`}{};\text{ and, for each }𝐱D,\text{ these }\hfill \\ \text{partial derivatives are holomorphic functions }\hfill \\ \text{of }\tau \text{ throughout }C(𝐱)\{r,s,r_0,s_0\}.\hfill \end{array}$$ (5B.6) Furthermore, $$H\text{ is }𝐂^{1,1}\text{ on }D.$$ (5B.7) Proof: From Thm. 28, statement (5A.3) and the fact that $`L^{(i)}`$ is $`𝐂^2`$, one concludes from Eq. (5B.2) that $$\begin{array}{c}𝒴_2^{(i)}(𝐱,\sigma )/r,𝒴_2^{(i)}(𝐱,\sigma )/s,\hfill \\ ^2𝒴_2^{(i)}(𝐱,\sigma )/rs,^2𝒴_2^{(i)}(𝐱,\sigma )/\sigma ^2,\hfill \\ ^2𝒴_2^{(i)}(𝐱,\sigma )/r\sigma \text{ and }^2𝒴_2^{(i)}(𝐱,\sigma )/s\sigma \hfill \end{array}$$ exist and are continuous functions of $`(𝐱,\sigma )`$ throughout $`\{(𝐱,\sigma ):𝐱D,\sigma \stackrel{ˇ}{}^{(i)}(x^{7i})\}`$. Then, from Eq. (5B.1), one obtains like conclusions for $`𝒴_1^{(i)}(𝐱,\sigma )`$, whereupon the statement (5B.5) follows. Statements (5B.6) and (5B.7) follow from Thm. 28, statement (5A.3), the known differentiability and holomorphy properties of $`\nu (𝐱,\tau )^1`$ on $`\mathrm{d}om\stackrel{`}{}`$, and the theorem on the holomorphy of functions given by Cauchy-type integrals. End of proof. The following two lemmas will be used to prove Thm. 32. ###### LEMMA 30 ($`d(\nu (𝐱,\tau )^1\stackrel{`}{}(𝐱,\tau ))`$) If $`𝐯K^3`$, then the first partial derivatives of $$\frac{\nu ^+(𝐱,\sigma ^{})^1𝒴_1(𝐱,\sigma ^{})W_2^T(\sigma ^{})J}{\sigma ^{}\tau }$$ (5B.8) with respect to $`r`$ and with respect to $`s`$ are summable over $`\sigma ^{}\overline{}(𝐱)`$; and $`d\left[\nu (𝐱,\tau )^1\stackrel{`}{}(𝐱,\tau )\right]`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{d\left[\nu ^+(𝐱,\sigma ^{})^1𝒴_1(𝐱,\sigma ^{})\right]W_2^T(\sigma ^{})J}{\sigma ^{}\tau }}`$ (5B.9) $`\text{for all }(𝐱,\tau )\mathrm{d}om\stackrel{`}{}.`$ Proof: We shall tacitly employ statements (5B.5) and (5B.6) of Thm. 29 in some steps of this proof. We shall supply the proof only for $`[\nu (𝐱,\tau )^1\stackrel{`}{}(𝐱,\tau )]/r`$ and leave the proof for the partial derivative with respect to $`s`$ for the reader. The summability over $`\overline{}(𝐱)`$ of the partial derivative with respect to $`r`$ of (5B.8) is seen from the facts that $`\nu ^+(𝐱,\sigma ^{})^1`$ $`=`$ $`^+(\sigma ^{}r)^+(\sigma ^{}s)\left[^+(\sigma ^{}r_0)^+(\sigma ^{}s_0)\right]^1`$ (5B.10) and $`{\displaystyle \frac{\nu ^+(𝐱,\sigma ^{})^1}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}^+(\sigma ^{}s)\left[^+(\sigma ^{}r)^+(\sigma ^{}r_0)^+(\sigma ^{}s_0)\right]^1,`$ (5B.11) where $$^+(\sigma ):=\{\begin{array}{ccc}\hfill \sqrt{\sigma }& \text{if}& \sigma 0,\hfill \\ \hfill i\sqrt{\sigma }& \text{if}& \sigma 0,\hfill \end{array}$$ are both summable over $`\overline{}(𝐱)`$, and a summable function times a continuous function over a bounded interval is summable. In the proofs of this lemma and the next lemma, we shall employ the shorthand notations $$\begin{array}{ccc}\hfill f(𝐱,\sigma ^{})& :=& 𝒴_1(𝐱,\sigma ^{})W_2^T(\sigma ^{})J,\hfill \\ \hfill g(𝐱,\sigma ^{})& :=& \nu (𝐱,\tau )^1(𝐱,\tau ),\hfill \end{array}$$ (5B.12) whereupon Eq. (2B.46) becomes $`g(𝐱,\tau )`$ $`=`$ $`I{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1{\displaystyle \frac{f(𝐱,\sigma ^{})}{\sigma ^{}\tau }}`$ (5B.13) $`=`$ $`Ig_1(𝐱,\tau )f(𝐱,r)g_2(𝐱,\tau ),`$ where $`g_1(𝐱,\tau )`$ $`:=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1{\displaystyle \frac{f(𝐱,\sigma ^{})f(𝐱,r)}{\sigma ^{}\tau }},`$ (5B.14) and $`g_2(𝐱,\tau )`$ $`:=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{\nu ^+(𝐱,\sigma ^{})^1}{\sigma ^{}\tau }}.`$ (5B.15) We shall first deal with the term $`f(𝐱,r)g_2(𝐱,\tau )`$. It is easy to show that $$g_2(𝐱,\tau )=\nu (𝐱,\tau )^11.$$ (5B.16) Therefore, for all $`(𝐱,\tau )\mathrm{d}om\stackrel{`}{}`$, $$\frac{g_2(𝐱,\tau )}{r}=\frac{1}{2(\tau r)}\nu (𝐱,\tau )^1.$$ (5B.17) Also, note that $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{\nu ^+(𝐱,\sigma ^{})^1/r}{\sigma ^{}\tau }}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{\nu ^+(𝐱,\sigma ^{})^1}{(\sigma ^{}r)(\sigma ^{}\tau )}}`$ (5B.18) $`=`$ $`{\displaystyle \frac{\nu (𝐱,\tau )^1}{2(\tau r)}}.`$ So, from Eqs. (5B.15), (5B.17) and (5B.18), $`{\displaystyle \frac{}{r}}\left[f(𝐱,r)g_2(𝐱,\tau )\right]`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{\left[\nu ^+(𝐱,\sigma ^{})^1f(𝐱,r)\right]/r}{\sigma ^{}\tau }}`$ (5B.19) $`\text{for all }(𝐱,\tau )\mathrm{d}om\stackrel{`}{}.`$ That takes care of the term $`f(𝐱,r)g_2(𝐱,\tau )`$. We shall next deal with the term $`g_1(𝐱,\tau )`$. From statement (5B.5) in Thm. 29 and from Eq. (5B.12), one can see that $$\frac{}{r}\left\{^+(\sigma ^{}r)^+(\sigma ^{}s)\left[f(𝐱,\sigma ^{})f(𝐱,r)\right]\right\}$$ (5B.20) exists and is a continuous function of $`(𝐱,\sigma ^{})`$ throughout $`\{(𝐱,\sigma ^{}):𝐱D,\sigma ^{}\overline{}(𝐱)\}`$. \[We leave details for the reader.\] No loss of generality will be incurred if we tentatively introduce a closed and bounded convex neighborhood $`𝒩`$ of the point $`𝐱_0`$ in the space $`D`$, whereupon it is seen that $$\{(𝐱,\sigma ^{}):𝐱𝒩,\sigma ^{}\overline{}(𝐱)\}$$ is a bounded closed subspace of $`R^3`$; and, therefore, $`M(𝒩):=`$ $`sup\left\{\right\|\left\|{\displaystyle \frac{}{r}}\left\{^+(\sigma ^{}r)^+(\sigma ^{}s)[f(𝐱,\sigma ^{})f(𝐱,r)]\right\}\right\|\|:𝐱𝒩,\sigma ^{}\overline{}(𝐱)\}`$ is finite; and the integrand in the expression for $`g_1(𝐱,\tau )`$ that is given by Eq. (5B.14) satisfies $$\frac{}{r}\left[\nu ^+(𝐱,\sigma ^{})^1\frac{f(𝐱,\sigma ^{})f(𝐱,r)}{\sigma ^{}\tau }\right]\left[\sqrt{\|\sigma ^{}r_0\|\|\sigma ^{}s_0\|}\|\sigma ^{}\tau \|\right]^1M(𝒩).$$ (5B.22) Since the right side of the above inequality is summable over $`\overline{}(𝐱)`$ and is independent of $`𝐱`$, a well-known theorem<sup>29</sup><sup>29</sup>29See Ref. 18, Sec. 39. on differentiation of a Lebesgue integral with respect to a parameter tells us that $`g_1(𝐱,\tau )/r`$ exists (which, it happens, we already know) and is given by $$\frac{g_1(𝐱,\tau )}{r}=\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\frac{\frac{}{r}\left\{\nu ^+(𝐱,\sigma ^{})^1\left[f(𝐱,\sigma ^{})f(𝐱,r)\right]\right\}}{\sigma ^{}\tau }$$ (5B.23) for all $`𝐱𝒩`$ and $`\tau C(𝐱)\{r,s,r_0,s_0\}`$, where we have used the fact that the contribution to $`g_1(𝐱,\tau )/r`$ due to differentiation of the integral with respect to the endpoint $`r\{a^3,b^3\}`$ of the integration interval $`\overline{}^{(3)}(𝐱)`$ vanishes, because the integrand in Eq. (5B.14) vanishes when $`\sigma ^{}=r`$. However, since $`𝒩`$ can always be chosen so that it covers any given point in $`D`$, Eq. (5B.23) holds for all $`(𝐱,\tau )\mathrm{d}om\stackrel{`}{}`$; and upon combining (5B.23), (5B.19) and (5B.13), one obtains $`{\displaystyle \frac{g(𝐱,\tau )}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{[\nu ^+(𝐱,\sigma ^{})^1f(𝐱,\sigma ^{})]/r}{\sigma ^{}\tau }}`$ (5B.24) $`\text{for all }(𝐱,\tau )\mathrm{d}om\stackrel{`}{},`$ which is the coefficient of $`dr`$ in Eq. (5B.9). End of proof. Before we give the next lemma, note that application of the Plemelj relations to Eq. (2B.46) yields $`{\displaystyle \frac{1}{2}}[^+(𝐱,\sigma )+^{}(𝐱,\sigma )]`$ $`=`$ $`𝒴_1(𝐱,\sigma )W_2^T(\sigma )J`$ (5B.25) $`\text{for all }𝐱D\text{ and }\sigma (𝐱),`$ and $`{\displaystyle \frac{1}{2}}\nu ^+(𝐱,\sigma )^1[^+(𝐱,\sigma )^{}(𝐱,\sigma )]`$ $`=`$ $`I{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1{\displaystyle \frac{𝒴_1(𝐱,\sigma ^{})W_2^T(\sigma ^{})J}{\sigma ^{}\sigma }}`$ (5B.26) $`\text{for all }𝐱D\text{ and }\sigma (𝐱).`$ ###### LEMMA 31 (Differentiability properties of $`^\pm `$ when $`𝐯K^3`$) As in the preceding lemma, suppose that $`𝐯K^3`$ and $``$ is the solution of the HHP corresponding to $`(𝐯,^M)`$. Then the following three statements hold: * The partial derivatives $`^\pm (𝐱,\sigma )/r`$, $`^\pm (𝐱,\sigma )/s`$ and $`^2^\pm (𝐱,\sigma )/rs`$ exist and are continuous functions of $`(𝐱,\sigma )`$ throughout $`\{(𝐱,\sigma ):𝐱D,\sigma (𝐱)\}`$. * The $`1`$-form $$\frac{d[\nu ^+(𝐱,\sigma ^{})^1𝒴_1(𝐱,\sigma ^{})]W_2^T(\sigma ^{})J}{\sigma ^{}\sigma }$$ (5B.27) is, for each $`𝐱D`$ and $`\sigma (𝐱)`$, summable over $`\overline{}(𝐱)`$ in the PV sense. * For all $`𝐱D`$ and $`\sigma (𝐱)`$, $`d\left\{{\displaystyle \frac{1}{2}}\nu ^+(𝐱,\sigma )^1[^+(𝐱,\sigma )^{}(𝐱,\sigma )]\right\}=`$ (5B.28) $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{d[\nu ^+(𝐱,\sigma ^{})^1𝒴_1(𝐱,\sigma ^{})]W_2^T(\sigma ^{})J}{\sigma ^{}\sigma }}.`$ Proofs: * This follows from statement (5B.5), Eq. (5B.25) and Eq. (2B.42). End of proof. The proofs of parts (ii) and (iii) will be supplied only for the coefficients of $`dr`$ in Eqs. (5B.27) and (5B.28). The proofs for the coefficients of $`ds`$ are left to the reader. * As functions of $`\sigma ^{}`$, $`W_2^T(\sigma ^{})`$ is $`𝐂^3`$, $`𝒴_1(𝐱,\sigma ^{})`$ is $`𝐂^2`$ and $`𝒴_1(𝐱,\sigma ^{})/r`$ is $`𝐂^1`$ on $`\overline{}(𝐱)`$; and $`\nu ^+(𝐱,\sigma ^{})^1`$ and $`\nu ^+(𝐱,\sigma ^{})^1/r`$ are summable over $`\overline{}(𝐱)`$. Therefore, for a sufficiently small $`ϵ>0`$, $$\frac{\frac{}{r}\left[\nu ^+(𝐱,\sigma ^{})^1𝒴_1(𝐱,\sigma ^{})\right]W_2^T(\sigma ^{})J}{\sigma ^{\prime \prime }\sigma }$$ (5B.29) is summable over $`\overline{}(𝐱)]\sigma ϵ,\sigma +ϵ[`$. Moreover, since the numerator of (5B.29) is a $`𝐂^1`$ function of $`\sigma ^{}`$, it is well known that (5B.29) is summable over $`[\sigma ϵ,\sigma +ϵ]`$ in the PV sense. Therefore, (5B.29) is summable over $`\overline{}(𝐱)`$ in the PV sense. End of proof. * In terms of the shorthand notations (5B.12), Eq. (5B.26) is expressible in the form $$\frac{1}{2}[g^+(𝐱,\sigma )+g^{}(𝐱,\sigma )]=I\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1\frac{f(𝐱,\sigma ^{})}{\sigma ^{}\sigma },$$ (5B.30) where Thm. 29 furnishes the following properties of $`f(𝐱,\sigma ^{})`$: $$\begin{array}{c}f(𝐱,\sigma ^{})/r,f(𝐱,\sigma ^{})/s,^2f(𝐱,\sigma ^{})/(\sigma ^{})^2,\hfill \\ ^2f(𝐱,\sigma ^{})/rs,^2f(𝐱,\sigma ^{})/r\sigma ^{}\text{ and }^2f(𝐱,\sigma ^{})/s\sigma ^{}\hfill \\ \text{exist and are continuous functions of }(𝐱,\sigma ^{})\hfill \\ \text{throughout }\{(𝐱,\sigma ^{}):𝐱D,\sigma ^{}\overline{}(𝐱)\}.\hfill \end{array}$$ (5B.31) Let us introduce the additional shorthand notations $`f_0(𝐱,\sigma ^{},\sigma )`$ $`:=`$ $`{\displaystyle \frac{f(𝐱,\sigma ^{})f(𝐱,\sigma )}{\sigma ^{}\sigma }},`$ (5B.32) $`f_1(𝐱,\sigma ^{},\sigma )`$ $`:=`$ $`f_0(𝐱,\sigma ^{},\sigma )f_0(𝐱,r,\sigma ),`$ (5B.33) $`g_1(𝐱,\sigma )`$ $`:=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1f_1(𝐱,\sigma ^{},\sigma ),`$ (5B.34) $`g_2(𝐱,\sigma )`$ $`:=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{\nu ^+(𝐱,\sigma ^{})^1}{\sigma ^{}\sigma }}`$ (5B.35) and $`g_3(𝐱,\sigma )`$ $`:=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1.`$ (5B.36) Then Eq. (5B.30) is expressible in the form $$\frac{1}{2}[g^+(𝐱,\sigma )+g^{}(𝐱,\sigma )]=Ig_1(𝐱,\sigma )f(𝐱,\sigma )g_2(𝐱,\sigma )f_0(𝐱,r,\sigma )g_3(𝐱,\sigma ).$$ (5B.37) Let us first consider the above terms that contain $`g_2`$ and $`g_3`$. A well-known formula yields $$g_2(𝐱,\sigma )=1,$$ (5B.38) while the usual contour integration technique yields $$g_3(𝐱,\sigma )=\frac{1}{2}(r+sr_0s_0).$$ (5B.39) Therefore, by using $$\frac{1}{\pi i}_\overline{}𝑑\sigma ^{\prime \prime }\nu ^+(\sigma ^{\prime \prime })^1(\sigma ^{\prime \prime }\sigma )^1(\sigma ^{}\sigma ^{\prime \prime })^1=0\text{ for all }\sigma \overline{}(𝐱)\{r_0,s_0\}.$$ (5B.40) and the fact that $$\nu ^+(𝐱,\sigma ^{})^1/r=\frac{1}{2}(\sigma ^{}r)^1\nu ^+(𝐱,\sigma ^{})^1,$$ the reader can prove that $`{\displaystyle \frac{g_2(𝐱,\sigma )}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{\nu ^+(𝐱,\sigma ^{})^1/r}{\sigma ^{}\sigma }}`$ (5B.41) and $`{\displaystyle \frac{g_3(𝐱,\sigma )}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1/r,`$ (5B.42) whereupon $`{\displaystyle \frac{[f(𝐱,\sigma )g_2(𝐱,\sigma )]}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{[f(𝐱,\sigma )\nu ^+(𝐱,\sigma ^{})^1]/r}{\sigma ^{}\sigma }},`$ (5B.43) and $`{\displaystyle \frac{[f(𝐱,\sigma )g_3(𝐱,\sigma )]}{r}}`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{}{r}}\left[f_0(𝐱,r,\sigma )\nu ^+(𝐱,\sigma ^{})^1\right].`$ (5B.44) That completes the analysis of the terms in Eq. (5B.37) that contain $`g_2`$ and $`g_3`$. We next consider $`g_1`$. From (5B.31) to (5B.33), one sees that $$\begin{array}{c}\text{For each }\sigma (𝐱),f_1(𝐱,\sigma ^{},\sigma )/r\text{ and }(\sigma ^{}r)^1f_1(𝐱,\sigma ^{},\sigma )\hfill \\ \text{exist and are continuous functions of }(𝐱,\sigma ^{})\hfill \\ \text{throughout }\{(𝐱,\sigma ^{}):𝐱D,\sigma ^{}\overline{}(𝐱)\}.\hfill \end{array}$$ (5B.45) Therefore, as regards the integrand in the definition (5B.34) of $`g_1(𝐱,\sigma )`$, one readily deduces (by an argument similar to the one used in the proof of the preceding lemma) that, corresponding to each closed and bounded neighborhood $`𝒩`$ of the point $`𝐱_0`$ in the space $`D`$, and each $`\sigma (𝐱)`$, there exists a positive real number $`M(𝒩,\sigma )`$ such that $`{\displaystyle \frac{}{r}}\left[\nu ^+(𝐱,\sigma ^{})^1f_1(𝐱,\sigma ^{},\sigma )\right]`$ $``$ $`{\displaystyle \frac{M(𝒩,\sigma )}{\sqrt{\|\sigma ^{}r_0\|\|\sigma ^{}s_0\|}}}`$ (5B.46) $`\text{for all }𝐱𝒩\text{ and }\sigma ^{}\overline{}(𝐱)\{r_0,s_0\}.`$ The remainder of the proof employs the same theorem on differentiation of a Lebesgue integral with respect to a parameter that was used in the proof of the preceding lemma. The result is $$\frac{g_1(𝐱,\sigma )}{r}=\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\frac{}{r}\left[\nu ^+(𝐱,\sigma ^{})^1f_1(𝐱,\sigma ^{},\sigma )\right].$$ (5B.47) Upon combining the results given by Eqs. (5B.43), (5B.44) and (5B.47), one obtains with the aid of Eqs. (5B.30), (5B.32) to (5B.34), and Eq. (5B.37), $`{\displaystyle \frac{}{r}}{\displaystyle \frac{1}{2}}[g^+(𝐱,\sigma )+g^{}(𝐱,\sigma )]`$ $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{[\nu ^+(𝐱,\sigma ^{})^1f(𝐱,\sigma ^{})]/r}{\sigma ^{}\sigma }}`$ (5B.48) $`\text{for all }𝐱D\text{ and }\sigma (𝐱).`$ End of proof. The point of the preceding two lemmas is the following crucial theorem. ###### THEOREM 32 (Limits of $`d`$ when $`𝐯K^3`$) Suppose $`𝐯K^3`$ and $``$ is the solution of the HHP corresponding to $`(𝐯,^M)`$. Then, the following three statements hold: * For each $`𝐱D`$ and $`\sigma (𝐱)`$, $`d\stackrel{`}{}(𝐱,\sigma \pm \zeta )`$ converges as $`\zeta 0`$ ($`\mathrm{I}m\zeta >0`$) and $$\underset{\zeta 0}{lim}d\stackrel{`}{}(𝐱,\sigma \pm \zeta )=d^\pm (𝐱,\sigma ).$$ (5B.49) Note: The existences of $`d\stackrel{`}{}(𝐱,\tau )`$ and $`d^\pm (𝐱,\sigma )`$ are guaranteed by Thm. 29 \[statement (5B.6)\] and by Lem. 31(i), respectively. * $`(𝐱,\tau )`$ converges as $`\tau r_0`$ and as $`\tau s_0`$ \[$`\tau C\overline{}(𝐱)`$\]; and $`\nu (𝐱,\tau )^1(𝐱,\tau )`$ converges as $`\tau r`$ and as $`\tau s`$. * For each $`i\{3,4\}`$, $$(\tau x^i)\frac{\stackrel{`}{}(𝐱,\tau )}{x^i}$$ (5B.50) converges as $`\tau r_0`$ and as $`\tau s_0`$, while $$\nu (𝐱,\tau )^1(\tau x^i)\frac{\stackrel{`}{}(𝐱,\tau )}{x^i}$$ (5B.51) converges as $`\tau r`$ and as $`\tau s`$. Proofs: * We shall prove statement (i) for the coefficient of $`dr`$ in $`d\stackrel{`}{}(𝐱,\tau )`$ and leave the proof for the coefficient of $`ds`$ to the reader. Employ the shorthand notation $$f(𝐱,\sigma ^{}):=𝒴_1(𝐱,\sigma ^{})W_2^T(\sigma ^{})J$$ (5B.52) in the integrand of Eq. (2B.46), which then becomes $`\nu (𝐱,\tau )^1(𝐱,\tau )`$ $`=`$ $`I{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu (𝐱,\sigma ^{})^1{\displaystyle \frac{f(𝐱,\sigma ^{})}{\sigma ^{}\tau }}`$ (5B.53) $`\text{for all }(𝐱,\tau )\mathrm{d}om,`$ whereupon, from Eq. (5B.9) in Lem. 30, and from Eq. (5B.18), $`{\displaystyle \frac{\nu (𝐱,\tau )^1}{2(\tau r)}}(𝐱,\tau )+\nu (𝐱,\tau )^1{\displaystyle \frac{(𝐱,\tau )}{r}}=`$ (5B.54) $`\mathrm{\Phi }(𝐱,\tau )+{\displaystyle \frac{\nu (𝐱,\tau )^1}{2(\tau r)}}f(𝐱,r)`$ $`\text{for all }(𝐱,\tau )\mathrm{d}om\stackrel{`}{},`$ where $`\mathrm{\Phi }(𝐱,\tau )`$ $`:=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1{\displaystyle \frac{\varphi (𝐱,\sigma ^{})}{\sigma ^{}\tau }}`$ (5B.55) and $`\varphi (𝐱,\sigma ^{})`$ $`:=`$ $`{\displaystyle \frac{f(𝐱,\sigma ^{})}{r}}{\displaystyle \frac{f(𝐱,\sigma ^{})f(𝐱,r)}{2(\sigma ^{}r)}}.`$ (5B.56) From Eq. (5B.56) and the properties of $`f(𝐱,\sigma ^{})`$ given by statement (5B.31) $$\begin{array}{c}\varphi (𝐱,\sigma ^{})/\sigma ^{}\text{ exists and is a continuous function of }\hfill \\ (𝐱,\sigma ^{})\text{ throughout }\{(𝐱,\sigma ^{}):𝐱D,\sigma ^{}\overline{}(𝐱)\}.\hfill \end{array}$$ (5B.57) Therefore, $`\nu (𝐱,\sigma ^{})^1\varphi (𝐱,\sigma ^{})`$ obeys a Hölder condition of index $`1`$ on each closed subinterval of $`(𝐱)`$; and it follows from the theorem in Sec. 16 in Muskhelishvili’s treatise<sup>30</sup><sup>30</sup>30See footnote 13. that (5B.55) satisfies $$\mathrm{\Phi }^\pm (𝐱,\sigma ):=\underset{\zeta 0}{lim}\mathrm{\Phi }(𝐱,\sigma \pm \zeta )\text{ exists for all }\sigma (𝐱).$$ (5B.58) Moreover, from the Plemelj relations \[Eq. (17.2) in Sec. 17 of Muskhelishvili’s treatise\], $$\mathrm{\Phi }^\pm (𝐱,\sigma )=\pm \nu ^+(𝐱,\sigma )^1\varphi (𝐱,\sigma )+\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1\frac{\varphi (𝐱,\sigma ^{})}{\sigma ^{}\sigma }.$$ (5B.59) \[The existence of the above PV integral is demonstrated in Sec. 12 of Muskhelishvili’s treatise.\] From Eq. (5B.54), condition (3) in the definition of the HHP \[the one about the existence of $`^\pm (𝐱)`$\] and statement (5B.58), $$\underset{\zeta 0}{lim}\frac{\stackrel{`}{}(𝐱,\sigma \pm \zeta )}{r}\text{ exists for each }𝐱D\text{ and }\sigma (𝐱);$$ (5B.60) and, with the aid of Eqs. (5B.25), (5B.52) and (5B.59), $$\underset{\zeta 0}{lim}\frac{1}{2}\left[\frac{\stackrel{`}{}(𝐱,\sigma +\zeta )}{r}+\frac{\stackrel{`}{}(𝐱,\sigma \zeta )}{r}\right]=\frac{f(𝐱,\sigma )}{r}$$ (5B.61) and $`\underset{\zeta 0}{lim}{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{r}}\left[\nu (𝐱,\sigma +\zeta )^1\stackrel{`}{}(𝐱,\sigma +\zeta )+\nu (𝐱,\sigma \zeta )^1\stackrel{`}{}(𝐱,\sigma \zeta )\right]`$ (5B.62) $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1{\displaystyle \frac{\varphi (𝐱,\sigma ^{})}{\sigma ^{}\sigma }}.`$ However, from Eq. (5B.40), $$\frac{1}{\pi i}_\overline{}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1\frac{f(𝐱,r)}{(\sigma ^{}r)(\sigma ^{}\sigma )}=0.$$ Therefore, from Eq. (5B.56), Eq. (5B.62) becomes $`\underset{\zeta 0}{lim}{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{r}}\left[\nu (𝐱,\sigma +\zeta )^1\stackrel{`}{}(𝐱,\sigma +\zeta )+\nu (𝐱,\sigma \zeta )^1\stackrel{`}{}(𝐱,\sigma \zeta )\right]`$ (5B.63) $`=`$ $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{[\nu ^+(𝐱,\sigma ^{})^1f(𝐱,\sigma ^{})]/r}{\sigma ^{}\sigma }}.`$ Next, from Eqs. (5B.25) and (5B.52), $$\frac{1}{2}\left[\frac{^+(𝐱,\sigma )}{r}+\frac{^{}(𝐱,\sigma )}{r}\right]=\frac{f(𝐱,\sigma )}{\sigma };$$ (5B.64) and, from Eq. (5B.28) in Lem. 31, $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{}{r}}\left\{\nu ^+(𝐱,\sigma )^1[^+(𝐱,\sigma )^{}(𝐱,\sigma )]\right\}=`$ (5B.65) $`{\displaystyle \frac{1}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}{\displaystyle \frac{[\nu ^+(𝐱,\sigma ^{})^1f(𝐱,\sigma ^{})]/r}{\sigma ^{}\sigma }}.`$ A comparison of the above Eqs. (5B.64) and (5B.65) with Eqs. (5B.64) and (5B.63), together with the fact that $$\underset{\zeta 0}{lim}\frac{\nu (𝐱,\sigma \pm \zeta )^1}{r}=\frac{\nu ^\pm (𝐱,\sigma )^1}{r},$$ now yields $$\underset{\zeta 0}{lim}\frac{\stackrel{`}{}(𝐱,\sigma \pm \zeta )}{r}=\frac{^\pm (𝐱,\sigma )}{r}.$$ (5B.66) Statements (5B.60) and (5B.66) complete the proof of part (i) of our theorem for $`(𝐱,\tau )/r`$. End of proof. * Since $$\nu ^+(𝐱,\sigma )^1=\frac{^+(\sigma r)^+(\sigma s)}{^+(\sigma r_0)^+(\sigma s_0)},$$ (5B.67) one has $$\nu ^+(𝐱,\sigma )^1f(𝐱,\sigma )=0\text{ when }\sigma =r\text{ and when }\sigma =s.$$ (5B.68) Therefore, from statement $`1^0`$ in Sec. 29 of Muskhelishvili’s treatise, and from our Eq. (5B.53), $$\nu (𝐱,\tau )^1(𝐱,\tau )\text{ converges as }\tau r\text{ and as }\tau s[\tau C\overline{}(𝐱)].$$ (5B.69) Furthermore, from Eqs. (5B.15), (5B.16) and (5B.53), $`(𝐱,\tau )=\nu (𝐱,\tau )I+[\nu (𝐱,\tau )1]f(𝐱,r_0)`$ (5B.70) $`{\displaystyle \frac{\nu (𝐱,\tau )}{\pi i}}{\displaystyle _\overline{}}𝑑\sigma ^{}\nu ^+(𝐱,\sigma ^{})^1\left[{\displaystyle \frac{f(𝐱,\sigma ^{})f(𝐱,r_0)}{\sigma ^{}r}}\right].`$ From statement (5B.31), $`f(𝐱,\sigma ^{})/\sigma ^{}`$ exists and is a continuous function of $`\sigma ^{}`$ throughout $`\overline{}(𝐱)`$. Therefore, as one can see from Eq. (5B.67), $$\nu ^+(𝐱,\sigma )^1[f(𝐱,\sigma ^{})f(𝐱,r_0)]=0\text{ when }\sigma =r_0;$$ (5B.71) and it then follows from Eq. (5B.70) and the same statement $`1^0`$ in Sec. 29 of Muskhelishvili that was used before that $$(𝐱,\tau )\text{ converges [to }f(𝐱,r_0)\text{] as }\tau r_0.$$ (5B.72) Similarly, one proves that $$(𝐱,\tau )\text{ converges [to }f(𝐱,s_0)\text{] as }\tau s_0.$$ (5B.73) Statements (5B.69), (5B.72) and (5B.73) together constitute part (ii) of our theorem. End of proof. * We shall prove this part of our theorem for $`i=3`$, and the proof for $`i=4`$ is left to the reader. We start with the definition (5B.55) of $`\mathrm{\Phi }(𝐱,\tau )`$. The proof that we have just given for part (ii) of this theorem is also applicable to $`\mathrm{\Phi }(𝐱,\tau )`$. Specifically, the proof of part (ii) remains valid if one makes all of the following substitutions in its wording and equations: $`f(𝐱,\sigma ^{})`$ $``$ $`\varphi (𝐱,\sigma ^{}),`$ Eq. (5B.53) $``$ $`\text{Eq. (}\text{5B.55}\text{)},`$ $`\nu (𝐱_0,𝐱,\tau )(𝐱,\tau )`$ $``$ $`\mathrm{\Phi }(𝐱,\tau ),`$ statement (5B.31) $``$ $`\text{condition (}\text{5B.57}\text{)}.`$ Therefore, the conclusion of part (ii) of our theorem remains valid if one makes the substitution ‘$`\nu (𝐱,\tau )^1(𝐱,\tau )`$$``$$`\mathrm{\Phi }(𝐱,\tau )`$’. So, for all $`(𝐱,\tau )\mathrm{d}om`$, $$\begin{array}{c}\mathrm{\Phi }(𝐱,\tau )\text{ converges as }\tau r\text{ and as }\tau s,\text{ and }\hfill \\ \nu (𝐱,\tau )\mathrm{\Phi }(𝐱,\tau )\text{ converges as }\tau r_0\text{ and as }\tau s_0.\hfill \end{array}$$ (5B.74) When the above statement (5B.74) is applied to Eq. (5B.54), one obtains the statement in part (iii) of our theorem for the case $`i=3`$. End of proof. Note: The meanings that we assigned above to ‘$`f(𝐱,\sigma ^{})`$’, ‘$`\varphi (𝐱,\sigma ^{})`$’ and ‘$`\mathrm{\Phi }(𝐱,\tau )`$’ will not be used in the remainder of these notes. They were temporary devices for the purpose of abbreviating the proofs of the preceding theorem and two lemmas. ## 6 Proof of the generalized Geroch conjecture ### A. Generalized Abel transforms of the initial data and the identification of the sets $`𝒮_{}^{\mathrm{}}`$ and $`𝒮_{}^{\mathrm{}}`$ In Sec. 1A, we introduced a linear system $`_{HE}`$ for the Ernst equation that is related to $`=_{KC}`$ by Eqs. (1A.24) to (1A.26). It will now be useful to introduce one more linear system $`\stackrel{~}{}_{HE}`$ such that $$\stackrel{~}{}_{HE}(𝐱,\tau ):=P^M(𝐱_0,\tau )_{HE}(𝐱,\tau )P^M(𝐱_0,\tau )^1,$$ (6A.1) whereupon Eq. (1A.24) and the fact that $$^M(𝐱,\tau )=P^M(𝐱,\tau )P^M(𝐱_0,\tau )^1$$ (6A.2) yields $$(𝐱,\tau )=A(𝐱)^M(𝐱,\tau )\stackrel{~}{}_{HE}(𝐱,\tau ),$$ (6A.3) where $$A:=\frac{1}{\sqrt{h_{22}}}\left(\begin{array}{cc}1& h_{12}\\ 0& h_{22}\end{array}\right).$$ (6A.4) Note that $$h=Ah^MA^T,h^M=\left(\begin{array}{cc}\rho ^2& 0\\ 0& 1\end{array}\right).$$ (6A.5) Therefore, from Eqs. (4B.2) to (4B.5), and the fact that $`A^T\mathrm{\Omega }A=(detA)\mathrm{\Omega }=\mathrm{\Omega }`$, $$\left[\stackrel{~}{}_{HE}(𝐱,\tau ^{})\right]^{}𝒜^M(𝐱_0,\tau )\stackrel{~}{}_{HE}(𝐱,\tau )=𝒜^M(𝐱_0,\tau ).$$ (6A.6) Obviously, $`d\stackrel{~}{}_{HE}=\stackrel{~}{\mathrm{\Gamma }}_{HE}\stackrel{~}{}_{HE}`$, where $$\stackrel{~}{\mathrm{\Gamma }}_{HE}(𝐱,\tau )=P^M(𝐱_0,\tau )\mathrm{\Gamma }_{HE}(𝐱,\tau )P^M(𝐱_0,\tau )^1,$$ (6A.7) and $`\mathrm{\Gamma }_{HE}`$ is given by Eq. (1A.13). Note that $`\stackrel{~}{\mathrm{\Gamma }}_{HE}`$ can be obtained by making the following substitutions in $`\mathrm{\Gamma }_{HE}`$: $`J`$ $``$ $`\stackrel{~}{J}(\tau ):=P^M(𝐱_0,\tau )JP^M(𝐱_0,\tau )^1=\left(\begin{array}{cc}i& 2(\tau z_0)\\ 0& i\end{array}\right),`$ (6A.10) $`𝒩(\tau )`$ $``$ $`\stackrel{~}{𝒩}(\tau ):=P^M(𝐱_0,\tau )𝒩(\tau )P^M(𝐱_0,\tau )^1=\left(\begin{array}{cc}\tau +z_0& i\rho _0^2\\ i& \tau z_0\end{array}\right),`$ (6A.13) where $$J:=i\sigma _2\text{ and }𝒩(\tau ):=\mu (𝐱_0,\tau )\sigma _3.$$ (6A.14) The properties of $`\stackrel{~}{}_{HE}`$ can be deduced from those of $`_{HE}`$. For example, consider the generalized Abel transforms<sup>31</sup><sup>31</sup>31In the Weyl case the $`\alpha ^{(i)}`$ are easily expressed in terms of Abel transforms. (our term) $$\begin{array}{ccc}\hfill \alpha ^{(3)}(\sigma )& :=& _{HE}^+((\sigma ,s_0),\sigma )^1\text{ for }\sigma ^{(3)}\text{ and }\hfill \\ \hfill \alpha ^{(4)}(\sigma )& :=& _{HE}^+((r_0,\sigma ),\sigma )^1\text{ for }\sigma ^{(4)}\hfill \end{array}$$ (6A.15) of the initial data functions $$\begin{array}{ccc}\hfill ^{(3)}(r)& =& (r,s_0)\text{ for }r^{(3)}\text{ and }\hfill \\ \hfill ^{(4)}(s)& =& (r_0,s)\text{ for }s^{(4)}.\hfill \end{array}$$ (6A.16) Analysis<sup>32</sup><sup>32</sup>32For details, see our Magnum Opus (gr-qc/9903104). yields $$\alpha ^{(i)}=I\alpha _0^{(i)}+J\alpha _1^{(i)}+𝒩^+(\sigma )[I\alpha _2^{(i)}+J\alpha _3^{(i)}],$$ (6A.17) where $`𝒩^+(\sigma )=\mu ^+(𝐱_0,\sigma )\sigma _3`$, $`\alpha _k^{(i)}:^{(i)}R^1(k=0,1,2,3),`$ (6A.18) $`\alpha _k^{(i)}\text{ is }H(1/2)\text{ on each closed subinterval of }^{(i)},`$ (6A.19) $$\alpha _k^{(i)}\text{ is }𝐂^{n1}\text{ if }^{(i)}\text{ is }𝐂^n\text{ and }\alpha _k^{(i)}\text{ is analytic if }^{(i)}\text{ is analytic},$$ (6A.20) and $$det\alpha ^{(i)}=[\alpha _0^{(i)}]^2+[\alpha _1^{(i)}]^2+(\sigma r_0)(\sigma s_0)\left\{[\alpha _2^{(i)}]^2+[\alpha _3^{(i)}]^2\right\}=1.$$ (6A.21) Instead of $`\alpha ^{(3)}`$ and $`\alpha ^{(4)}`$, we shall be employing $$\begin{array}{ccc}\hfill V^{(3)}(\sigma )& :=& \stackrel{~}{}_{HE}^+((\sigma ,s_0),\sigma )^1\text{ for }\sigma ^{(3)}\text{ and }\hfill \\ \hfill V^{(4)}(\sigma )& :=& \stackrel{~}{}_{HE}^+((r_0,\sigma ),\sigma )^1\text{ for }\sigma ^{(4)},\hfill \end{array}$$ (6A.22) whose pertinent properties are easily deduced from those of $`\alpha ^{(3)}`$ and $`\alpha ^{(4)}`$ by using Eq. (6A.1). For example, $$V^{(i)}=I\alpha _0^{(i)}+\stackrel{~}{J}(\sigma )\alpha _1^{(i)}+\stackrel{~}{𝒩}(\sigma )[I\alpha _2^{(i)}+\stackrel{~}{J}(\sigma )\alpha _3^{(i)}].$$ (6A.23) Furthermore, with the aid of Eq. (6A.6) and the definitions of $`K`$ and $`K^{\mathrm{}}`$ by Eqs. (1C.3) to (1C.6), one readily deduces from Eqs. (6A.18) to (6A.21) that $$𝐕K\text{ where }𝐕:=(V^{(3)},V^{(4)}),$$ (6A.24) and $`𝐕K^{n1}\text{ if }^{(3)}\text{ and }^{(4)}\text{ are }𝐂^n,`$ $`𝐕K^{\mathrm{}}\text{ if }^{(3)}\text{ and }^{(4)}\text{ are }𝐂^{\mathrm{}}\text{ and}`$ $`𝐕K^{an}\text{ if }^{(3)}\text{ and }^{(4)}\text{ are }𝐂^{an}.`$ (6A.25) Defining $`𝒮_𝐕`$ $`:=`$ $`\text{the set of all ordered pairs }𝐕=(V^{(3)},V^{(4)}),`$ where $`V^{(i)}`$ is a $`2\times 2`$ matrix function with the $`\text{domain }^{(i)}\text{ and there exists }𝒮_{}`$ such that Eqs. (6A.22) hold, $`B(^{(i)})`$ $`:=`$ the multiplicative group of all $`\mathrm{exp}\left(\stackrel{~}{J}𝝋^{(i)}\right)=I\mathrm{cos}𝝋^{(i)}+\stackrel{~}{J}\mathrm{sin}𝝋^{(i)}`$ such that $`𝝋^{(i)}`$ is any real-valued function that has the domain $`^{(i)}`$ and is $`H(1/2)`$ on every closed subinterval of $`^{(i)}`$, it will turn out to be possible to identify the sets $`𝒮_{}^{\mathrm{}}`$ involved in the generalized Geroch conjecture in terms of the more fundamental sets $$𝒮_𝐕^{\mathrm{}}:=\{𝐕𝒮_𝐕:\text{ there exists }𝐰B(^{(3)})\times B(^{(4)})\text{ for which }\mathrm{𝐕𝐰}k^{\mathrm{}}\},$$ (6A.28) where $$k^{\mathrm{}}=kK^{\mathrm{}},k:=\{\mathrm{𝐕𝐰}:𝐕𝒮_𝐕,𝐰B(^{(3)})\times B(^{(4)})\},$$ (6A.29) and, for any members $`𝐯=(v^{(3)},v^{(4)})`$ and $`𝐯^{}=(v^{(3)},v^{(4)})`$ of $`K`$, $$\mathrm{𝐯𝐯}^{}:=(v^{(3)}v^{(3)},v^{(4)}v^{(4)}).$$ Specifically, we let $$𝒮_{}^{\mathrm{}}:=\text{ the set of all }𝒮_{}\text{ for which }𝐕𝒮_𝐕^{\mathrm{}}.$$ (6A.30) Having defined $`𝒮_{}^{\mathrm{}}`$, we can easily identify the remaining important sets. Thus, $$𝒮_{}^{\mathrm{}}:=\text{ the set of all }𝒮_{}\text{ for which }𝒮_{}^{\mathrm{}},$$ (6A.31) with a like definition of $`𝒮_H^{\mathrm{}}`$. We leave the proof of the following theorem, which actually motivated how we formulated our HHP corresponding to $`(𝐯,_0)`$, to the reader: ###### THEOREM 33 (Motivation) For all $`𝐯K`$ and for all $`𝒮_{}`$ members $``$ and $`_0`$ whose corresponding $`𝒮_𝐕`$ members are $`𝐕`$ and $`𝐕_0`$, respectively, the following statements (i) and (ii) are equivalent to one another: * There exists $`𝐰B(^{(3)})\times B(^{(4)})`$ such that $$𝐯=\mathrm{𝐕𝐰𝐕}_0^1.$$ (6A.32a) * For each $`𝐱D`$, $`i\{3,4\}`$ and $`\sigma ^{(i)}(𝐱)`$, $$^+(𝐱,\sigma )v^{(i)}(\sigma )[_0^+(𝐱,\sigma )]^1=^{}(𝐱,\sigma )v^{(i)}(\sigma )[_0^{}(𝐱,\sigma )]^1.$$ (6A.32b) Moreover, if $`^{(i)}`$ and $`_0^{(i)}`$ are $`𝐂^{n_i}`$ (resp. analytic) and $`w^{(i)}`$ is $`𝐂^{n_i1}`$ (resp. analytic), then the function of $`\sigma `$ that equals each side of Eq. (6A.32b) has a $`𝐂^{n_i1}`$ (resp. analytic) extension $`Y^{(i)}(𝐱)`$ to the interval $$\mathrm{d}omY^{(i)}(𝐱)=\stackrel{ˇ}{}^{(i)}(x^{7i})$$ (6A.32c) and, if $`𝐯K^{\mathrm{}}`$ and $`_0𝒮_{}^{\mathrm{}}`$, then $`𝐕𝒮_𝐕^{\mathrm{}}`$ and $`𝒮_{}^{\mathrm{}}`$. ###### THEOREM 34 (Relation of $`_0`$ and $`𝐕_0`$) For each $`_0𝒮_{}`$ whose corresponding member of $`𝒮_𝐕`$ is $`𝐕_0`$, and for each $`𝐰B(^{(3)})\times B(^{(4)})`$, $`_0`$ is a solution of the HHP corresponding to $`(𝐕_0𝐰,^M)`$. Proof: For each $`𝐱D`$ and $`i\{3,4\}`$, * Thm. 1(i) states that $`_0(𝐱)`$ is holomorphic throughout its domain $`C\overline{}(𝐱)`$, * Thm. 2 states that $`^\pm (𝐱)`$ exist and, from Thm. 33 and the fact that $`𝐕^M=(I,I)`$, $`Y_0^{(i)}(𝐱,\sigma )`$ $`:=`$ $`_0^+(𝐱,\sigma )V_0^{(i)}(\sigma )w^{(i)}(\sigma )[^{M+}(𝐱,\sigma )]^1`$ (6A.33a) $`=`$ $`_0^{}(𝐱,\sigma )V_0^{(i)}(\sigma )w^{(i)}(\sigma )[^M(𝐱,\sigma )]^1`$ $`\text{for all }\sigma ^{(i)}(𝐱);`$ and Thm. 2 and Thm. 1(iii) imply that $`_0(𝐱)`$ is bounded at $`𝐱_0`$ and $`\nu (𝐱)^1_0(𝐱)`$ is bounded at $`𝐱`$, while the function $`Y_0(𝐱)`$ whose domain is $`(𝐱)`$ and whose values are given by $`Y_0^{(i)}(𝐱,\sigma )`$ at each $`\sigma ^{(i)}(𝐱)`$ satisfies the condition $$Y_0(𝐱)\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{.}$$ (6A.33b) Thus, $`_0`$ is a solution of the HHP corresponding to $`(\mathrm{𝐕𝐰},^M)`$. End of proof. ###### THEOREM 35 (Reduction theorem) For each $`𝐱D`$ and $`2\times 2`$ matrix function $`(𝐱)`$ with the domain $`C\overline{}(𝐱)`$, for each $`𝐯K`$ and $`_0𝒮_{}`$ whose corresponding member of $`𝒮_𝐕`$ is $`𝐕_0`$, and for each $`𝐰B(^{(3)})\times B(^{(4)})`$, the following two statements are equivalent to one another: * The function $`(𝐱)`$ is a solution of the HHP corresponding to $`(𝐯,_0,𝐱)`$. * The function $`(𝐱)`$ is a solution of the HHP corresponding to $`(\mathrm{𝐯𝐕}_0𝐰,^M,𝐱)`$. Proof: Suppose that statement (i) is true. Then $`(𝐱)`$ satisfies all four conditions (1) through (4) in the definition of the HHP corresponding to $`(𝐯,_0,𝐱)`$. In particular, from conditions (3) and (4), $`Y^{(i)}(𝐱,\sigma )`$ $`:=`$ $`^+(𝐱,\sigma )v^{(i)}(\sigma )[_0^+(𝐱,\sigma )]^1`$ (6A.34a) $`=`$ $`^{}(𝐱,\sigma )v^{(i)}(\sigma )[_0^{}(𝐱,\sigma )]^1`$ $`\text{for all }i\{3,4\}\text{ and }\sigma ^{(i)}(𝐱);`$ and $`Y(𝐱)\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{.}`$ (6A.34b) So, from the preceding Thm. 34 and Eqs. (6A.33a) and (6A.34a), $`X^{(i)}(𝐱,\sigma )`$ $`:=`$ $`^+(𝐱,\sigma )u^{(i)}(\sigma )[^{M+}(𝐱,\sigma )]^1`$ (6A.34c) $`=`$ $`_0^{}(𝐱,\sigma )u^{(i)}(\sigma )[^M(𝐱,\sigma )]^1`$ $`\text{for all }i\{3,4\}\text{ and }\sigma ^{(i)}(𝐱),`$ where $$𝐮:=𝐯(𝐕_0𝐰)$$ (6A.34d) (which is a member of $`K`$, since $`𝒮_𝐕K`$ and $`BK`$); and, furthermore, $$X(𝐱)=Y(𝐱)Y_0(𝐱)$$ (6A.34e) and, from (6A.33b), (6A.34b) and (6A.34e), $$X(𝐱)\text{ is bounded at }𝐱\text{ and }𝐱_0\text{.}$$ (6A.34f) Therefore, we have proved that statement (ii) is true if statement (i) is true. Next, suppose statement (ii) is true. Then $`(𝐱)`$ satisfies all four conditions in the definition of the HHP corresponding to $`(𝐮,^M,𝐱)`$, where $`𝐮`$ is defined by Eq. (6A.34d). In particular, from conditions (3) and (4), Eq. (6A.34c) and the statement (6A.34f) hold. Since $`detV_0^{(i)}=detw^{(i)}=1`$ and since $`det_0(𝐱)=det^M(𝐱)=\nu (𝐱)`$ \[Thm. 1(iii)\], Eq. (6A.33a) yields $`detY^{(i)}(𝐱)=1`$. Therefore, both sides of Eq. (6A.33a) are invertible, and $`[Y_0^{(i)}(𝐱,\sigma )]^1`$ $`=`$ $`^{M+}(𝐱,\sigma )[V_0^{(i)}(\sigma )w^{(i)}(\sigma )]^1[_0^+(𝐱,\sigma )]^1`$ (6A.34g) $`=`$ $`^M(𝐱,\sigma )[V_0^{(i)}(\sigma )w^{(i)}(\sigma )]^1[_0^{}(𝐱,\sigma )]^1`$ $`\text{for all }i\{3,4\}\text{ and }\sigma ^{(i)}(𝐱);`$ and, from (6A.33b), $$Y_0(𝐱)^1\text{ is bounded at }𝐱\text{ and at }𝐱_0\text{.}$$ (6A.34h) So, by multiplying both sides of Eq. (6A.34c) by the corresponding sides of Eq. (6A.34g), and then using (6A.34d), (6A.34f) and (6A.34h), we establish that $``$ is a solution of the HHP corresponding to $`(𝐯,_0,𝐱)`$. End of proof. ### B. The HHP solution $``$ is a member of $`𝒮_{}^{\mathrm{}}`$ when $`𝐯K^{\mathrm{}}`$ and $`\mathrm{}`$ is $`n3`$, $`n+(n3)`$, $`\mathrm{}`$ or ‘an’ ###### THEOREM 36 ($`\stackrel{`}{}/x^i=\mathrm{\Gamma }_i\stackrel{`}{}`$) When $`𝐯K^3`$, $``$ is the solution of the HHP corresponding to $`(𝐯,^M)`$ and $`H`$ is the function defined by Eq. (4D.1) in Thm. 23, then \[from Thm. 29\] $`d\stackrel{`}{}(𝐱,\tau )`$ and $`dH(𝐱)`$ exist; and, for each $`i\{3,4\}`$, $$\frac{\stackrel{`}{}(𝐱,\tau )}{x^i}=\mathrm{\Gamma }_i(𝐱,\tau )\stackrel{`}{}(𝐱,\tau )\text{ for all }(𝐱,\tau )\mathrm{d}om\stackrel{`}{},$$ (6B.1) where $$\mathrm{\Gamma }_i(𝐱,\tau ):=\frac{1}{2(\tau x^i)}\frac{H(𝐱)}{x^i}\mathrm{\Omega }.$$ (6B.2) Proof: From Thm. 3(ii), $`(𝐱,\tau )^1`$ exists for all $`(𝐱,\tau )\mathrm{d}om`$; and, for the continuous extension of $`Y`$ that is defined by Cor. 10 (also, see the beginning of Sec. 4F) and Eq. (2B.35), $`Y(𝐱,\sigma )^1`$ exists for all $`𝐱D`$ and $`\sigma \overline{}(𝐱)`$. From Thm. 29, $`d\stackrel{`}{}(𝐱,\tau )`$, $`dY(𝐱,\sigma )`$ and $`dH(𝐱)`$ exist and are continuous functions of $`(𝐱,\tau )`$, $`(𝐱,\sigma )`$ and $`𝐱`$ throughout $`\mathrm{d}om\stackrel{`}{}`$, $`\mathrm{d}omY:=\{(𝐱,\sigma ):𝐱D,\sigma \overline{}(𝐱)\}`$ and $`D`$, respectively; and, for each $`𝐱D`$, $`d\stackrel{`}{}(𝐱,\tau )`$ is a holomorphic function of $`\tau `$ throughout $`C(𝐱)\{r,s,r_0,s_0\}`$. It then follows, with the aid of conditions (1) through (3) in the definition of the HHP, Eq. (4D.1) in Thm. 23, and Thm. 32(i) that, for each $`𝐱D`$, $`Z_i(𝐱,\tau )`$ $`:=`$ $`(\tau x^i){\displaystyle \frac{\stackrel{`}{}(𝐱,\tau )}{x^i}}\stackrel{`}{}(𝐱,\tau )^1\text{ exists and is a holomorphic}`$ (6B.3) $`\text{function of }\tau \text{ throughout }C(𝐱)\{r,s,r_0,s_0\}`$ $`Z_i(𝐱,\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{H(𝐱)}{x^i}}\mathrm{\Omega }+O(\tau ^1)\text{ in at least}`$ (6B.4) $`\text{one neighborhood of }\tau =\mathrm{},`$ $`Z_i^\pm (𝐱,\sigma )`$ $`\text{exists for each }\sigma (𝐱)`$ (6B.5) and $`Z_i^+(𝐱,\sigma )`$ $`=`$ $`Z_i^{}(𝐱,\sigma )`$ (6B.6) $`=`$ $`(\sigma x^i){\displaystyle \frac{Y(𝐱,\sigma )}{x^i}}Y(𝐱,\sigma )^1+Y(𝐱,\sigma ){\displaystyle \frac{1}{2}}{\displaystyle \frac{H^M(𝐱)}{x^i}}\mathrm{\Omega }Y(𝐱,\sigma )^1`$ $`\text{for all }\sigma \overline{}(𝐱),`$ where we have used the fact that the defining equation in condition (3) for the HHP corresponding to $`(𝐯,^M,𝐱)`$ is expressible in the form $$^\pm (𝐱,\sigma )=Y^{(j)}(𝐱,\sigma )^{M\pm }(𝐱,\sigma )[v^{(j)}(\sigma )]^1\text{ for all }\sigma ^{(j)}(𝐱);$$ (6B.7) and we have used the fact that, since $`^M𝒮_{}`$, $$\frac{\stackrel{`}{}^M(𝐱,\tau )}{x^j}=\mathrm{\Gamma }_i^M(𝐱,\tau )\stackrel{`}{}^M(𝐱,\tau )\text{ for all }\tau C(𝐱)\{r,s,r_0,s_0\}.$$ (6B.8) We next define a continuous extension of $`Z_i(𝐱)`$ \[which we also denote by $`Z_i(𝐱)`$\] to the domain $`C\{r,s,r_0,s_0\}`$ by letting $$Z_i(𝐱,\sigma ):=Z_i^\pm (𝐱,\sigma ).$$ (6B.9) Then, from the statement (6B.3) and the theorem of Riemann that we have already used in a different context, $$Z_i(𝐱,\tau )\text{ is a holomorphic function of }\tau \text{ throughout }C\{r,s,r_0,s_0\}.$$ (6B.10) However, from Eq. (6B.3) and Thms. 32(ii) and (iii), $$\begin{array}{c}\nu (𝐱,\tau )Z_i(𝐱,\tau )\text{ converges as }\tau r_0\text{ and as }\tau s_0,\hfill \\ \text{and }\nu (𝐱,\tau )^1Z_i(𝐱,\tau )\text{ converges as }\tau r\text{ and as }\tau s.\hfill \end{array}$$ Also, from Eq. (6B.6) and the continuity on $`\overline{}(𝐱)`$ of $`dY(𝐱,\sigma )`$ and $`Y(𝐱,\sigma )^1=\mathrm{\Omega }Y(𝐱,\sigma )^T\mathrm{\Omega }`$, $$Z_i(𝐱,\sigma )\text{ converges as }\sigma r_0,\sigma s_0,\sigma r,\sigma s.$$ (6B.11) Combining (6B.10), (B.) and (6B.11), one obtains, by reasoning that should now be familiar to us, $`Z_i(𝐱,\tau )=Z_i(𝐱,\mathrm{})`$, whereupon the conclusion of our theorem follows from Eqs. (6B.3) and (6B.4). End of proof. ###### COROLLARY 37 ($`d\stackrel{`}{}=\mathrm{\Gamma }\stackrel{`}{}`$) For each $`(𝐱,\tau )\mathrm{d}om\stackrel{`}{}`$, $$d\stackrel{`}{}(𝐱,\tau )=\mathrm{\Gamma }(𝐱,\tau )\stackrel{`}{}(𝐱,\tau ),$$ (6B.12) where $`\mathrm{\Gamma }(𝐱,\tau )`$ $`:=`$ $`{\displaystyle \frac{1}{2}}(\tau z+\rho )^1dH(𝐱)\mathrm{\Omega }`$ (6B.13) $`=`$ $`{\displaystyle \underset{i}{}}dx^i\mathrm{\Gamma }_i(𝐱,\tau ).`$ Proof: Obvious. End of proof. ###### THEOREM 38 ($`𝒜\mathrm{\Gamma }=\frac{1}{2}\mathrm{\Omega }dH\mathrm{\Omega }`$) Suppose $`𝐯K^3`$ and $``$ is the solution of the HHP corresponding to $`(𝐯,^M)`$. Then $$𝒜\mathrm{\Gamma }=\frac{1}{2}\mathrm{\Omega }dH\mathrm{\Omega },$$ (6B.14) where $`H`$, $`𝒜`$ and $`\mathrm{\Gamma }`$ are defined by Eqs. (4D.1), (4D.9) and (6B.13), respectively. Proof: The proof will be given in three parts: * For each $`H^{}𝒮_H`$, note that $$\mathrm{R}eH^{}=h^{}$$ (6B.15) and that the defining differential equation for $`\mathrm{I}mH^{}`$ in terms of $`\mathrm{R}eH^{}`$ is expressible in the form $$h^{}\mathrm{\Omega }d(\mathrm{R}eH^{})=\rho (i\mathrm{I}mH^{}).$$ (6B.16) Recall that $`h^{}`$ is symmetric and $`deth^{}=\rho ^2`$. So, $$(h^{}\mathrm{\Omega })^2=\rho ^2I.$$ (6B.17) From Eq. (6B.17), Eq. (6B.16) is equivalent to the equation $$h^{}\mathrm{\Omega }d(i\mathrm{I}mH^{})=\rho d(\mathrm{R}eH^{});$$ (6B.18) and, therefore, Eq. (6B.16) is equivalent to the equation $$h^{}\mathrm{\Omega }dH^{}=\rho dH^{}.$$ (6B.19) Furthermore, the above Eq. (6B.19) yields $$𝒜^{}\mathrm{\Gamma }^{}=[(\tau z)\mathrm{\Omega }+\mathrm{\Omega }h^{}\mathrm{\Omega }]\frac{1}{2}(\tau z+\rho )^1dH^{}\mathrm{\Omega }=\frac{1}{2}(\tau z+\rho )^1[(\tau z)\mathrm{\Omega }dH^{}\mathrm{\Omega }+\rho \mathrm{\Omega }dH^{}\mathrm{\Omega }].$$ So, Eq. (6B.19) implies $$𝒜^{}\mathrm{\Gamma }^{}=\frac{1}{2}\mathrm{\Omega }dH^{}\mathrm{\Omega }\text{ for each }H^{}𝒮_H.$$ (6B.20) The reader should have no difficulty in proving that, conversely, Eq. (6B.20) implies (6B.19). We shall use the above result later in our proof. * We now obtain a second result that we shall need for the proof. From Eq. (4D.10) in Thm. 24, $$[^{}(𝐱,\sigma )]^{}𝒜(𝐱,\sigma )^\pm (𝐱,\sigma )=𝒜(𝐱_0,\sigma )\text{ for all }\sigma (𝐱).$$ (6B.21) Now, recall that $`𝒜(𝐱_0,\sigma )=𝒜^M(𝐱_0,\sigma )`$ in our gauge. Also, recall that $$[^{}(𝐱,\tau ^{})]^{}𝒜^{}(𝐱,\tau )^{}(𝐱,\tau )=𝒜^M(𝐱_0,\sigma )\text{ for all }^{}𝒮_{}.$$ (6B.22) Therefore, we obtain the following result by using Eqs. (6B.7) \[condition (3) in the definition of the HHP corresponding to $`(𝐯,^M,𝐱)`$\], (6B.21), (1C.4) and (6B.22) \[for $`^{}=^M`$\]: $$Y^{}(𝐱,\sigma )𝒜(𝐱,\sigma )Y(𝐱,\sigma )=𝒜^M(𝐱,\sigma )\text{ for all }\sigma (𝐱).$$ However, recall that $`Y(𝐱,\sigma )`$ is now a continuous function of $`\sigma `$ throughout $`\overline{}(𝐱)`$. Therefore, $$Y^{}(𝐱,\sigma )𝒜(𝐱,\sigma )Y(𝐱,\sigma )=𝒜^M(𝐱,\sigma )\text{ for all }\sigma \overline{}(𝐱).$$ (6B.23) We shall use this result below. * From the definition of $`𝒜`$ and $`\mathrm{\Gamma }`$, each component of $`𝒜(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )`$ is a holomorphic function of $`\tau `$ throughout $`C\{r,s\}`$ and has no essential singularity at $`\tau =r`$ and at $`\tau =s`$. In fact, if there are any singularities at these points, they are simple poles. That much is obvious. From Eqs. (6B.12), (6B.7) and (6B.8), $`𝒜(𝐱,\sigma )\mathrm{\Gamma }(𝐱,\sigma )`$ $`=`$ $`𝒜(𝐱,\sigma )[d^\pm (𝐱,\sigma )][^\pm (𝐱,\sigma )]^1`$ (6B.24) $`=`$ $`𝒜(𝐱,\sigma )\left[dY(𝐱,\sigma )+Y(𝐱,\sigma )\mathrm{\Gamma }^M(𝐱,\sigma )\right][Y(𝐱,\sigma )]^1`$ $`\text{for all }\sigma (𝐱).`$ The above equation becomes, after using Eq. (6B.23), $$𝒜(𝐱,\sigma )\mathrm{\Gamma }(𝐱,\sigma )=\left\{𝒜(𝐱,\sigma )dY(𝐱,\sigma )+[Y(𝐱,\sigma )^{}]^1𝒜^M(𝐱,\sigma )\mathrm{\Gamma }^M(𝐱,\sigma )\right\}[Y(𝐱,\sigma )]^1,$$ which becomes, after using Eq. (6B.20) with $`H^{}=H^M`$, $$𝒜(𝐱,\sigma )\mathrm{\Gamma }(𝐱,\sigma )=\left\{𝒜(𝐱,\sigma )dY(𝐱,\sigma )+[Y(𝐱,\sigma )^{}]^1\frac{1}{2}\mathrm{\Omega }dH^M(𝐱)\mathrm{\Omega }\right\}[Y(𝐱,\sigma )]^1.$$ (6B.25) From Thm. 29 and the fact that $`detY(𝐱,\sigma )=1`$, the right side of the above equation is a continuous function of $`\sigma `$ throughout $`\overline{}(𝐱)`$. Therefore, $`𝒜(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )`$ is extendable to a holomorphic function of $`\tau `$ throughout $`C`$; and it follows that $$𝒜(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )=[𝒜(𝐱,\tau )\mathrm{\Gamma }(𝐱,\tau )]_{\tau =\mathrm{}}=\frac{1}{2}\mathrm{\Omega }dH(𝐱)\mathrm{\Omega }.$$ End of proof. ###### COROLLARY 39 ($`h\mathrm{\Omega }dH=\rho dH`$) Suppose $`H`$ is defined as in the preceding theorem. Then $$h\mathrm{\Omega }dH=\rho dH.$$ (6B.26) Proof: Multiply both sides of Eq. (6B.14) through by $`2\mathrm{\Omega }(\tau z+\rho )`$ on the left, and by $`\mathrm{\Omega }`$ on the right; and then set $`\tau =z`$. End of proof. ###### THEOREM 40 (HHP solution $`𝒮_{}^{\mathrm{}}`$) * For each $`𝐯K^{\mathrm{}}`$, where $`\mathrm{}`$ is $`n3`$, $`n+`$ ($`n3`$), $`\mathrm{}`$ or ‘an’, and, for each $`_0𝒮_{}^{\mathrm{}}`$, there exists exactly one solution $``$ of the HHP corresponding to $`(𝐯,_0)`$. * Let $`H`$ be defined in terms of $``$ by Eq. (4D.1), and let $`:=H_{22}`$. Then $`𝒮_{}`$, and $`H`$ is identical with the unique member of $`𝒮_H`$ that is constructed from $``$ in the usual way. * Furthermore, $``$ is identical with the member of $`𝒮_{}`$ that is defined in terms of $`H`$ in Sec. 1 \[and whose existence and uniqueness for a given $`H𝒮_H`$ is asserted in Thm. 1.\] * Finally, let $`_{HE}`$ be defined in terms of $`_{KC}=`$ by Eq. (1A.24), and let $`𝐕`$ denote the member of $`𝒮_𝐕`$ that is defined in terms of $`\stackrel{~}{}_{HE}`$ by Eq. (6A.22). Then $`𝐕𝒮_𝐕^{\mathrm{}}`$ and, therefore (by definition), $`𝒮_{}^{\mathrm{}}`$, $`H𝒮_H^{\mathrm{}}`$ and $`𝒮_{}^{\mathrm{}}`$. Proofs: * Let $`𝐕_0`$ denote the member of $`𝒮_𝐕`$ that corresponds to $`_0`$. Since $`_0𝒮_{}^{\mathrm{}}`$, there exists (by definition of $`𝒮_{}^{\mathrm{}}`$) $`𝐰B(^3)\times B(^{(4)})`$ such that $$𝐕_0𝐰k^{\mathrm{}}K^{\mathrm{}};$$ (6B.27) and, since $`K^{\mathrm{}}`$ is a group, $$\mathrm{𝐯𝐕}_0𝐰K^{\mathrm{}}.$$ (6B.28) From Thm. 22, there exists exactly one solution $``$ of the HHP corresponding to $`(\mathrm{𝐯𝐕}_0𝐰,^M)`$; and, from Thm. 35, it then follows that $``$ is also a solution of the HHP corresponding to $`(𝐯,_0)`$. Finally, from Thm. 3(iv), there is no other solution of the HHP corresponding to $`(𝐯,_0)`$. End of proof. * From the premises of this theorem, $`𝐯K^3`$. Therefore, from statement (5B.7) in Thm. 29, $$dH\text{ exists and is continuous }$$ (6B.29) throughout $`D`$; and since $`(d^2H)(𝐱)`$ $`=`$ $`drds\left[{\displaystyle \frac{^2H(𝐱)}{rs}}{\displaystyle \frac{^2H(𝐱)}{sr}}\right]`$ (6B.30) and $`(ddH)(𝐱)`$ $`=`$ $`drds\left[{\displaystyle \frac{^2H(𝐱)}{rs}}+{\displaystyle \frac{^2H(𝐱)}{sr}}\right],`$ (6B.31) it is true that $$d^2H\text{ exists and vanishes }$$ (6B.32) and $$ddH\text{ exists and is continuous }$$ (6B.33) throughout $`D`$. Also, Eq. (6B.26) in Cor. 39 asserts that $$\rho dH=h\mathrm{\Omega }dH,$$ (6B.34) where we recall from Eq. (4D.2) in Thm. 23 that $$h:=\mathrm{R}eH=h^T$$ (6B.35) and, from Thm. 25, $$deth=\rho ^2\text{ and }f:=\mathrm{R}e=g_{22}<0,$$ (6B.36) where $`g_{ab}`$ denotes the element of $`h`$ in the $`a`$th row and $`b`$th column. Since $`=1`$, Eq. (6B.34) is equivalent to the equation $$\rho dH=h\mathrm{\Omega }dH$$ (6B.37) from which we obtain $$\rho dH^{}\mathrm{\Omega }dH=dH^{}\mathrm{\Omega }h\mathrm{\Omega }dH.$$ (6B.38) Upon taking the hermitian conjugates of the terms in the above equation, and upon noting that $`\mathrm{\Omega }^{}=\mathrm{\Omega }`$, $`h^{}=h`$, $$(\omega \eta )^T=\eta ^T\omega ^T\text{ and }\omega \eta =(\omega )\eta \text{ for any }n\times n\text{ matrix }1\text{-forms, }$$ (6B.39) one obtains $$\rho dH^{}\mathrm{\Omega }dH=dH^{}\mathrm{\Omega }h\mathrm{\Omega }dH.$$ (6B.40) From Eqs. (6B.40) and (6B.38), $$dH^{}\mathrm{\Omega }dH=0.$$ (6B.41) \[The above result (6B.41) was first obtained by the authors in a paper<sup>33</sup><sup>33</sup>33I. Hauser and F. J. Ernst, J. Math. Phys. 21, 1116-1140 (1980). See Eq. (37). which introduced an abstract geometric definition of the Kinnersley potential $`H`$ and which derived other properties of $`H`$ that we shall not need in these notes.\] We next consider the $`(2,2)`$ matrix elements of Eqs. (6B.34) and (6B.41). With the aid of Eqs. (6B.35) and (6B.36), one obtains $$\rho d=i(g_{12}d+fdH_{12})$$ (6B.42) and $$dH_{12}^{}dd^{}dH_{12}=0.$$ (6B.43) From Eq. (6B.42), $`fd(\rho d)\rho dd`$ $`=`$ $`if(dg_{12}d+dfdH_{12}ddH_{12})`$ $`=`$ $`if\left[{\displaystyle \frac{1}{2}}(dH_{12}+dH_{12}^{})d+{\displaystyle \frac{1}{2}}(d+d^{})dH_{12}ddH_{12}\right]`$ $`=`$ $`{\displaystyle \frac{if}{2}}(dH_{12}^{}d+d^{}dH_{12}).`$ Therefore, from Eq. (6B.43), $$fd(\rho d)\rho dd=0.$$ (6B.44) Furthermore, from Eqs. (1B.3), (4D.2) and (6B.36), $$(𝐱_0)=1\text{ and }\mathrm{R}e<0.$$ (6B.45) Therefore, $$𝒮_{},$$ (6B.46) since $``$ satisfies the Ernst equation (6B.44) and the requisite gauge conditions (6B.45) Next, let $$\chi :=\mathrm{I}m\text{ and }\omega :=g_{12}/g_{22}.$$ (6B.47) Then, by taking the imaginary parts of the terms in Eq. (6B.42), one deduces $$d\omega =\rho f^2d\chi .$$ (6B.48) Furthermore, Eqs. (6B.35) and (6B.36) enable us to express $`h`$ in the form $$h=A\left(\begin{array}{cc}\rho ^2& 0\\ 0& 1\end{array}\right)A^T,$$ (6B.49) where $$A:=\left(\begin{array}{cc}1& \omega \\ 0& 1\end{array}\right)\left(\begin{array}{cc}1/\sqrt{f}& 0\\ 0& \sqrt{f}\end{array}\right).$$ (6B.50) Finally, the imaginary parts of the terms in Eq. (6B.37) give us $$\rho d(\mathrm{I}mH)=hJdh.$$ (6B.51) A comparison of Eqs. (4D.2), (4D.4) and (6B.48) to (6B.51) with the definition of $`𝒮_H`$ that is given in Sec. 1 demonstrates that $`H`$ is precisely that member of $`𝒮_H`$ that is computed from $``$ in the usual way. End of proof. * From statement (5B.6) in Thm. 29, $$d\stackrel{`}{}(𝐱,\tau )\text{ exists for all }𝐱D\text{ and }\tau C(𝐱)\{r,s,r_0,s_0\};$$ (6B.52) and, from Cor. 37, $$d\stackrel{`}{}(𝐱,\tau )=\mathrm{\Gamma }(𝐱,\tau )\stackrel{`}{}(𝐱,\tau )\text{ for all }𝐱D\text{ and }\tau C(𝐱)\{r,s,r_0,s_0\}.$$ (6B.53) Furthermore, from Thm. 3(v), $$(𝐱_0,\tau )=I\text{ for all }\tau C.$$ (6B.54) Finally, consider that $`\overline{}(𝐱)=\overline{}^{(3)}(𝐱)`$ when $`s=s_0`$, and $`\overline{}(𝐱)=\overline{}^{(4)}(𝐱)`$ when $`r=r_0`$; and, from condition (1) in the definition of the HHP, $`(𝐱,\tau )`$ is a holomorphic function of $`\tau `$ throughout $`C\overline{I}(𝐱)`$. Therefore, $$\begin{array}{c}((r,s_0),\tau )\text{ and }((r_0,s),\tau )\hfill \\ \text{are continuous functions of }\tau \text{ at }\hfill \\ \tau =s_0\text{ and at }\tau =r_0\text{, respectively. }\hfill \end{array}$$ (6B.55) From the above statements (6B.52) to (6B.55) and from the definition of $`𝒮_{}`$ in Sec. 1, it follows that $`𝒮_{}`$. End of proof. * From condition (3) in the definition of the HHP and from Thm. 33, there exists $`𝐰^{}B(^{(3)})\times B(^{(4)})`$ such that $$𝐯=\mathrm{𝐕𝐰}^{}𝐕_0^1.$$ (6B.56) Therefore, $$𝐕=\mathrm{𝐯𝐕}_0(𝐰^{})^1.$$ (6B.57) However, from the proof of part (i) this theorem \[see Eq. (6B.28)\], there then exists $`𝐰B(^{(3)})\times B(^{(4)})`$ such that $$𝐕(𝐰^{}𝐰)=\mathrm{𝐯𝐕}_0𝐰K^{\mathrm{}}.$$ (6B.58) Therefore, since $`𝐰^{}𝐰B(^{(3)})\times B(^{(4)})`$, it follows from the definition of $`𝒮_𝐕^{\mathrm{}}`$ given by Eq. (6A.28) that $`𝐕𝒮_𝐕^{\mathrm{}}`$. Hence, by definition, $`𝒮_{}^{\mathrm{}}`$, $`H𝒮_H^{\mathrm{}}`$ and $`𝒮_{}^{\mathrm{}}`$. End of proof. ###### COROLLARY 41 ($`k^{\mathrm{}}=K^{\mathrm{}}`$) Suppose that $`\mathrm{}`$ is $`n3`$, $`n+`$ ($`n3`$), $`\mathrm{}`$ or ‘an’. Then $$k^{\mathrm{}}:=\{\mathrm{𝐕𝐰}:𝐕𝒮_𝐕,𝐰B(^{(3)})\times B(^{(4)}),\mathrm{𝐕𝐰}K^{\mathrm{}}\}=K^{\mathrm{}}.$$ (6B.59) Proof: From its definition $$k^{\mathrm{}}K^{\mathrm{}}.$$ (6B.60) Now, suppose $`𝐯K^{\mathrm{}}`$. Since $$𝐕^M=(\delta ^{(3)},\delta ^{(4)}),$$ (6B.61) where $$\delta ^{(i)}(\sigma ):=I\text{ for all }\sigma ^{(i)},$$ (6B.62) we know that $$^M𝒮_{}^{an}𝒮_{}^{\mathrm{}}.$$ (6B.63) Therefore, from the preceding theorem, there exists $`𝒮_{}^{\mathrm{}}`$ such that $``$ is the solution of the HHP corresponding to $`(𝐯,^M)`$; and, if $`𝐕`$ denotes the member of $`𝒮_𝐕^{\mathrm{}}`$ that corresponds to $``$, Eq. (6B.56) in the proof of the preceding theorem informs us that $`𝐰^{}B(^{(3)})\times B(^{(4)})`$ exists such that $`𝐯=\mathrm{𝐕𝐰}^{}`$. So $`𝐯k^{\mathrm{}}`$. We have thus proved that $$K^{\mathrm{}}k^{\mathrm{}},$$ (6B.64) whereupon (6B.60) and (6B.61) furnish us with the conclusion $`k^{\mathrm{}}=K^{\mathrm{}}`$. End of proof. ### C. The generalized Geroch group $`𝒦^{\mathrm{}}`$ Dfn. of $`Z^{(i)}`$ * Let $`Z^{(i)}`$ denote the subgroup of $`K(𝐱_0,^{(i)})`$ that is given by $$Z^{(i)}:=\{\delta ^{(i)},\delta ^{(i)}\},$$ (6C.1) where $`\delta ^{(i)}`$ is defined by Eq. (6B.62). End of Dfn. ###### THEOREM 42 (Center of $`K`$) The center of $`K(𝐱_0,^{(i)})`$ is $`Z^{(i)}`$. Hence the center of $`K`$ is $`Z^{(3)}\times Z^{(4)}`$. Proof: Left for the reader. Hint: See the proof of Lem. 43(i). End of proof. Dfn. of $`[𝐯]`$ for each $`𝐯K^3`$ * For each $`𝐯K^3`$, let $`[𝐯]`$ denote the function such that $$\mathrm{d}om[𝐯]:=𝒮_{}^3$$ (6C.2) and, for each $`_0𝒮_{}^3`$, $$[𝐯](_0):=\text{ the solution of the HHP corresponding to }(𝐯,_0).$$ (6C.3) Note that the existence of $`[𝐯]`$ is guaranteed by Thm. 35 and Thm. 22(iii). End of Dfn. Dfn. of $`𝒦^{\mathrm{}}(𝐱_0,^{(3)},^{(4)})`$ when $`\mathrm{}`$ is $`n3`$, $`n+`$ ($`n3`$), $`\mathrm{}`$ or ‘an’ * Let $$𝒦^{\mathrm{}}:=\{[𝐯]:𝐯K^{\mathrm{}}\}.$$ (6C.4) End of Dfn. The following lemma concerns arbitrary members $`𝐯`$ and $`𝐯^{}`$ of $`K`$, and arbitrary members $`_0`$ and $``$ of $`𝒮_{}`$. Therefore, the lemma could have been given as a theorem in Sec. 1. However, we have saved it for now, because the lemma is directly applicable in the proof of the next theorem. ###### LEMMA 43 (Properties of $`K`$) * Suppose that $`𝐯K`$, $`_0𝒮_{}`$ and $`𝒮_{}`$. Then $``$ is the solution of the HHP corresponding to $`(𝐯,_0)`$ if and only if $`𝐕^1\mathrm{𝐯𝐕}_0B(^{(3)})\times B(^{(4)})`$, where $`𝐕_0`$ and $`𝐕`$ are the members of $`𝒮_𝐕`$ corresponding to $`_0`$ and $``$, respectively. In particular, the solution of the HHP corresponding to $`(𝐯,_0)`$ is $`_0`$ if and only if $`𝐕_0^1\mathrm{𝐯𝐕}_0B(^{(3)})\times B(^{(4)})`$; and the solution of the HHP corresponding to $`(𝐯,^M)`$ is $`^M`$ if and only if $`𝐯B(^{(3)})\times B(^{(4)})`$. * In addition to the premises of part (i) of this lemma, suppose that $`𝐯^{}K`$. Thereupon, if $``$ is the solution of the HHP corresponding to $`(𝐯,_0)`$, and $`^{}`$ is the solution of the HHP corresponding to $`(𝐯^{},)`$, then $`^{}`$ is the solution of the HHP corresponding to $`(𝐯^{}𝐯,_0)`$. If $``$ is the solution of the HHP corresponding to $`(𝐯,_0)`$, then $`_0`$ is the solution of the HHP corresponding to $`(𝐯^1,)`$. Proofs: * This theorem follows from Thm. 33 and the properties of members of $`𝒮_{}`$ that are given in Thm. 1 \[specifically, the properties $`(𝐱,\mathrm{})=I`$ and (iv)\] and Thm. 2. The reader can easily fill in the details of the proof. End of proof. * This follows from the obvious facts that the equations $`Y^{(i)}(𝐱,\sigma )`$ $`:=`$ $`^+(𝐱,\sigma )v^{(i)}(\sigma )[_0^+(𝐱,\sigma )]^1`$ $`=`$ $`^{}(𝐱,\sigma )v^{(i)}(\sigma )[_0^{}(𝐱,\sigma )]^1`$ and $`Y^{(i)}(𝐱,\sigma )`$ $`:=`$ $`^+(𝐱,\sigma )v^{(i)}(\sigma )[^+(𝐱,\sigma )]^1`$ $`=`$ $`^{}(𝐱,\sigma )v^{(i)}(\sigma )[^{}(𝐱,\sigma )]^1`$ imply $`Y^{(i)}(𝐱,\sigma )Y^{(i)}(𝐱,\sigma )`$ $`=`$ $`^+(𝐱,\sigma )v^{(i)}(\sigma )v^{(i)}(\sigma )[_0^+(𝐱,\sigma )]^1`$ $`=`$ $`^{}(𝐱,\sigma )v^{(i)}(\sigma )v^{(i)}(\sigma )[_0^{}(𝐱,\sigma )]^1`$ and $`[Y^{(i)}(𝐱,\sigma )]^1`$ $`=`$ $`_0^+(𝐱,\sigma )[v^{(i)}(\sigma )]^1[F^+(𝐱,\sigma )]^1`$ $`=`$ $`_0^+(𝐱,\sigma )[v^{(i)}(\sigma )]^1[F^+(𝐱,\sigma )]^1`$ for all $`i\{3,4\}`$ and $`\sigma ^{(i)}`$. End of proof. Finally, we prove the following generalized Geroch conjecture: ###### THEOREM 44 * The mapping $`[𝐯]`$ is the identity map on $`𝒮_{}^{\mathrm{}}`$ iff $`𝐯Z^{(3)}\times Z^{(4)}`$. * The set $`𝒦^{\mathrm{}}`$ is a group of permutations of $`𝒮_{}^{\mathrm{}}`$ such that the mapping $`𝐯[𝐯]`$ is a homomorphism of $`K^{\mathrm{}}`$ onto $`𝒦^{\mathrm{}}`$; and the mapping $`\{\mathrm{𝐯𝐰}:𝐰Z^{(3)}\times Z^{(4)}\}[𝐯]`$ is an isomorphism \[viz, the isomorphism of $`K^{\mathrm{}}/(Z^{(3)}\times Z^{(4)})`$ onto $`𝒦^{\mathrm{}}`$\]. * The group $`𝒦^{\mathrm{}}`$ is transitive \[i.e., for each $`_0,𝒮_{}^{\mathrm{}}`$ there exists at least one element of $`𝒦^{\mathrm{}}`$ that transforms $`_0`$ into $``$\]. Proofs: * The statement that $`[𝐯]`$ is the identity mapping on $`𝒮_{}^{\mathrm{}}`$ means that each $`_0𝒮_{}`$ is the solution of the HHP corresponding to $`(𝐯,_0)`$; and, from Lem. 43(i), this is equivalent to the following statement: $$\text{For each }𝐕_0𝒮_𝐕^{\mathrm{}},𝐕_0^1\mathrm{𝐯𝐕}_0B(^{(3)})\times B(^{(4)}).$$ (6C.5) Since $`k^{\mathrm{}}=K^{\mathrm{}}`$ (Cor. 41), each $`𝐯^{}K^{\mathrm{}}`$ is also a member of $`k^{\mathrm{}}`$, and this means that there exist $`𝐕^{}𝒮_𝐕`$ and $`𝐰^{}B(^{(3)})\times B(^{(4)})`$ such that $`𝐯^{}=𝐕^{}𝐰^{}`$. Therefore, from statement (6C.5), $$\begin{array}{c}\text{For each }𝐯^{}K^{\mathrm{}},\text{ there exists }𝐰^{}B(^{(3)})\times B(^{(4)})\hfill \\ \text{such that }𝐰^{}(𝐯^{})^1\mathrm{𝐯𝐕}^{}(𝐰^{})^1B(^{(3)})\times B(^{(4)}).\hfill \end{array}$$ So, since $`B(^{(3)})\times B(^{(4)})`$ is a group, $$\text{For each }𝐯^{}K^{\mathrm{}},(𝐯^{})^1\mathrm{𝐯𝐯}^{}B(^{(3)})\times B(^{(4)}).$$ (6C.6) In particular, since $`𝐕^MK^{\mathrm{}}`$ \[see Eqs. (6B.61) to (6B.63)\] and $`(𝐕^M)^1\mathrm{𝐯𝐕}^M=𝐯`$, $$𝐯B(^{(3)})\times B(^{(4)}).$$ (6C.7) Therefore, there exist $$\alpha _0^{(i)}:^{(i)}R^1\text{ and }\alpha _1^{(i)}:^{(i)}R^1$$ (6C.8) such that $$v^{(i)}=I\alpha _0^{(i)}+\stackrel{~}{J}\alpha _1^{(i)};$$ (6C.9) and, since $`𝐯K^{\mathrm{}}`$ and $`detv^{(i)}=1`$, $$\alpha _0^{(i)}\text{ and }\alpha _1^{(i)}\text{ are }𝐂^{\mathrm{}}$$ (6C.10) and $$(\alpha _0)^2+(\alpha _1)^2=1.$$ (6C.11) Also \[see Eq. (6A.13)\], the function whose domain is $`^{(i)}`$ and whose values are given by $$u^{(i)}(\sigma )=\mathrm{exp}\stackrel{~}{𝒩}(\sigma )$$ is a member of $`K^{an}(^{(i)})K^{\mathrm{}}(^{(i)})`$, where $`K^{\mathrm{}}=K^{\mathrm{}}(^{(3)})\times K^{\mathrm{}}(I^{(4)})`$. Upon letting $`𝐯^{}=(u^{(3)},u^{(4)})`$ in Eq. (6C.6), and upon using Eq. (6C.9), one obtains $$I\alpha _0^{(i)}+J\alpha _1^{(i)}[u^{(i)}]^2B(^{(i)});$$ and this is true if and only if $$\alpha _1^{(i)}(\sigma )\mathrm{sinh}[2\stackrel{~}{𝒩}(\sigma )]=0\text{ for all }\sigma ^{(i)}.$$ (6C.12) However, $`\alpha _1^{(i)}`$ is continuous. Therefore, the condition (6C.12) can hold if and only if $`\alpha _1^{(i)}`$ is identically zero, whereupon (6C.9) and (6C.11) yield $`v^{(i)}=\pm \delta ^{(i)}`$. Hence $`𝐯Z^{(3)}\times Z^{(4)}`$ is a necessary and sufficient condition for $`[𝐯]=`$ the identity map on $`𝒮_{}^{\mathrm{}}`$. End of proof. * Suppose $`𝐕K^{\mathrm{}}`$ and suppose $`𝒮_{}^{\mathrm{}}`$ and the corresponding member of $`𝒮_𝐕^{\mathrm{}}`$ is $`𝐕`$. From the definition of $`𝒮_{}^{\mathrm{}}`$ and Cor. 41, there exists $`𝐰B(^{(3)})\times B(^{(4)})`$ such that $$\mathrm{𝐕𝐰}k^{\mathrm{}}=K^{\mathrm{}}.$$ Therefore, since $`𝐯K^{\mathrm{}}`$ and $`K^{\mathrm{}}`$ is a group, $$𝐯^1\mathrm{𝐕𝐰}K^{\mathrm{}}=k^{\mathrm{}}.$$ Therefore, from the definition of $`k^{\mathrm{}}`$, there exist $`𝐕_0𝒮_𝐕`$ and $`𝐰^{}B(^{(3)})\times B(^{(4)})`$ such that $$𝐯^1\mathrm{𝐕𝐰}=𝐕_0𝐰^{}.$$ So, since $`B(^{(3)})\times B(^{(4)})`$ is a group, $$𝐕^1\mathrm{𝐯𝐕}_0B(^{(3)})\times B(^{(4)}).$$ It then follows from Lem. 43(i) that $``$ is the solution of the HHP corresponding to $`(𝐯,_0)`$, where $`_0`$ is the member of $`𝒮_{}^{\mathrm{}}`$ that corresponds to $`𝐕_0`$. We have thus shown that every member $``$ of $`𝒮_{}^{\mathrm{}}`$ is in the range of $`[𝐯]`$; i.e., $$[𝐯]\text{ is a mapping of }𝒮_{}^{\mathrm{}}\text{ onto }𝒮_{}^{\mathrm{}}.$$ (6C.13) Next, suppose $`_0`$ and $`_0^{}`$ are members of $`𝒮_{}^{\mathrm{}}`$ such that $$:=[𝐯](_0)=[𝐯](_0^{}).$$ Then, $``$ is the solution of the HHP’s corresponding to $`(𝐯,_0)`$ and to $`(𝐯,_0^{})`$, whereupon Lem. 43(ii) informs us that $`_0^{}`$ is the solution of the HHP corresponding to $`(𝐯^1𝐯,_0)`$. Hence, $`_0^{}=_0`$. We have thus shown that $`[𝐯]`$ is one-to-one. Upon combining this result with (6C.13), we obtain $$\begin{array}{c}\text{For each }𝐯K^{\mathrm{}},[𝐯]\text{ is a permutation of }𝒮_{}^{\mathrm{}}\hfill \\ \text{{ i.e., }[𝐯]\text{ is a one-to-one mapping of }𝒮_{}^{\mathrm{}}\text{ onto }𝒮_{}^{\mathrm{}}\text{ }. }\hfill \end{array}$$ (6C.14) Furthermore, the reader can easily show from Lem. 43(ii) that, if $$[𝐯^{}][𝐯]:=\text{ the composition of the mappings }[𝐯^{}]\text{ and }[𝐯],$$ (6C.15) then $$[𝐯^{}][𝐯]=[𝐯^{}𝐯].$$ (6C.16) Lemma 43(ii) also yields $$[𝐯]^1=[𝐯^1].$$ (6C.17) Therefore, since $`K^{\mathrm{}}`$ is a group, $`𝒦^{\mathrm{}}`$ is a group with respect to composition of mappings. The remainder of the proof is straightforward and is left to the reader. End of proof. * Let $`_0`$ and $``$ be any members of $`𝒮_{}^{\mathrm{}}`$ such that the corresponding members of $`𝒮_𝐕^{\mathrm{}}`$ are $`𝐕_0`$ and $`𝐕`$, respectively. By definition of $`𝒮_𝐕^{\mathrm{}}`$, there exist members $`𝐰_0`$ and $`𝐰`$ of the group $`B(^{(i)})\times B(^{(4)})`$ such that $$𝐕_0𝐰_0\text{ and }\mathrm{𝐕𝐰}\text{ are members of }k^{\mathrm{}}=K^{\mathrm{}}.$$ Then, from Lem. 43(i), $``$ is the solution of the HHP corresponding to $`(𝐯,_0)`$, where $$𝐯:=\mathrm{𝐕𝐰}(𝐰_0)^1𝐕_0^1,$$ and is clearly a member of $`K^{\mathrm{}}`$. So, for each $`_0𝒮_{}^{\mathrm{}}`$ and $`𝒮_{}^{\mathrm{}}`$, there exists $`[𝐯]𝒦^{\mathrm{}}`$ such that $`=[𝐯](_0)`$; and that is what is meant by the statement that $`𝒦^{\mathrm{}}`$ is transitive. End of proof. As a final note, the K–C subgroup of $`𝒦^3`$ is $$\{[𝐯]𝒦^{an}:v^{(3)}\text{ and }v^{(4)}\text{ have equal analytic extensions to the domain }]r_1,s_1[\}.$$ ## Acknowledgement Research supported in part by grants PHY-93-07762, PHY-96-01043 and PHY-98-00091 from the National Science Foundation to FJE Enterprises.
warning/0002/math0002140.html	ar5iv	text	# On quadratic and higher normality of small codimension projective varieties ## 1. Introduction A variety $`X𝐏^n`$ is called $`j`$-normal if the restriction map $`H^0(𝐏^n,𝒪(j))H^0(X,𝒪(j))`$ is surjective. Hartshorne’s conjecture implies that smooth varieties $`X𝐏^n`$ of small codimension are $`j`$-normal. Peternell, Le Potier , Schneider and Ein proved indipendently that smooth codimension $`2`$ varieties $`X𝐏^n`$ are $`2`$-normal if $`n10`$. This bound is probably not sharp (Hartshorne’s conjecture implies $`n6`$) but it is interesting because it does not depend on the degree of $`X`$ (for similar bounds depending on the degree, see ).Ein’s results were extended to higher normality by Alzati and Ottaviani in , but the techniques of those papers seem not to work in codimension $`3`$ because the Koszul complexes appearing in the proof have greater length and are difficult to control. On the other hand Ran in proved, with different techniques, that smooth codimension $`2`$ varieties $`X𝐏^n`$ are $`j`$-normal if $`n3j^2+2j+2`$. Ran constructs explicitly, for any $`YH^0(X,𝒪(j))`$, a hypersurface $`F`$ in $`𝐏^n`$ of degree $`j`$ as the union of lines which intersects $`Y`$ with multiplicity $`j+1`$. This works because the assumption implies that the locus of $`j+1`$-secants is not empty. In our doctoral thesis, we expanded all the details of Ran’s paper and we were able to prove the following theorem which gives bounds for $`j`$-normality also in codimension $`r3`$. Denote by $`\mathrm{\Sigma }_{(j+1)}`$ the set of $`(j+1)`$-secants to $`X`$ through a (generic) external point. ###### Theorem 1.1. Let $`X`$ be a $`r`$ codimension subvariety of $`𝐏^{m+r}`$; if $$\mathrm{\Sigma }_{(j+1)}\mathrm{}$$ $$2(r+1)jmr\text{and}(j+1)((r+1)j1)m1$$ then: $$\rho _j:H^0(𝐏^{m+r},𝒪_{𝐏^{m+r}}(j))H^0(X,𝒪_X(j))$$ is surjective. If $`r=2`$, the numerical assumptions of theorem $`1.1`$ are exactly as in , while Ran is able to show that in this bound if $`\mathrm{\Sigma }_{j+1}=\mathrm{}`$ then $`X`$ is a complete intersection. Ran himself pointed out in a remark at the end of the paper that his proof could also be extended to higher codimension. When $`X`$ is the zero locus of a section of a vector bundle, then the numeric assumption is more explicit. ###### Theorem 1.2. Let $`X`$ be a $`m`$ dimension variety in $`𝐏^{m+r}`$ given by the zero locus of a section of a rank $`r`$ vector bundle $`E`$ on $`𝐏^{m+r}`$. We have $$\text{deg }\mathrm{\Sigma }_{j+1}=\frac{1}{(j+1)!}\underset{i=0}{\overset{j}{}}c_r(E(i))$$ ###### Corollary 1.3. With the assumptions of the theorem $`1.2`$, if $$c_r(E(i))0i=1\mathrm{}j$$ $$2(r+1)jmrand(j+1)((r+1)j1)m1$$ then: $$\rho _j:H^0(𝐏^{m+r},𝒪_{𝐏^{m+r}}(j))H^0(X,𝒪_X(j))$$ is surjective. In section $`4`$ we get a new proof of Zak theorem about linear normality with the assumption $`n4r`$. In the same range there is still another proof due to Faltings . Moreover in this paper we prove the following result on quadratical normality where the numeric assumption is easier checked. This is a partial answer to problem $`12`$ in Schneider list . ###### Theorem 1.4. Let $`X`$ be a $`m`$ dimension variety in $`𝐏^{m+r}`$. If $$c_r(N(2))0and6rm4$$ then $`X`$ is $`2`$-normal. I thank G.Ottaviani for the precious suggestions and the useful discussions and L.Göttsche for some ideas used for the proof of theorem $`1.2`$. ## 2. Proof of theorem $`1.1`$ Consider a branched covering that is a finite surjective morphism between two irreducible and nonsingular algebraic varieties $`V`$ and $`W`$ $`f:VW`$; let $`d`$ be the degree of $`f`$. As we are assuming that $`V`$ and $`W`$ are non-singular, $`f`$ is flat and consequently the direct image $`f_{}𝒪_V`$ is locally free of rank $`d`$ on $`W`$. The trace $`Tr_{V/W}:f_{}𝒪_V𝒪_W`$ gives rise to a splitting: $`f_{}𝒪_X=𝒪_WF`$, where $`F=ker(Tr_{V/W}).`$ We shall be concerned with the rank $`d1`$ vector bundle on $`W`$: $`E=F^{}.`$ E will be termed vector bundle associated with the covering $`f`$. Let $`e_f(x)=dim_𝐂(𝒪_xX/f^{}m_{f(V)})`$ be the local degree of $`f`$ in $`x`$ which counts the number of sheets of covering that come together at $`x`$. ###### Theorem 2.1 (Gaffney-Lazarsfeld). Let $`V`$ and $`W`$ be varieties of dimension $`n`$ and $`f:VW`$ a branched covering of degree $`d`$; if the vector bundle associated with a branched covering is ample, then there exists at least one point $`xV`$ at which $$e_f(x)min(d,n+1).$$ Proof See . Lazarsfeld himself points out that smoothness of $`W`$ is not essential. Thanks to this theorem, we are able to prove the following Lemma: ###### Lemma 2.2. Let $`X`$ be a $`r`$-codimensional subvariety of $`𝐏^n`$, if $`rkn`$ and the set of $`k`$-secant lines to $`X`$ through an external point $`P`$ is not empty, then there exists at least a $`k`$-secant through this point at which the $`k`$ points coincide. Proof We consider the projection from $`P`$ of $`k`$-secants on a generic hyperplane $`𝐏^{n1}`$; let $`f`$ be its restriction to the points of $`X`$, and $`Y`$ the image of $`f`$. $`X^{}=f^1(Y)`$, $`X^{}`$ is the set of points in $`X`$ lying on a $`k`$-secant. The dimension $`n^{}`$ of $`X^{}`$ and $`Y`$ is $`n1k(r1)`$ and $`f:X^{}Y`$ is a finite covering with degree $`k`$: by our assumptions the degree of the covering is less than or equal to $`n^{}+1.`$ If we prove that the vector bundle associated with the covering is ample, then we can use the theorem of Gaffney and Lazarsfeld to prove that there exists a point at which the sheets of covering come together. We denote $`C`$ as the cone of $`k`$-secants through an external point $`P`$; since there are $`k`$ points of $`X^{}`$ for each $`k`$-secant, we observe that $`X^{}`$ is a divisor of $`C`$ and since the point $`P`$ is external to $`X`$ this divisor is disjoint from singularities of $`C`$. Let $`C^{}`$ be the desingularization of $`C`$, we have: $$C^{}=𝐏(𝒪_Y𝒪_Y(1)),$$ then $`X^{}`$ is isomorphic to a divisor of $`C^{}`$. $`f_{}𝒪_X^{}`$ is a vector bundle of rank $`k`$; we want to prove that: $$f_{}𝒪_X^{^{}}=𝒪_Y𝒪_Y(1)\mathrm{}𝒪_Y(1k).$$ $`X^{}`$ is the zero locus of a section of $`𝒪_{𝐏(𝒪𝒪(1))}(k)`$; in fact, from we have: $`Pic(C^{})=Pic(Y)H`$, where $`H`$ is hyperplane section. $`X^{}`$ is a divisor which meets the generic fibre in $`k`$ points and it is disjoint to the infinite section, and so $`X`$ is linearly equivalent to $`kH`$. Now we consider the associated exact sequence: $$0𝒪_𝐏(k)𝒪_𝐏𝒪_X^{^{}}0$$ Let $`\pi `$ be the projection from $`𝐏(𝒪𝒪(1))`$ to $`Y`$; applying $`\pi _{}`$ to the sequence we obtain: $$0𝒪_Y\pi _{}𝒪_X^{}R^1\pi _{}𝒪_𝐏(k)0.$$ Using the exercise $`8.4`$ of , (page $`253`$) we prove that: $$R^1\pi _{}𝒪(k)\pi _{}(𝒪(k2))^{}𝒪_Y(1)$$ and from the same exercise we have: $$\pi _{}𝒪(k2)S^{k2}(𝒪𝒪(1))=𝒪𝒪(1)\mathrm{}𝒪(k2)$$ then $$R^1\pi _{}𝒪(k)=𝒪(1)𝒪(2)\mathrm{}𝒪(k+1)$$ substituting in the exact sequence we get: $$\pi _{}𝒪_X^{^{}}=𝒪_Y𝒪_Y(1)𝒪_Y(2)\mathrm{}𝒪_Y(1k)$$ then $$\pi _{}𝒪_X^{^{}}=𝒪_YF$$ $$\pi _{}𝒪_X^{^{}}=f_{}𝒪_X^{^{}}$$ where $`F`$ is a vector bundle whose dual is ample. We can now use the Gaffney-Lazarsfeld’s theorem to obtain the thesis. ###### Lemma 2.3. Let $`G`$ be a generic hypersurface of $`𝐏^n`$ of degree $`j`$ passing through a point $`P`$, then the variety of lines through $`P`$ lying in $`G`$ is a complete intersection of $`𝐏^{n2}`$ with dimension $`nj1`$ and degree $`j!`$. Proof. We can choose a coordinate system such that P is the point $`(a,0,0,\mathrm{},0)`$. Let $`\pi `$ be the hyperplane $`x_0=0`$; for every point $`Q`$ of $`\pi `$ we consider the line $`r`$ through $`P`$ and $`Q`$ that is $`(a(1t),tx_1,\mathrm{},tx_n).`$ G is given by $`F(y_0,\mathrm{},y_n)=0`$ with $`F(y_0,\mathrm{},y_n)=by_0^j+f_1(y_1,\mathrm{},y_n)y_0^{j1}+\mathrm{}+f_j(y_1,\mathrm{},y_n)`$ where $`f_i`$ are polynomials of degree $`i`$; since $`PG`$ we have $`b=0`$. A line $`r`$ lie on G if and only if: $$F(a(1t),tx_1,\mathrm{},tx_n)=ty_0^{j1}f_1(x_1,\mathrm{}x_n)+\mathrm{}t^jf_j(x_1,\mathrm{},x_n)=0$$ for every $`t`$, and so we must have: $`f_i(x_1,\mathrm{},x_n)=0i=1,\mathrm{},n.`$ Since $`G`$ is generic and $`f_1`$ is linear, this gives a transversal intersection contained in $`𝐏^{n2}`$. Finally we get that the variety of lines of $`G`$ through a point $`PG`$ is a complete intersection of degree $`j!`$ and dimension $`n1j`$. Let $`X`$ be a subvariety of $`𝐏^{m+r}`$; we denote by $`\mathrm{\Sigma }_j`$ the cycle of $`j`$-secant lines to $`X`$ through an external point. Proof of theorem 1.1 Consider a generic element $`Y`$ of the linear system $`𝒪_X(j)`$. Since the locus of $`(j+1)`$-secants through a generic point is not empty, then $`X`$ can not be included in a hypersurface of degree $`j`$ and so $`H^0(_X(j))=0`$. In order to prove the theorem we just have to find one hypersurface of degree $`j`$ which contains $`Y`$. We define $`R^k`$=$`\{(y,z)Y\times 𝐏^{m+r}`$ : $``$ a line L from $`z𝐏^{m+r}`$ such that $`LY`$ has multiplicity $`k`$ in $`y`$}. Let $`p`$ and $`q`$ the projections of $`R^k`$ to $`Y`$ and to $`𝐏^{m+r}`$ respectively: $${}_{z}{}^{}R_{}^{k}=p(q^1(z))R_y^k=q(p^1(y))$$ $`R_y^k`$ is the set of points on lines from $`y`$ intersecting Y with multiplicity $`k`$ and it is a cone of vertex $`y`$. In a neighborhood of $`y`$ we can identify $`𝐏^{m+r}`$ with $`𝐂^{m+r}`$ where $`y`$ is the origin, $`Y`$ is defined in an appropriate neighborhood of $`y`$ by $`(r+1)`$ polynomials $`f_1\mathrm{}f_{r+1}`$. $`R_y^k`$ is given by vanishing of the homogeneous components of degree $`k1`$; and so if a generic line $`L`$ of $`𝐏^{m+r}`$ meets $`R_y^k`$ in $`k`$ points, then $`LR_y^k`$. Moreover: $$\text{dim}R_y^km+r(r+1)(k1)yY.$$ Let $`F=q(R^{j+1})`$ be the set of points of $`𝐏^{m+r}`$ on lines which intersect $`Y`$ with multiplicity $`j+1`$ in one point: we want to prove that $`F`$ is the hypersurface we looked for. $`YF`$ because dim $`R_y^{j+1}0`$ $`yY.`$ The first step is to prove that $$F𝐏^{m+r}.$$ Let $`Y^{}=XG`$ where $`G`$ is a generic hypersurface of degree $`j`$. $`Y^{}`$ is obtained by $`Y`$ by semicontinuity, so the dimension of $`F`$ passing from $`Y`$ to $`Y^{}`$ cannot decrease, and since the (j+1)-secants to $`Y^{}`$ are contained in $`G`$ we obtain: $`\text{ dim}Fm+r1.`$ Next step is to prove that : $$\text{dim}Fm+r1.$$ The set of (j+1)-secants to $`Y^{}=XG`$ through an external point $`PG`$ is given by the intersection of $`\mathrm{\Sigma }_{j+1}`$ with the variety of Lemma 2.3, and so, by the assumption, we obtain a variety with degree different to $`0`$; $`j`$ times this degree gives the virtual degree of $`(j+1)`$ secants intersecting a generic line of $`𝐏^{m+r}`$. Since $`Y`$ is a degeneration of $`Y^{}`$, this virtual degree is the same and it is different from $`0`$ as stated previously. Let $`B`$ the locus of $`(j+1)`$-secants to $`Y`$ interecting a generic line, $`B`$ has dimension $`0`$ in the grassnammian of lines in $`𝐏^{m+r}`$ and it is given by $`AS`$ where: $$A=\left\{\text{lines of }𝐏^{m+r}\text{ that are }(j+1)\text{-secant to }Y\right\}$$ $$S=\left\{\text{lines of }𝐏^{m+r}\text{ intersecting a given line}\right\}$$ $$\text{dim}\{AS\}0\text{codim}\{AS\}2(m+r1).$$ Since the line is generic, we have: $$\text{codim}\{AS\}=\text{codim A}+\text{ codim S}$$ $$\text{codim S}=m+r2\text{ codim}Am+r$$ $$\text{dim}Am+r2$$ . Let $`A^{}`$ be the variety of points of $`A`$, then we have: $$\text{ dim}A^{}m+r1.$$ Now we have to prove that $`A^{}=F`$. The inclusion $`FA^{}`$ is trivial; we want to prove that if $`pA^{}`$, then $`pF`$. From Lemma $`2.2`$ we have that if $`p`$ lies on a $`(j+1)`$-secant to $`Y`$ then it lies also on a line intersecting $`Y`$ with multiplicity $`(j+1)`$ in a point of $`Y`$. Finally we have to prove that $$degFj.$$ Let suppose that a generic line L of $`𝐏^{m+r}`$ meets F in (j+1) points $`z_1\mathrm{}z_{j+1}LF`$. Let’s compute $`c_i`$= codim $`(_{z_i}R^{j+1},Y)`$: $$\text{dim }R^{j+1}=\text{dim }Y+\text{dim }p^1(y)=\text{dim}F+\text{dim }q^1(z),$$ since dim $`p^1(y)=`$ dim $`R_y^{j+1}`$ and dim $`q^1(z)=`$dim $`{}_{z}{}^{}R_{}^{j+1}`$, as previously stated, we have: $$\text{dim }R^{j+1}2m1+r(r+1)j$$ then $$\text{dim }_{z_i}R^{j+1}m(r+1)j$$ $$c_i=\text{codim }(_{z_i}R^{j+1},Y)(r+1)j1$$ By the Lefschetz-Barth’s theorem and by the assumption. we have $$𝐂=H^{2c_i}(𝐏^{m+r},𝐂)=H^{2c_i}(Y,𝐂)\underset{i=1}{\overset{j+1}{}}{}_{z_i}{}^{}R_{}^{j+1}\mathrm{}$$ in fact: $$2c_imr2e(j+1)((r+1)j1)m1.$$ Let $`y_{i=1}^{j+1}{}_{z_i}{}^{}R_{}^{j+1}`$ then $`z_iLR_y^{j+1}peri=1\mathrm{}j+1`$ and so $`LR_y^{j+1}.`$ This is a contradiction as $`L`$ is generic. We deduce that $`degFj`$. ## 3. Proof of theorem $`1.2`$ Proof of theorem 1.2 Let $`P`$ be the fixed point and $`QG(𝐏^1,𝐏^n)`$ the space of lines from $`P`$, $`Q𝐏^{n1}`$; let: $$T=\{(q,l)q𝐏^nlQql\}$$ and $`\alpha `$ and $`\beta `$ be the projections of $`T`$ on $`𝐏^n`$ and $`Q`$ respectively. $`T`$ is a $`𝐏^1`$-bundle on $`Q`$ and the fibre is given by all the points lying on lines $`l`$, we can view it as the projectivised of $`𝒪_Q𝒪_Q(1)`$. Let $`(T/Q)^{k+1}`$ be the $`(k+1)`$-power of fibre of $`T`$ on $`Q`$,that is: $$(T/Q)^{k+1}=\underset{k+1\text{t}imes}{\underset{}{T\times _QT\times _Q\mathrm{}\times _QT}}$$ We call $`Z(T/Q)^{k+2}`$ the incidence variety in $`T\times _Q(T/Q)^{k+1}`$, that is: $$Z=\{(x_0,\mathrm{},x_{k+1}(T/Q)^{k+2}x_0=x_i\text{ for same }i(1,\mathrm{},k+1)\}$$ Let $`p`$ and $`q`$ be the projections of $`Z`$ on $`T`$ and $`(T/Q)^{k+1}`$ respectively; we denote: $$E^{(k+1)}=q_{}(p^{}\alpha ^{}(E))$$ $`E^{(k+1)}`$ is a vector bundle on $`(T/Q)^{k+1}`$ of rank $`r(k+1)`$. Let $`s`$ be the section of $`E`$ such that $`X`$ is the zero locus of $`s`$; $`s^{(k+1)}=q_{}(p^{}\alpha ^{}(s))`$ is a section of $`E^{(k+1)}`$ which vanishes in the set: $`\{(x_1,\mathrm{},x_{k+1}(T/Q)^{k+1}\alpha (x_i)X\}`$. The line through $`\alpha (x_1),\mathrm{},\alpha (x_{k+1})`$ is a $`(k+1)`$-secant to $`X`$. Considering that the rearrangement of those points gives the same $`(k+1)`$-secant to $`X`$, from Portous’ formula we have that, if the dimension is zero, the number of $`(k+1)`$-secant is given by the degree of the top Chern-class $`c_{(k+1)r}(E^{(k+1)})`$ divided by $`(k+1)!`$. So if we want to know the degree of $`(k+1)`$-secants we have to compute $`c_{(k+1)}(E^{(k+1)})`$. Let $`q_1,\mathrm{}q_{k+1}`$ be the projections of $`(T/Q)^{k+1}`$ on $`𝐏^n`$; we have the following exact sequence: $$0q^{}E𝒪(\mathrm{\Delta }_{1,k+1}\mathrm{}\mathrm{\Delta }_{k,k+1})E^{(k+1)}E^{(k)}0$$ with $`\mathrm{\Delta }_{i,j}=\{(x_1,\mathrm{}x_{k+1}(T/Q)^{k+1}x_i=x_j\}`$. From sequence we have: $$c_{(k+1)r}E^{(k+1)}=c_{kr}E^{(k)}c_r(q^{}E𝒪(\mathrm{\Delta }_{1,k+1}\mathrm{}\mathrm{\Delta }_{k,k+1})).$$ It is necessary to determine the cohomology of $`T`$ and $`(T/Q)^{k+1}`$. Let $`\alpha `$ and $`\beta `$ be the projections of $`T`$ on $`𝐏^n`$ and $`Q`$ respectively: $`T`$ is blow-up of $`𝐏^n`$ in $`P`$; we call $`D`$ the exceptional divisor and $`H=\alpha ^{}(𝒪_{𝐏^n}(1))`$, then we have: $`HD=\beta ^{}(𝒪_Q(1)),D=𝒪_T(1).`$ The Wu-Chern’s equation gives: $`D^2+\beta ^{}𝒪_Q(1)D=0.`$ The intersection ring of $`T`$ is generated by two elements: $$D,HD=D,\beta ^{}𝒪_Q(1).$$ For the next degrees we have: $$(\beta ^{}𝒪_Q(1))^2D=\beta ^{}𝒪_Q(1)D^2=D^3$$ $$\mathrm{}$$ $$(\beta ^{}𝒪_Q(1))^nD=(1)^{n1}\beta ^{}𝒪_Q(1)D^n=D^{n+1}.$$ We observe that $`H^n=1`$ and $`D^n=(1)^{n1}`$ in fact $`D_D=𝒪_D(1)`$ and $`D𝐏^{n1}`$. Consider now the fibred product $`T\times _QT`$: $`H^{}(T)`$ is generated by $`D`$ as $`H^{}(Q)`$-module; $`H^{}(T\times _QT)=H^{}(T)\times _{H^{}(Q)}H^{}(T)`$ is generated by $`D1=D_1`$ and $`1D=D_2`$ as $`H^{}(Q)`$-module; if we consider it as a vector space we have: $$H^2(T\times _QT)=D_1,D_2,\beta ^{}𝒪_Q(1)$$ $$H^4(T\times _QT)=D_1^2,D_2^2,(\beta ^{}𝒪_Q(1))^2,D_1D_2$$ $$\mathrm{}$$ $$H^{2j}(T\times _QT)=D_1^j,D_2^j,(\beta ^{}𝒪_Q(1))^j,D_1D_2\beta ^{}𝒪_Q(1)^{j2}$$ we denote: $`H_1=q_1^{}(H)H_2=q_2^{}(H)`$ $$H_1=D_1+\beta ^{}𝒪_Q(1)H_2=D_2+\beta ^{}𝒪_Q(1).$$ We prove that: $`\mathrm{\Delta }_{1,2}=D_1+D_2+\beta ^{}𝒪_Q(1)`$ . Let $`p_1`$ and $`p_2`$ be the projection of $`T\times _QT`$ on the two factors. $`\mathrm{\Delta }_{1,2}`$ is given by the zero locus of a section of $`p_1^{}𝒪(1)p^{}Q_{rel}`$ (see ,page 242). We know that $`c_1(p_1^{}𝒪(1))=D_1`$; consider now the exact section $$0𝒪(1)\beta ^{}𝒪\beta ^{}𝒪(1)Q_{rel}0.$$ We get $`c_1(p_2^{}Q_{rel})=D_2+\beta ^{}𝒪(1)`$ and so we have $`\mathrm{\Delta }_{1,2}=D_1+D_2+\beta ^{}𝒪_Q(1).`$ For the general case we have that: $`H^2(T/Q)^{k+1}`$ is generated by $`D_1,D_2,\mathrm{}D_{k+1},\beta ^{}𝒪_Q(1)`$, $`H^{2m}(T/Q^{k+1})`$ is generated by $`(\beta ^{}𝒪_Q(1))^m,D_{i_1}\mathrm{}D_{i_t}\beta ^{}𝒪_Q(1)^{mt}.`$ Moreover: $`\mathrm{\Delta }_{i,j}=D_i+D_j+\beta ^{}𝒪_Q(1).`$ Now we prove the theorem proceeding by induction on $`k`$: for $`k=1`$ the exact sequence is: $$0q_2^{}(E)O(\mathrm{\Delta }_{1,2})E^{(2)}q_1^{}(E)0$$ $`c_{2r}(E)^{(2)}`$ $`=`$ $`c_r((q_1(E))c_r(q_2^{}(E)O(D_1D_2\beta ^{}𝒪_Q(1))`$ $`=`$ $`c_r(E)H_1^r[c_r(E)H_2^r+c_{r1}(E)H_2^{r1}(D_1D_2\beta ^{}𝒪_Q(1))+`$ $`\mathrm{}+c_{ri}(E)H_2^{ri}(D_1D_2\beta ^{}𝒪_Q(1))^i+`$ $`\mathrm{}+(D_1D_2\beta ^{}𝒪_Q(1))^r]`$ since: $`H_1D_1=0`$ and $`D_2+\beta ^{}𝒪_Q(1)=H_2`$ we have: $`c_{2r}=c_r(E)c_r(1)H_1^rH_2^r`$. Now we suppose the statement true for $`nk`$ and we try to prove it for $`n=k+1`$. $`c_{(k+1)r}`$ $`=`$ $`c_{kr}E^{(k)}c_r(q_{k+1}^{}(E)O(D_1D_2\mathrm{}D_kkD_{k+1}k\beta ^{}𝒪_Q(1))`$ $`=`$ $`c_r(E)c_r(E(1))\mathrm{}c_r(E(k+1))H_1^r\mathrm{}H_{k1}^r[c_r(E)H_{k+1}^r+`$ $`c_{r1}(E)H_{k+1}^{r1}(D_1D_2\mathrm{})+\mathrm{}(D_1\mathrm{}k\beta ^{}𝒪_Q(1))^r]`$ since we know that: $`H_iD_i=0`$ and $`D_{k+1}+\beta ^{}𝒪_Q(1)=H_{k+1}`$ we obtain: $$c_{(k+1)r}(E^{(k+1)})=c_r(E)c_r(E(1))\mathrm{}c_r(E(k))H_1^r\mathrm{}H_{k+1}^r$$ and so the theorem is proved. Remark 1 We observe that if the dimension of the locus of $`k`$-secants through a generic point is smaller than expected, then the class of the formula has to be zero (see Remark $`2.2`$). Remark 2 In the case $`r=2`$ the theorem has been already proved by Ran in \[R\]. By the Hartshorne-Serre correspondence every subcanonical subvariety of codimension $`2`$ is a zero locus of a section of a rank $`2`$ vector bundlon $`𝐏^n`$; moreover if $`n10`$ by Larsen’s theorem we have that every subvariety is subcanonical. In this case the formula for $`j+1`$-secant is true for every subvariety. ## 4. A new proof of Zak theorem on linear normality Let $`X`$ be a $`r`$ codimensional subvariety of $`𝐏^n`$; from Barth theorem we have that if $`rn/4`$ then $`H^{2i}(X,)`$; in particular we can write $`c_i(N)=c_iH^i`$ with $`i=1\mathrm{}r`$ where $`c_i𝐙`$. From now, we consider $`c_i(N)`$ as a integer. ###### Lemma 4.1. Let $`X`$ be a $`r`$ codimensional subvariety of $`𝐏^n`$. If $`n4r`$, then the degree of set of bisecant to $`X`$ through an external point is $$c_r(N)c_r(N(1)).$$ Proof. Let $`P`$ be the fixed point, if we project $`X`$ from $`P`$ to a generic hyperplane we can use the double point formula to get the set of bisecant to $`X`$ from $`P`$. $$2\mathrm{\Sigma }_2=f^{}f_{}[X](c(f^{}T𝐏^{n1})c(TX)^1)_{r1}[X]$$ and from the exact sequence $$0T_XT𝐏_S^nN_{X,𝐏^n}0$$ we have $`c(TX)^1=c(T𝐏^n)^1c(N)`$ and substituting we get $$2\mathrm{\Sigma }_2=H^{r1}(dc_{r1}+c_{r2}+\mathrm{}(1)^ic_i\mathrm{})=c_r(N(1))H^{r1}.$$ From theorem $`1.1`$ and from Lemma we get a different proof of Zak’s theorem. ###### Theorem 4.2 (Zak). Let $`X`$ be a $`r`$ codimensional subvariety of $`𝐏^n`$, if $`n4r`$, then $`X`$ is linearly normal. Proof. We prove the theorem proceeding by induction on $`r`$. If $`r=1`$ is trivially true. Now we suppose that it is true for $`r1`$. If $`c_r(N(1))0`$ from theorem $`1.1`$ and from lemma $`4.1`$ we have the thesis. If $`c_r(N(1))=0`$ from lemma $`4.1`$ we have that there are not bisecant to $`X`$ through an external point $`P`$; if we project $`X`$ from $`P`$ to a generic hyperplane, we get a smooth subvariety in $`𝐏^{n1}`$ of codimension $`r1`$ that is linearly normal by induction. This is a contradiction. ## 5. Proof of theorem $`1.4`$ ###### Lemma 5.1. Let $`l,m,p𝐍`$ such that $`l+p=m`$ then we have $$\left(\genfrac{}{}{0pt}{}{l}{t}\right)=\underset{i=0}{\overset{k}{}}(1)^i\left(\genfrac{}{}{0pt}{}{m}{ti}\right)\left(\genfrac{}{}{0pt}{}{p1+i}{i}\right)$$ $$\underset{1=0}{\overset{t}{}}(1)^i\left(\genfrac{}{}{0pt}{}{n}{ti}\right)\left(\genfrac{}{}{0pt}{}{n+1+i}{i}\right)=(1)^t$$ Proof. We consider an exact sequence $$0ABC0$$ where A, B, C are vector spaces with dimension respectively $`l,m,p`$. From this exact sequence we obtain other two exact sequences: (1) $$0^tA^tB^{t1}BC\mathrm{}^{ti}S^iC\mathrm{}S^nC0$$ (2) $$0^tA^{t1}AB\mathrm{}^{ti}AS^iB\mathrm{}S^nBS^nC0$$ considering that $`^t(𝐂^m)=\left(\genfrac{}{}{0pt}{}{m}{t}\right)`$ and $`S^t(𝐂^𝐧)=\left(\genfrac{}{}{0pt}{}{n1+t}{t}\right)`$ we have $$\left(\genfrac{}{}{0pt}{}{l}{t}\right)=\underset{1=0}{\overset{t}{}}(1)^i\left(\genfrac{}{}{0pt}{}{m}{ti}\right)\left(\genfrac{}{}{0pt}{}{p1+i}{i}\right)$$ From $`(2)`$ if $`m=n+1`$ and $`l=n`$ we have: $$\underset{1=0}{\overset{t}{}}(1)^i\left(\genfrac{}{}{0pt}{}{n}{ti}\right)\left(\genfrac{}{}{0pt}{}{n+1+i}{i}\right)=(1)^t$$ ###### Lemma 5.2. Let $`X`$ a $`r`$ codimensional subvariety of $`𝐏^n`$ then if $`n4r`$ the locus of trisecant is $$\mathrm{\Sigma }_3=\frac{1}{2}H^{2r2}c_r(N(1))c_r(N(2))$$ Proof Göttsche’s formula for trisecant through a fixed point is $$\mathrm{\Sigma }_3=(a)+(b)(c)$$ where: $$(a)=H^{2r2}\left(\frac{n}{2}d^2\underset{k=0}{\overset{nr}{}}\left(\left(\genfrac{}{}{0pt}{}{2n2r+2}{k}\right)\left(\genfrac{}{}{0pt}{}{n}{kn+2r2}\right)\right)_XH^ks_{nrk}/2\right)$$ $$(b)=\underset{k=0}{\overset{2r2}{}}\underset{t=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{t}\right)\left(\genfrac{}{}{0pt}{}{n+1}{kt}\right)\underset{j=rt1}{\overset{2r2k}{}}2^{j+tr+1}s_j(X)s_{2r2kj}(X)H^k$$ and $$(c)=\underset{k=0}{\overset{2r2}{}}d\left(\genfrac{}{}{0pt}{}{n+r}{k}\right)s_{2r2k}(X)H^k.$$ We prove the Lemma when $`r`$ is even (the case $`r`$ odd is the same). It is well known: $$s_k=\underset{i=0}{\overset{n}{}}(1)^{k+1}H^{ki}c_i(N)\left(\genfrac{}{}{0pt}{}{n+ki}{ki}\right).$$ Let $`c_i=c_i(N)`$; substituting we have: $`(a)`$ $`=`$ $`H^{2r2}({\displaystyle \frac{n}{2}}d^2{\displaystyle \frac{1}{2}}({\displaystyle \underset{k=0}{\overset{nr}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n2r+2}{k}}\right)`$ $`{\displaystyle \underset{i=0}{\overset{r}{}}}(1)^{nrk+i}c_iH^{nri}\left({\displaystyle \genfrac{}{}{0pt}{}{n+nrki}{nrki}}\right)+`$ $`{\displaystyle \underset{k=0}{\overset{nr}{}}}{\displaystyle \underset{i=0}{\overset{r}{}}}(1)^{nrk+i}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{kn+2r2}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n+nrki}{nrki}}\right)))`$ we put $`k^{}=nrki`$ and so we have: $`(a)`$ $`=`$ $`H^{2r2}({\displaystyle \frac{n}{2}}d^2{\displaystyle \frac{1}{2}}({\displaystyle \underset{i=0}{\overset{r}{}}}c_iH^{nri}{\displaystyle \underset{k^{}=0}{\overset{nri}{}}}(1)^k^{}\left({\displaystyle \genfrac{}{}{0pt}{}{n+k^{}}{k^{}}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n2r+2}{nrik^{}}}\right)+`$ $`{\displaystyle \underset{k=0}{\overset{r2i}{}}}(1)^k^{}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{r2ik^{}}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n+k^{}}{k^{}}}\right)))`$ now we can use the Lemma $`5.1`$ and we obtain (3) $$(a)=H^{2r2}\left(\frac{n}{2}d^2\frac{1}{2}\left(d^2(n2r+1)\underset{i=0}{\overset{r1}{}}(1)^{ri}c_iH^{nri}\right)\right)$$ $`(b)`$ $`=`$ $`{\displaystyle \underset{t=0}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{t}}\right){\displaystyle \underset{j=rt1}{\overset{2r2}{}}}2^{j+t+1r}s_j{\displaystyle \underset{k=t}{\overset{2r2j}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{kt}}\right){\displaystyle \underset{m=0}{\overset{r}{}}}(1)^{2r2kj+m}`$ $`H^{2r2jm}c_m\left({\displaystyle \genfrac{}{}{0pt}{}{n+2r2kjm}{2r2kjm}}\right)`$ If we denote $`k^{}=2r+2kjm`$ we get $`(b)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{r}{}}}c_m{\displaystyle \underset{t=0}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{t}}\right){\displaystyle \underset{j=rt1}{\overset{2r2}{}}}2^{j+t+1r}s_jH^{2r2jm}`$ $`{\displaystyle \underset{k^{}=0}{\overset{2r2jtm}{}}}(1)^k^{}\left({\displaystyle \genfrac{}{}{0pt}{}{n+1}{2r2jmtk^{}}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n+k^{}}{k^{}}}\right)`$ From Lemma $`5.1`$ we have that the last sum is equal to $`1`$ if $`2r2jmt=0`$ and equal to $`0`$ in the other cases; this fact implies also that $`j=2r2mtrt1`$ and so we obtain $`mr1`$. $$(b)=\underset{m=0}{\overset{r1}{}}c_m\underset{t=0}{\overset{n1}{}}\left(\genfrac{}{}{0pt}{}{n}{t}\right)2^{r1m}s_{2r2mt}H^t$$ $`(b)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{r1}{}}}c_m2^{r1m}{\displaystyle \underset{t=0}{\overset{n1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n}{t}}\right){\displaystyle \underset{i=0}{\overset{r}{}}}(1)^{2r2mt+i}`$ $`H^{2r2im}c_i\left({\displaystyle \genfrac{}{}{0pt}{}{n+2r2mti}{2r2mti}}\right)`$ let $`t^{}=2rsmti`$ $`(b)`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{r1}{}}}{\displaystyle \underset{i=0}{\overset{2r2m}{}}}c_mc_iH^{2r2im}2^{r1m}`$ $`{\displaystyle \underset{t^{}=0}{\overset{2r2mi}{}}}(1)^t^{}\left({\displaystyle \genfrac{}{}{0pt}{}{n+t^{}}{t^{}}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2r2mti}}\right)`$ and again from the Lemma $`5.1`$ we have (4) $$(b)=\underset{m=0}{\overset{r1}{}}\underset{i=0}{\overset{2r2m}{}}(1)^{m+i}2^{r1m}c_mc_iH^{2r2im}$$ $$(c)=\underset{k=0}{\overset{2r2}{}}d\left(\genfrac{}{}{0pt}{}{n+r}{k}\right)\underset{i=0}{\overset{r}{}}(1)^{2r2k+i}H^{2r2i}c_i\left(\genfrac{}{}{0pt}{}{n+2r2ki}{2r2ki}\right)$$ let $`k^{}=2r2ki`$ $`(c)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{r}{}}}dH^{2r2i}c_i{\displaystyle \underset{k^{}=0}{\overset{2r2}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{n+r}{2r2ik^{}}}\right)(1)^k^{}\left({\displaystyle \genfrac{}{}{0pt}{}{n+k^{}}{k^{}}}\right)`$ $`=`$ $`{\displaystyle \underset{i=0}{\overset{r}{}}}dH^{2r2i}c_i\left({\displaystyle \genfrac{}{}{0pt}{}{r1}{2r2i}}\right)`$ (5) $$(c)=dc_{r1}H^{r1}+d^2(r1)H^{2r2}.$$ Supposing that we are in the range of Barth’s theorem, we have that $`c_i=c_iH^i`$ where $`c_i`$. Finally we get from (3), (4) and (5): $$\mathrm{\Sigma }_3=(a)+(b)(c)=\frac{1}{2}H^{2r2}\underset{m=0}{\overset{r}{}}\underset{i=0}{\overset{r}{}}(1)^{m+i}2^{rm}c_mc_i$$ that is $$\mathrm{\Sigma }_3=H^{2r2}\frac{1}{2}c_r(N(1))c_r(N(2))$$ .
warning/0002/nucl-th0002051.html	ar5iv	text	# Transition from hadronic to partonic interactions for a composite spin-1/2 model of a nucleon ## I Introduction At low energies and momentum transfers, nuclei are described in term of nucleons . Interactions between the nucleons are modelled successfully by exchange of mesons , or more simply, by potentials. When nuclei are probed at very high momentum transfer, e.g., in electron scattering, partons within the nucleons and mesons become the dominant scatterers . Interactions between the partons are described by QCD. Between the high and low momentum transfer regimes, there is a transition region where a good description is lacking. The meson-exchange dynamics does not account in a satisfactory way for the compositeness of the nucleons and mesons. Therefore, it is of interest to study quark-based composite models of hadrons in order to get some insight on the limits of validity of a hadronic description. Electron scattering data for momentum transfer Q $``$ 1 GeV/c often meet dual descriptions: models based on hadrons on one hand and models based on quark phenomenology on the other . Moreover, the two kinds of description generally are not reconciled to one another in the sense that there is no smooth transition from one to the other as Q increases. Perturbative QCD descriptions are mainly qualitative and not properly normalized at low energy . In the mesonic description, the mechanism of hard scattering from quarks that predominates in the perturbative QCD description is hidden or absent. In this paper, we develop a simple model of a nucleon as a bound state of a fermion and a boson with the goal of gaining some insight into the transition region where, as Q increases, one passes from the dominance of hadronic processes to the dominance of scattering from the constituents of a nucleon. One may think of this model as having a quark and a spin-0 diquark bound together to make a nucleon and its excited states. The model is covariant and gauge invariant, but it lacks confinement. Excited states of the nucleon are a continuum of quark and diquark scattering states. Thus, it is mainly useful for processes where nucleon resonances do not play an important role. One such case is mesonic-exchange currents in nuclei. An essential feature arises from compositeness: there are contact-like terms in second-order interactions. These are required by gauge invariance and they play a small but significant role at low energy, for example, in low-energy theorems . For very large momentum transfer, the contact-like terms become dominant. They contain the leading-order mechanism for the external probe to scatter from the partons without any intervening hadronic state. When a hadronic state exists between interactions, it produces form factors that fall rapidly with increasing Q, thus quenching the scattering. This is the fate of hadron-like terms in the second-order interactions, i.e., the terms that provide a hadronic interpretation at low momentum transfer. In the limit that one of the interactions transfers a large momentum, the contact-like terms tend to the off-forward parton distributions for the composite nucleon model . For the simple model that we consider, there is a clean separation of the hadron-like and contact-like contributions to second-order interactions. Interactions of the model nucleon with an electromagnetic probe have some realistic features. By introducing cutoff parameters, the nucleon’s charge and magnetic form factors can be described reasonably. At low momentum transfers, interactions of the model nucleon can be interpreted in terms of hadron dynamics. For asymptotically large momentum transfer Q, at fixed $`x=Q^2/(2M\nu )`$, scaling obtains. We calculate the resulting parton distribution $`f(x)`$. In Sec. 2 we formulate the model in terms of a lagrangian for a fermion and a boson interacting via a contact interaction. The model is not renormalizable: it is regulated by introducing subtraction terms of the Pauli-Villars type. We consider only the simplest subset of contributions to the fermion-boson correlator. This produces a spin-1/2 propagator with a single bound state pole (“the nucleon”) at mass M. Electromagnetic and pionic interactions are introduced in Sec. 3 as couplings to the fermion constituent (“the quark”). For simplicity, couplings to the boson (“the diquark”) are omitted. For pseudoscalar coupling of the pion to the quark, the model produces mostly pseudovector coupling to the nucleon. It would be purely pseudovector if the mass of the quark were zero and there were no regulators of fermion type. In Sec. 4 we consider a virtual photopion amplitude involving second-order interactions with the composite nucleon. Two types of interaction occur: first, interactions with intervening propagation of a nucleon or its excited states and second, contact-like contributions where photon and pion interactions with the nucleon occur within the same vertex. A standard analysis based upon elementary particles with form factors is compared with the composite nucleon analysis. In Sec. 5 we consider deep inelastic scattering from the nucleon for finite Q and as Q $`\mathrm{}`$. In Sec. 6 we present calculations of the virtual photo-pion production amplitude for a kinematical situation that arises in meson-exchange contributions to electron-deuteron scattering. Calculations show that contact-like terms can become dominant for Q $``$ 1 GeV/c for some processes. Conclusions are presented in Sec. 7. A more complete description of the details of the calculations is given in four appendices. ## II Composite nucleon model A fermion and a boson interacting via a contact interaction can generate a composite spin-1/2 particle. For this purpose, the following lagrangian is used . $`L=\overline{\psi }(x)(i/m)\psi (x)+{\displaystyle \frac{1}{2}}[_\nu \varphi (x)^\nu \varphi (x)\mu ^2\varphi ^2(x)]+g\overline{\psi }(x)\psi (x)\varphi ^2(x),`$ (1) where $`\psi (x)`$ is the field for a fermion of mass m and $`\varphi (x)`$ is the field for a boson of mass $`\mu `$. The fermion-boson contact interaction with coupling constant g is not renormalizable; finite results are obtained by introducing a Pauli-Villars regulator of mass $`\mathrm{\Lambda }_1`$. A bound state appears as a pole in the fermion-boson correlator, $$G(p)=id^4xe^{ipx}0\|T\left(\psi (x)\varphi (x)\overline{\psi }(0)\varphi (0)\right)\|0.$$ (2) Figure 1 shows the sequence of elementary bubble graphs that contribute to G(p) in a perturbative expansion. Because this sequence is sufficient to exhibit a bound state, contributions beyond those shown in Fig. 1 are not considered. Summing the bubble graphs of Fig. 1 produces $$G(p)=\frac{1}{1\mathrm{\Sigma }(p)}.$$ (3) Here, $`\mathrm{\Sigma }(p)`$ is the contribution of a single fermion-boson loop, $`\mathrm{\Sigma }_b(p;m,\mu ,\mathrm{\Lambda }_1)=ig{\displaystyle \frac{d^4k}{(2\pi )^4}S(pk;m)D(k;\mu ,\mathrm{\Lambda }_1)},`$ (4) where the propagator for the fermion is $`S(p;m)=1/(p/m+i\eta )`$. With a Pauli-Villars regulator of mass $`\mathrm{\Lambda }_1`$ included, the propagator for the boson line is $$D(k;\mu ,\mathrm{\Lambda }_1)=\frac{1}{k^2\mu ^2+i\eta }\frac{1}{k^2\mathrm{\Lambda }_1^2+i\eta }.$$ (5) A generalization of the model that is suitable for describing a nucleon’s form factor is obtained by including additional regulator terms as follows, $`\mathrm{\Sigma }(p)=ig{\displaystyle \frac{d^4k}{(2\pi )^4}\left[S(pk;m)\alpha S(pk;m_1)(1\alpha )S(pk;m_2)\right]}`$ (6) $`\times \left[D(k;\mu ,\mathrm{\Lambda }_1)+\beta D(k;\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)\right].`$ (7) (8) where $$\alpha =\frac{m_2m}{m_2m_1},$$ (9) and $$\beta =\frac{\mu ^2\mathrm{\Lambda }_1^2}{\mathrm{\Lambda }_2^2\mathrm{\Lambda }_1^2}.$$ (10) The constants $`\alpha `$ and $`\beta `$ are selected so that high loop momentum is cut off as k<sup>-9</sup>. It is evident that the generalized form for $`\mathrm{\Sigma }(p)`$ is equal to a linear combination of the elementary bubble graph terms, $`\mathrm{\Sigma }_b`$, defined above. $`\mathrm{\Sigma }(p)=\mathrm{\Sigma }_b(p;m,\mu ,\mathrm{\Lambda }_1)+\beta \mathrm{\Sigma }_b(p;m,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)`$ (11) $`\alpha \left[\mathrm{\Sigma }_b(p;m_1,\mu ,\mathrm{\Lambda }_1)+\beta \mathrm{\Sigma }_b(p;m_1,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)\right]`$ (12) $`(1\alpha )\left[\mathrm{\Sigma }_b(p;m_2,\mu ,\mathrm{\Lambda }_1)+\beta \mathrm{\Sigma }_b(p;m_2,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)\right]`$ (13) When a bound state of mass M is present, the pole in the composite system propagator G(p) has the form $$G(p)=\frac{Z_2}{p/M+i\eta }+R(p),$$ (14) where Z<sub>2</sub> is a wave-function renormalization factor. A renormalized propagator $`\stackrel{~}{\mathrm{G}}`$(p) is obtained by dividing G(p) by Z<sub>2</sub> such that there is unit residue for the nucleon pole. The remainder R(p) is regular at $`p/=M`$ and it represents excited state contributions. In the model considered, the excited state spectrum is a continuum of quark-diquark scattering states. This is an unrealistic feature for a nucleon so the model should be used where the effects of resonances are not important. The most general form for $`\mathrm{\Sigma }`$ that is allowed by Lorentz invariance is $$\mathrm{\Sigma }(p)=A(p^2)p/+B(p^2).$$ (15) Presence of the bound state pole means that $`\mathrm{\Sigma }(p)=1`$ at $`p/=M`$. This condition leads to $`Z_2^1`$ $`=`$ $`\left({\displaystyle \frac{d\mathrm{\Sigma }(p)}{dp/}}\right)_{p/=M}`$ (16) $`=`$ $`[A_0+2M(A_0^{}M+B_0^{})],`$ (17) where $`A_0A(M^2)`$, $`A_0^{}dA(p^2)/dp^2\|_{p^2=M^2}`$, and similarly for $`B_0^{}`$. For later use, we introduce covariant projection operators, $$L^\rho (p)\frac{W_p+\rho p/}{2W_p},$$ (18) where $`\rho `$ = $`+`$ or $``$, W<sub>p</sub> = $`\sqrt{\mathrm{p}^2}`$, and L<sup>+</sup> (p) + L<sup>-</sup>(p) = 1. Projecting the propagator to the $`\rho `$ = $`+`$ and $``$ subspaces in which $`p/`$ takes the values $`\pm `$W<sub>p</sub>, leads, for the renormalized propagator, to $$\stackrel{~}{G}(p)=G^+(p)+G^{}(p)$$ (19) where $$G^\rho (p)=\frac{L^\rho (p)}{Z_2[1B(p^2)\rho W_pA(p^2)]}.$$ (20) To summarize this section, the composite model of a nucleon is formulated in a covariant way as a bound state of a spin-1/2 quark and a spin-0 diquark. Details of the calculation of A(p<sup>2</sup>), B(p<sup>2</sup>) and Z<sub>2</sub> are given in Appendix A. ## III Photon and pion interactions Electromagnetic interactions are introduced via a fermion-photon coupling term in the lagrangian: $`_\mu `$ $`\overline{\psi }`$(x) $`\gamma ^\mu \widehat{e}`$A<sub>μ</sub>(x) $`\psi `$(x), where $`\widehat{e}=\frac{1}{2}e(1+\tau _3)`$ is the charge operator for the quark. Photon coupling to the boson is omitted in order to keep the model simple. Consequently, the model proton is composed of a quark of charge e and a neutral diquark. The model neutron is composed of a neutral quark and diquark and thus has no electromagnetic interations. Inserting a photon into the propagator as indicated in Figure 2 produces the form $$G(p_f)\widehat{e}\mathrm{\Lambda }_\mu (p_f,p_i)G(p_i),$$ (21) where $`\mathrm{\Lambda }_\mu `$ describes the photon-nucleon vertex. One extracts the photon-nucleon (dressed) interaction as the residue of the two poles at $`p/_i`$ = M and $`p/_f`$ = M, which leads to $$\overline{u}(p_f;M)\widehat{e}Z_2\mathrm{\Lambda }_\mu (p_f,p_i)u(p_i;M),$$ (22) where u(p;M) is the Dirac spinor for mass M and momenta p<sub>i</sub> and p<sub>f</sub> are on the mass shell. The Z<sub>2</sub> factor and Dirac spinor factors arise from the parts of the initial- and final-state propagators that attach to the vertex $`\mathrm{\Lambda }_\mu `$. It is convenient to absorb the Z<sub>2</sub> factor into $`\mathrm{\Lambda }_\mu `$ to obtain a renormalized vertex $`\stackrel{~}{\mathrm{\Lambda }}_\mu `$. For momenta p<sub>i</sub> and p<sub>f</sub> that are either on-shell or off-shell, the renormalized vertex involves a fermion-boson loop with a photon insertion in the fermion propagator as follows, $`\stackrel{~}{\mathrm{\Lambda }}_\mu (p_f,p_i)=`$ $`igZ_2{\displaystyle \frac{d^4k}{(2\pi )^4}S(p_fk;m)\gamma _\mu S(p_ik;m)D(k;\mu ,\mathrm{\Lambda })}.`$ (23) In the generalized model with additional Pauli-Villars regulators, the vertex is a sum of such terms, one for each term in Eq. (13). In each $`\mathrm{\Sigma }_b`$, the fermion propagator $`S(pk;m_n)`$ for mass m<sub>n</sub> is replaced by $`S(p_fk;m_n)\gamma _\mu S(p_ik;m_n)`$. Gauge invariance requires that the vertex satisfy the following Ward-Takahashi identity , $`(p_fp_i)^\mu \mathrm{\Lambda }_\mu (p_f,p_i)`$ $`=`$ $`G^1(p_f)G^1(p_i)`$ (24) $`=`$ $`\mathrm{\Sigma }(p_i)\mathrm{\Sigma }(p_f).`$ (25) This is satisfied when the photon couples to all fermion propagators in the same way, including those introduced as Pauli-Villars regulators. In general, the vertex function can be decomposed in terms of charge and magnetic form factors F<sub>1</sub> and F<sub>2</sub>. For the off-mass-shell case, there is an additional form factor F<sub>3</sub>. Moreover, all form factors depend upon $`p/_i`$ and $`p/_f`$. In order to have scalar form factors, it is necessary to project with the operators L<sup>±</sup> and to commute the $`p/_i`$ and $`p/_f`$ toward the projectors so that they may be replaced by $`\rho _i`$W<sub>i</sub> or $`\rho _f`$W<sub>f</sub>, where W<sub>i</sub> = $`\sqrt{p_i^2}`$ and W<sub>f</sub> = $`\sqrt{p_f^2}`$. This analysis is carried out in Appendix B. It produces $$\stackrel{~}{\mathrm{\Lambda }}_\mu (p_f,p_i)=\underset{\rho _f,\rho _i=\pm }{}L^{\rho _f}(p_f)\mathrm{\Lambda }_\mu ^{\rho _f,\rho _i}(p_f,p_i)L^\rho (p_i),$$ (26) where $`\mathrm{\Lambda }_\mu ^{\rho _f,\rho _i}(p_f,p_i)`$ $`=\gamma _\mu F_1^{\rho _f,\rho _i}(p_f,p_i)+i\sigma _{\mu \nu }q^\nu F_2^{\rho _f,\rho _i}(p_f,p_i)+q_\mu F_3^{\rho _f,\rho _i}(p_f,p_i),`$ (27) and q = p<sub>i</sub> \- p<sub>f</sub>. Each scalar form factor is a different function depending on the values of $`\rho _f`$ and $`\rho _i`$, e.g., F$`{}_{}{}^{+}{}_{1}{}^{}`$ is different from F$`{}_{}{}^{++}{}_{1}{}^{}`$. We shall return to this point shortly. For the on-shell situation, owing to time-reversal invariance, one only has F$`{}_{}{}^{++}{}_{1}{}^{}`$(q<sup>2</sup>) and F$`{}_{}{}^{++}{}_{2}{}^{}`$(q<sup>2</sup>), which are the usual charge and magnetic form factors of the proton. With three fermion masses and three boson masses as parameters, the generalized model allows a reasonable fit to the proton’s electromagnetic form factors. Figure 3 shows F$`{}_{}{}^{++}{}_{1}{}^{}`$(q<sup>2</sup>) and F$`{}_{}{}^{++}{}_{2}{}^{}`$(q<sup>2</sup>) in comparison with the dipole form F<sub>dipole</sub> = (1 + Q<sup>2</sup>/0.71 GeV<sup>2</sup>)<sup>-2</sup> that often is used to characterize experimental form factors. The parameter values used are: m= .38, m<sub>1</sub> = .56, m<sub>2</sub> = .61, $`\mu `$ = .79, $`\mathrm{\Lambda }_1`$ = .85 and $`\mathrm{\Lambda }_2`$ = .90, all in GeV. The bound state is at M = .93826 GeV. The anomalous magnetic moment of the composite nucleon is $`\kappa =2MF_2^{++}(0)=2.086`$, which may be compared with $`\kappa _{\mathrm{proton}}=1.79`$. A parallel analysis may be made for couplings of an elementary pion to the quark by adding a pseudoscalar $`\pi `$-quark interaction $`_5`$g$`{}_{\pi }{}^{}\overline{\psi }`$(x)$`\gamma _5\stackrel{}{\tau }\psi `$(x) $`\stackrel{}{\pi }`$(x) to the lagrangian. Figure 2 shows one pion insertion into the propagator. This produces $$G(p_f)g_\pi \stackrel{}{\tau }\widehat{\varphi }\mathrm{\Lambda }_5(p_f,p_i)G(p_i),$$ (28) where $`\stackrel{}{\tau }\widehat{\varphi }=\tau _+\varphi _{}+\tau _{}\varphi _++\tau _3\varphi _0`$, with $`\varphi _\pm `$ and $`\varphi _0`$ being isospin wave functions for $`\pi ^\pm `$ and $`\pi ^0`$ mesons. A renormalized pion-nucleon vertex function is calculated from a fermion-boson loop graph with a pseudoscalar insertion on the fermion, as follows, $`\stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i)=`$ $`igZ_2{\displaystyle \frac{d^4k}{(2\pi )^4}S(p_fk;m)\gamma _5S(p_ik;m)D(k;\mu ,\mathrm{\Lambda }_1)}.`$ (29) In the generalized model, the pion-nucleon vertex function is a sum of such terms, one for each term in Eq. (13). In each $`\mathrm{\Sigma }_b`$, the fermion propagator $`S(pk;m_n)`$ for mass m<sub>n</sub> is replaced by $`S(p_fk;m_n)\gamma _5S(p_ik;m_n)`$. Again it is necessary to rearrange terms and to project in order to have scalar form factors. This produces (see Appendix B for details) $$\stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i)=\underset{\rho _f,\rho _i=\pm }{}L^{\rho _f}(p_f)\gamma _5F_5^{\rho _f,\rho _i}(p_f,p_i)L^{\rho _i}(p_i).$$ (30) Figure 3 shows the resulting $`\pi `$N form factor F$`{}_{}{}^{++}{}_{5}{}^{}`$(q<sup>2</sup>) for on-mass-shell nucleon momenta. It is quite similar to the magnetic form factor. When one leg of the vertex function is off the mass shell, the form factors differ from the on-shell results. We wish to relate the off-shell effects to those appropriate to a hadronic vertex that is sandwiched between elementary Dirac propagators. For this purpose, it is necessary to incorporate off-shell effects from the propagators into the off-shell vertex function and form factors. In general, one encounters an off-shell vertex function sandwiched between propagators, as follows, $$\stackrel{~}{G}(p_f)\stackrel{~}{\mathrm{\Lambda }}(p_f,p_i)\stackrel{~}{G}(p_i).$$ (31) The renormalized propagator of the composite system may be written as $$\stackrel{~}{G}(p)=\frac{Z^{(+)}(p)L^{(+)}(p)}{Z_2(W_pM)}+\frac{Z^{()}(p)L^{()}(p)}{Z_2(W_pM)},$$ (32) where $`Z^{(\pm )}(p)=(\pm W_pM)/[1A(p^2)\pm W_pB(p^2)]`$ are scalar functions. In the limit that $`W_p+M`$, $`Z^{(+)}Z_2`$, and in the limit that $`W_pM`$, $`Z^{()}Z_2`$. For a point particle the factors $`Z^{(\pm )}/Z_2`$ are unity, i.e., an elementary Dirac propagator may be written in the same way with $`Z^{(\pm )}/Z_2`$ factors replaced by unity. Thus, these factors carry off-shell effects due to the propagator. A factor $`\sqrt{Z^{(\pm )}(p)/Z_2}`$ from each propagator in Eq. (31) is redistributed to the vertex function in order to obtain a vertex function that is suitable for use with the elementary Dirac propagator. The remaining $`\sqrt{Z^\pm (p)/Z_2}`$ factor in the propagators should be distributed to vertex functions preceding or following the ones indicated in Eq. (31). Figure 4 shows the variation with off-shell momentum p<sup>2</sup> for the F$`{}_{1}{}^{++}(M^2,Q^2,p^2)`$ form factor, with $`p^2`$ being the off-shell momentum. A factor $`\sqrt{Z^{(+)}(p)/Z_2}`$ is included for the off-shell leg. Similar results are obtained for the F$`{}_{}{}^{++}{}_{2}{}^{}`$ and F$`{}_{}{}^{++}{}_{5}{}^{}`$ form factors. Roughly, when p<sup>2</sup> varies from .8 M<sup>2</sup> to 1.2 M<sup>2</sup>, the form factor varies from 0.8 to 1.4 times the on-shell form factor. The off-shell variation of form factors is stronger than has been found in the work of Tiemeijer and Tjon or that of Naus and Koch . For couplings between $`+`$ and $``$ states, the form factors generally are off shell because the momentum $`p`$ of the negative state differs from $`W_p=M`$, where $`W_p=\sqrt{p^2}`$. Typically, $`+`$ to $``$ couplings are evaluated near $`W_p=+M`$, and thus they should include a factor $`\sqrt{Z^{()}(p)/Z_2}`$ from the negative-energy propagator in order to be compared with elementary couplings. Off-shell dependence of the F$`{}_{}{}^{+}{}_{1}{}^{}`$, F$`{}_{}{}^{+}{}_{2}{}^{}`$ and F$`{}_{}{}^{+}{}_{5}{}^{}`$ form factors is different in each case. It is shown in Figures 5, 6 and 7. In each case, the F<sup>+-</sup> form factor is shown as the ratio to the on-shell F<sup>++</sup> form factor, and a factor $`\sqrt{Z^{()}(p)/Z_2}`$ is included. The composite nucleon model gives nontrivial modifications of the form factors with off-shell momentum. Although the pure pseudoscalar operator $`\gamma _5`$ appears for each $`\rho _f`$ and $`\rho _i`$ value in the pion-nucleon vertex function, the form factors differ, i.e., F$`{}_{5}{}^{+}`$F$`{}_{}{}^{+}{}_{5}{}^{}`$, as mentioned above. It is instructive to compare with an elementary vertex that contains a fraction $`\lambda `$ of pseudovector and $`1\lambda `$ of pseudoscalar couplings as follows, $$\mathrm{\Lambda }_5^{\mathrm{elem}}(p_f,p_i)=\lambda \gamma _5\frac{p/_ip/_f}{2M}+(1\lambda )\gamma _5.$$ (33) Expanding by use of the projection operators and specializing to on-mass-shell kinematics yields $`\mathrm{\Lambda }_5^{\mathrm{elem}}(p_f,p_i)={\displaystyle \underset{\rho _f,\rho _i=\pm }{}}L^{\rho _f}(p_f)\gamma _5\left[\lambda {\displaystyle \frac{\rho _f+\rho _i}{2}}+1\lambda \right]L^{\rho _i}(p_i).`$ (34) On mass shell, the $`++`$ vertex is $`\gamma _5`$ independent of the mixing parameter $`\lambda `$. The $`+`$ vertex is proportional to $`(1\lambda )`$ and thus is suppressed for pseudovector coupling. A measure of the fraction of pseudovector coupling shows up in the ratio of $`+`$ and $`++`$ form factors. For the composite nucleon model, we define an equivalent pseudovector fraction in order to give a simple interpretation of the different $`++`$ and $`+`$ couplings as follows, $$\lambda =1\left(\frac{F_5^+(p^{},p)}{F_5^{++}(p^{},p)}\sqrt{\frac{Z^{()}(p)}{Z_2}}\right)_{W_p^{}=W_p=+M}.$$ (35) Figure 8 shows this ratio for the model nucleon. In the low Q range, the composite model produces 75% pseudovector coupling of the pion starting from a pseudoscalar coupling to the quark. (If the factor $`\sqrt{Z^{()}(p)/Z_2}`$ were omitted, it would be 94% pseudovector.) At Q $``$ 1 GeV/c, the vertex becomes closer to pseudoscalar. To summarize this section, the model nucleon has realistic charge and magnetic form factors. The pion form factor is similar to the magnetic one and the $`\pi `$N vertex is about 75% pseudovector and 25% pseudoscalar. Couplings between $`+`$ and $`++`$ states differ, which is a general feature of off-shell vertices. ## IV Second-order interactions - the virtual photopion amplitude Using the couplings discussed in the previous section, we consider a virtual photopion production process. This involves inserting a photon and a pion in all possible ways into the propagator and extracting the scattering amplitude as the residue of the poles in G(p<sub>f</sub>) and G(p<sub>i</sub>) as before. We also consider a standard hadronic treatment of the same process for comparison. ### A Composite nucleon analysis For the process in which a nucleon with initial momentum p<sub>i</sub> absorbs a photon of momentum q, propagates with momentum p<sub>i</sub> \+ q, and subsequently emits a pion of momentum r, ending up with momentum p<sub>f</sub>, where p<sub>i</sub> \+ q = p<sub>f</sub> \+ r, the resulting amplitude is shown in Fig. 9 and is given by (omitting isospin factors) $`V_{5,\mu }(p_f,p_i`$ $`+q,p_i){\displaystyle }_\rho \overline{u}(p_f)g_\pi \gamma _5F^{+,\rho }_5(p_f,p_f+r)G^\rho (p_i+q)\stackrel{~}{\mathrm{\Lambda }}^{\rho ,+}_\mu (p_i+q,p_i)u(p_i).`$ (36) For the crossed process in which the nucleon first emits a pion of momentum r and subsequently absorbs a photon of momentum q, ending up with the same momentum p<sub>f</sub>, the amplitude is (omitting isospin factors) $`V_{\mu ,5}(p_f,p_i+r,p_i){\displaystyle \underset{\rho }{}}\overline{u}(p_f)\stackrel{~}{\mathrm{\Lambda }}_\mu ^{+,\rho }(p_f,p_fq)G^\rho (p_ir)g_\pi \gamma _5F_5^{\rho ,+}(p_ir,p_i)u(p_i)`$ (37) These contributions to the photopion amplitude will be referred to as “Born” terms. In the analysis, two factors of Z<sub>2</sub> arise, one from the external wave functions and another from the pole term of the propagator. These factors are absorbed into the two vertex functions so that all quantities appearing in Eq. (36) and (37) are renormalized. Renormalized photon vertex function $`\stackrel{~}{\mathrm{\Lambda }}_\mu ^{\rho _f,\rho _i}`$ is defined in Eq. (27) in terms of form factors F$`{}_{1}{}^{}{}_{}{}^{\rho _f,\rho _i}`$, F$`{}_{2}{}^{}{}_{}{}^{\rho _f,\rho _i}`$, and F$`{}_{3}{}^{}{}_{}{}^{\rho _f,\rho _i}`$. Propagation has been split into separate factors for $`\rho =+`$ and $`\rho =`$ states using covariant projection operators. Note that G<sup>+</sup> contains the nucleon pole term and the excited states, which in this case are quark and diquark scattering states. Similarly, G<sup>-</sup> is the negative-energy propagation that occurs in Z-graphs. However, the standard Z-graph is based on noncovariant projection of the propagator and this causes some differences when nucleon momenta are not close to the mass shell. All of these elements arise also in a hadronic description. The variation of V<sub>5,μ</sub> with momentum transfer is characterized roughly by F(q<sup>2</sup>) G<sup>+</sup>(p<sub>i</sub> +q) F(r<sup>2</sup>), where F(q<sup>2</sup>) is a typical form factor and G<sup>+</sup> is the positive-energy propagator. At large q<sup>2</sup> and r<sup>2</sup>, these contributions become small owing to the form factors involved. A similar estimate holds for V<sub>μ,5</sub>. Excited states of the nucleon do little to alter this behavior because they involve form factors that typically fall faster with increasing momentum transfer than the nucleon’s form factors. In addition, there are contact-like terms as indicated in Fig. 9 that differ from those that arise in a hadronic description. They correspond to the two orders in which the photon and pion interact with a constituent fermion within a single vertex. Omitting isospin factors, they are defined by, $`C_{5,\mu }(p_f,p_i+q,p_i)\overline{u}(p_f)[igZ_2{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}S(p_fk;m)g_\pi \gamma _5S(p_i+qk;m)`$ (38) $`\gamma _\mu S(p_ik;m)D(k;\mu ,\mathrm{\Lambda }_1)]u(p_i)`$ (39) (40) $`C_{\mu ,5}(p_f,p_ir,p_i)\overline{u}(p_f)[igZ_2{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}S(p_fk;m)\gamma _\mu S(p_irk;m)`$ (41) $`g_\pi \gamma _5S(p_ik;m)D(k;\mu ,\mathrm{\Lambda }_1)]u(p_i).`$ (42) Initial and final states are on-shell positive-energy states, i.e., $`p_i^2=p_f^2=M^2`$ and $`\rho _i=\rho _f=+`$. In the generalized model with additional Pauli-Villars regulators, $`C_{5,\mu }`$ and $`C_{\mu ,5}`$ terms become sums of terms of the form given in Eqs. (40) and (42). In each $`\mathrm{\Sigma }_b`$ of Eq. (13), the fermion propagator $`S(p_fk;m_n)`$ is replaced by $`S(p_fk;m_n)g_\pi \gamma _5S(p_i+qk;m_n)\gamma _\mu S(p_ik;m_n)`$ to obtain $`C_{5,\mu }`$, or by $`S(p_fk;m_n)\gamma _\mu S(p_ir;m_n)g_\pi \gamma _5S(p_ik;m_n)`$ to obtain $`C_{\mu ,5}`$. Finally, there is an amplitude that results from the photon coupling to the charged pion. This is referred to as the pion-in-flight amplitude, and it takes the form $`A_\mu ^\pi (p_f,p_i)=g_\pi e\stackrel{}{\tau }T_3\widehat{\varphi }\overline{u}(p_f)\stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i)u(p_i)G_\pi (rq)J_\mu ^\pi ,`$ (43) where $`eT_3`$ is the charge operator for the pion, and $$J_\mu ^\pi =2r_\mu q_\mu .$$ (44) The total amplitude for photopion production is the sum of Born and contact-like parts, with appropriate isospin factors included, and the pion-in-flight term. $`A_\mu (p_f,q,p_i)=\stackrel{}{\tau }\widehat{\varphi }\widehat{e}V_{5,\mu }(p_f,p_i+q,p_i)+\widehat{e}\stackrel{}{\tau }\widehat{\varphi }V_{\mu ,5}(p_f,p_ir,p_i)`$ (45) $`+\stackrel{}{\tau }\widehat{\varphi }\widehat{e}C_{5,\mu }(p_f,p_i+q,p_i)+\widehat{e}\stackrel{}{\tau }\widehat{\varphi }C_{\mu ,5}(p_f,p_ir,p_i)+A_\mu ^\pi (p_f,p_i)`$ (46) Note that the order of isospin factors is important as they do not commute. Gauge invariance implies conservation of the EM current, viz., $`q^\mu A_\mu =0`$ when the pion is on mass shell, i.e., $`r^2=m_\pi ^2`$. When the pion is off mass shell, there is in general a nonzero result proportional to $`G_\pi ^1(r)=r^2m_\pi ^2`$. The required form is realized in the photo-pion amplitude A<sub>μ</sub> because of the following Ward-Takahashi identities. $`q^\mu V_{\mu ,5}=\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i+q)u(p_i)`$ (47) $`q^\mu V_{5,\mu }=\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_ir,p_i)u(p_i),`$ (48) $`q^\mu C_{\mu ,5}=\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i)u(p_i)\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i+q)u(p_i)`$ (49) $`q^\mu C_{5,\mu }=\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_ir,p_i)u(p_i)\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i)u(p_i).`$ (50) These identities may be derived by use of Eqs. (25), (40) and (42). In the contact-like terms, one needs to use the elementary Ward-Takahashi identity for Dirac propagators $$q^\mu S(p+q;m)\gamma _\mu S(p;m)=S(p;m)S(p+q;m).$$ (51) The pion-in-flight term is rewritten in terms of a commutator involving the nucleon’s charge operator, $`\widehat{e}`$, using the isospin identity $`e\stackrel{}{\tau }T_3\widehat{\varphi }=[\widehat{e},\stackrel{}{\tau }\widehat{\varphi }]`$. Its contribution to the divergence of the amplitude then is $`q^\mu A_\mu ^\pi (p_f,p_i)=[\widehat{e},\stackrel{}{\tau }\widehat{\varphi }]\overline{u}(p_f)g_\pi \stackrel{~}{\mathrm{\Lambda }}_5(p_f,p_i)u(p_i)G_\pi (rq)\left[G_\pi ^1(r)G_\pi ^1(rq)\right]`$ (52) Contributions to $`q^\mu A_\mu `$ from the Born terms are cancelled exactly by the contributions from the contact-like terms that have the same isospin factors, and the remaining contributions from the contact terms are cancelled by the second term from the pion-in-flight contribution. This leaves only a term proportional to $`G_\pi ^1(r)`$ that vanishes for an on-shell pion. The full amplitude is gauge invariant and the presence of the contact-like terms is essential for this result. The distinguishing feature of the contact-like terms is that no propagator for the composite system occurs between interactions. Thus, there is not a separate form factor for each interaction. However, the contact-like terms do depend upon the momentum transfer. They differ from a form factor mainly by the presence of an extra fermion propagator in the loop integrals of Eqs. (40) and (42). If the extra propagator lines were shrunk to a point, the contact-like terms would be related to form factors at momentum transfer q-r. This suggests that the contact terms should behave like F((q-r)<sup>2</sup>) s(q), where s(q) accounts for the extra propagator. Our calculations show that s(q) is given roughly by s(q) = $`\kappa ^2`$/($`\kappa ^2`$\- q<sup>2</sup> ), where $`\kappa `$ is a typical fermion mass. Comparing with V<sub>5,μ</sub> and V<sub>μ,5</sub>, the contact-like terms fall more slowly with increasing momentum transfer and ultimately they dominate the scattering. ### B Elementary particle with form factors analysis A standard treatment of meson-exchange currents in nuclear physics is to construct graphs corresponding to elementary particles and then to insert form factors at the vertices. The form factors are obtained from on-shell matrix elements, e.g., from phenomenological fits to electron scattering data for a free proton target. Treating the composite nucleon in this way, there are Born contributions of the $`V_{5,\mu }`$ and $`V_{\mu ,5}`$ types, which are evaluated using the F<sup>++</sup> form factors, and the pion-in-flight term. We consider both pseudoscalar and pseudovector pion-nucleon coupling in the elementary particle amplitude, and there is an additional contact term $`C_\mu ^{Elem(PV)}`$ in the pseudovector case that results from gauging the derivative of the pion field. The elementary amplitude with pseudovector pion coupling is defined as, $$A_\mu ^{\mathrm{Elem}}=\stackrel{}{\tau }\widehat{\varphi }\widehat{e}V_{5,\mu }^{\mathrm{Elem}}(p_f,p_i+q,p_i)+\widehat{e}\stackrel{}{\tau }\widehat{\varphi }V_{\mu ,5}^{\mathrm{Elem}}(p_f,p_ir,p_i)+C_\mu ^{Elem(PV)}(p_f,p_i)+A_\mu ^\pi (p_f,p_i),$$ (53) where $`V_{5,\mu }^{\mathrm{Elem}}(p_f,p_i+q,p_i)\overline{u}(p_f)g_\pi F_5^{+,+}(r)\gamma _5\left({\displaystyle \frac{p/_i+q/+M}{2M}}\right){\displaystyle \frac{1}{p/_i+q/M+i\eta }}`$ (54) $`\times \left[F_1^{+,+}(q)\gamma _\mu +i\sigma _{\mu ,\nu }q^\nu F_2^{+,+}(q)\right]u(p_i).`$ (55) Similarly, the crossed contribution is $`V_{\mu ,5}^{\mathrm{Elem}}(p_f,p_ir,p_i)\overline{u}(p_f)[F_1^{+,+}(q)\gamma _\mu +i\sigma _{\mu ,\nu }q^\nu F_2^{+,+}(q)]{\displaystyle \frac{1}{p/_ir/M+i\eta }}`$ (56) $`\times g_\pi \gamma _5\left({\displaystyle \frac{M+r/p/_i}{2M}}\right)F_5^{+,+}(r)u(p_i),`$ (57) In Eq. (55), a pseudovector vertex factor $`(p/_i+q/p/_f)/2M`$ has been evaluated by use of $`\overline{u}(p_f)\gamma _5p/_f=\overline{u}(p_f)\gamma _5M`$. Similarly, in Eq. (57), a $`p/_i`$ in the pseudovector factor $`(p/_i+r/p/_i)/2M`$ has been replaced by M by use of the Dirac equation. When pseudoscalar pion coupling is used, these factors are omitted. Note that the transition matrix elements to an intermediate negative-energy state in Eqs. (55) and(57) are based upon the same form factor as for the on-shell transition to an intermediate positive-energy state in the elementary amplitudes. However, when the pion vertex is pseudovector there is reduced coupling to the negative-energy states. In the corresponding Born amplitudes of Eqs. (36) and (37), transitions are based upon off-shell vertex functions that differ in general for the two transitions. Hadronic contact terms are implied by the off-shell factors $`(p/_i+q/+M)/(2M)`$ in Eqs. (55) and $`(M+r/p/_i)/(2M)`$ in (57). For example the first factor may be rewritten as $`1+(p/_i+q/M)/(2M)`$, where the numerator of the second part cancels the propagator. A corresponding rearrangement applies to $`V_{\mu ,5}^{\mathrm{elem}}`$. As a consequence, the pseudovector vertex can be replaced by a pseudoscalar one plus hadronic contact terms in which the factor $`1/(2M)`$ replaces the nucleon propagator between the photon absorption and pion emission. However, the form factors at the pion and electromagnetic vertices remain, which is why we refer to such terms as hadronic contact terms. They are distinct from the contributions we refer to as contact-like terms in the composite model. The resulting hadronic contact terms for pseudovector pion coupling have parts in which the $`F_1^{++}`$ form factor appears at the electromagnetic vertex. These contact terms exactly cancel with the $`C_\mu ^{Elem(PV)}`$ term. This leaves only the parts of hadronic contact terms that involve the magnetic form factor $`F_2^{++}`$. For pseudoscalar pion coupling, the situation is simpler because neither the contact terms from the off-shell vertices nor the one from gauging the derivative of the pion field are present. Consequently, there is a near equivalence of the pseudoscalar and pseudovector elementary amplitudes. After the cancellations in the pseudovector elementary amplitude, the only surviving differences from the pseudoscalar elementary amplitude are the parts of hadronic contact terms involving $`F_2^{++}`$. Vertex functions in the elementary amplitude cannot be defined precisely because of a basic conflict between the use of on-shell form factors and the conservation of four-momentum. Except when q happens to be equal to the difference of two on-shell momenta, one cannot have an on-shell vertex. In the elementary Born amplitudes of Eqs. (55) and (57) we have used on-shell form factors at vertices. This is consistent with having a vertex function that does not depend on $`p_i^2`$ or $`p_f^2`$ and therefore can be evaluated with on shell initial and final momenta whose difference is the momentum transfer, i.e., $`p_i=(E_Q,0,0,Q/2)`$ and $`p_i=(E_Q,0,0,Q/2)`$, where E<sub>Q</sub> = $`\sqrt{\mathrm{M}^2+\mathrm{Q}^2/4}`$, even though these are not the four-momenta that occur in the process. The assumption that the vertex depends only on $`Q^2`$, not the off-shell momenta, is often used when the off-shell dependence of the vertex function is unknown, but it does not have a sound theoretical basis for a composite particle. A standard nonrelativistic analysis would be similar to the elementary analysis described above. Relativistic kinematics would be used but the Z-graph parts of amplitudes would be omitted and typically a pseudoscalar pion coupling would be used. Calculations based on such a definition of a nonrelativistic amplitude will be discussed in Sec. 6. To summarize this section, the composite nucleon model has features which are similar to those of a hadronic theory, which also has, at least in principle, vertex functions that are off shell, and that are different functions for the different $`\rho `$-spins. Because the off-shell vertex functions are not known, the standard hadronic analysis uses the on-shell form factors in their place. The main feature that distinguishes the composite particle analysis from an elementary particle analysis is the presence of contact-like terms. They describe scattering from the partons and are related to off-forward parton distributions. ## V Deep inelastic scattering In order to obtain a rough normalization of the contact-like terms, we consider deep inelastic scattering from the composite nucleon. It is characterized by a hadronic tensor that takes a gauge-invariant form as follows , $`{\displaystyle \frac{W^{\mu \nu }}{4\pi M}}=\left(g^{\mu \nu }{\displaystyle \frac{q^\mu q^\nu }{q^2}}\right)W_1+\left(p^\mu q^\mu {\displaystyle \frac{pq}{q^2}}\right)\left(p^\nu q^\nu {\displaystyle \frac{pq}{q^2}}\right){\displaystyle \frac{W_2}{M^2}},`$ (58) Structure functions W<sub>1</sub> and W<sub>2</sub> depend on two scalar invariants: Q<sup>2</sup> = $``$q<sup>2</sup>, and $`\nu `$ = p $``$ q/M. In the limit Q$`{}_{}{}^{2}\mathrm{}`$, with x $``$ Q<sup>2</sup>/(2M $`\nu `$) held fixed, these functions become dependent only on x as follows , $$MW_1(x,Q^2)F_1(x)=\frac{1}{2}f(x),$$ (59) and $$\nu W_2(x,Q^2)F_2(x)=xf(x).$$ (60) This scaling behavior is a consequence of scattering from point-like constituents of the nucleon, with f(x) being the the probability of scattering from a parton that carries a fraction x of the nucleon’s momentum. A simple way to obtain the forward parton distribution f(x) is to calculate W<sup>xx</sup>. Either in the lab frame, where p = (M, 0, 0, 0,) and q = ($`\nu `$, 0, 0, $`\sqrt{Q^2+\nu ^2}`$), or in the c.m. frame of the final state, the x-components of q and p vanish. For finite Q we have $$2MW_1(x,Q^2)=\frac{W^{xx}}{2\pi },$$ (61) and in the asymptotic limit $$f(x)=\underset{Q^2\mathrm{}}{lim}2MW_1(x,Q^2).$$ (62) Batiz and Gross have analyzed the scaling limit for a composite nucleon model that is essentially similar to the one used in this paper. Their analysis is for one space and one time dimension. They show that scaling in a general gauge involves a cancellation between a gauge-dependent part of the impulse approximation graph of Fig. 10 and a gauge-dependent part of the final-state interaction. This is related to the Ward identities of Eqs. (48) and (50) which imply that gauge invariance requires cutting both the Born and contact-like terms. Using the Landau prescription, Batiz and Gross split the impulse amplitude into a gauge-invariant part and a remainder. The gauge-invariant part of the impulse graph provides the scaling result. The gauge variant remainder cancels with part of the final-state interaction such that the resultant contribution of these parts vanishes at least as fast as 1/Q<sup>2</sup>. In this section, we follow Batiz and Gross by using the Landau prescription for three space dimensions and one time dimension. The results differ because of integrals over the angles of final state particles and because the phase space in 3D differs from that in 1D. The hadronic tensor based on the impulse graph of Fig. 10 is calculated in the c.m. frame of the final state, $$W^{\mu \nu }=\frac{1}{2}\underset{s,s_1}{}\frac{\|\stackrel{}{p}_1\|}{4\pi W}\frac{1}{4\pi }𝑑\mathrm{\Omega }_1𝒯^\mu 𝒯^\nu ,$$ (63) where W is the total energy of the final state that contains an on-shell quark of momentum p<sub>1</sub> = (E<sub>1</sub>, p<sub>1</sub>) and an on-shell boson of momentum (W - E<sub>1</sub>, -p<sub>1</sub>). Amplitude $`𝒯^\mu `$ describes the impulse approximation graph for scattering from the fermion constituent. Using the Landau prescription as in Ref. to obtain a gauge-invariant current, this is $`𝒯^\mu ={\displaystyle \frac{\sqrt{\|gZ_2\|4mM}}{(p_1q)^2m^2}}\overline{u}_1(p_1;m)[2p_1^\mu {\displaystyle \frac{2p_1q}{q^2}}q^\mu \gamma ^\mu q/+q^\mu ]u(p;M),`$ (64) where u<sub>1</sub>(p<sub>1</sub>;m) is a Dirac spinor for the quark of mass m and u(p;M) is a Dirac spinor for the nucleon of mass M. The factor $`\sqrt{\|gZ_2\|}`$u(p;M) is the vertex function for the nucleon to fragment into a quark and a boson. An equivalent form of the hadronic tensor, which we use, is $$W^{\mu \nu }=\frac{\|\stackrel{}{p}_1\|}{4\pi W}\frac{\|gZ_2\|}{4\pi }𝑑\mathrm{\Omega }_1\frac{1}{[(p_1q)^2m^2]^2}^{\mu \nu },$$ (65) where we define $`^{\mu \nu }={\displaystyle \frac{1}{2}}Tr\{(2p_1^\mu {\displaystyle \frac{2p_1q}{q^2}}q^\mu \gamma ^\mu q/+q^\mu )(p/+M)`$ (66) $`\times (2p_1^\nu {\displaystyle \frac{2p_1q}{q^2}}q^\nu \gamma ^\nu q/+q^\nu )(p/_1+m)\}.`$ (67) Specializing to $`^{xx}`$, we find $`^{xx}=2(4p_1^xp_1^x+Q^2)(pp_1+Mm)+4M\nu (qp_12p_1^xp_1^x),`$ (68) where terms not involving x-components have arisen from use of $`\gamma ^x\gamma ^x=1`$. Expressing the vectors in the c.m. frame, where p<sub>1</sub> = (E<sub>1</sub>, $`\|𝐩_1\|`$sin$`\theta _1`$cos$`\varphi _1`$, $`\|𝐩_1\|`$sin$`\theta _1`$sin$`\varphi _1`$ , $`\|𝐩_1\|`$cos$`\theta _1`$ ), leads to the following expression $`^{xx}=8\|𝐩_1\|^2sin^2\theta _1cos^2\varphi _1\left[p^0E_1p^z\|𝐩_1\|cos\theta _1+MmM\nu \right]`$ (69) $`+4M\nu \left[q^0E_1q^z\|𝐩_1\|cos\theta _1\right]+2Q^2\left[p^0E_1p^z\|𝐩_1\|cos\theta _1+Mm\right]`$ (70) Carrying out the elementary angle integrations produces, $$\frac{W^{xx}}{2\pi }=\frac{\|gZ_2\|\|𝐩_1\|}{8\pi ^2W}\left(\frac{C_0}{a^2b^2}+C_1\mathrm{ln}\left(\frac{a+b}{ab}\right)+C_2\right),$$ (71) $`a`$ $`=`$ $`Q^22q^0E_1`$ (73) $`{\displaystyle \frac{Q^2}{2x}}+{\displaystyle \frac{2x1}{2(1x)}}[m^2\mu ^2]`$ $`b`$ $`=`$ $`2q^z\|𝐩_1\|`$ (75) $`{\displaystyle \frac{Q^2}{2x}}+{\displaystyle \frac{2M^2x(1x)m^2\mu ^2}{2(1x)}}`$ $`C_0`$ $`=4M\nu q^0E_1+2Q^2(p^0E_1+Mm)`$ (78) $`+[4M\nu q^z\|𝐩_1\|+2Q^2p^z\|𝐩_1\|]{\displaystyle \frac{a}{b}}`$ $`{\displaystyle \frac{Q^2}{x}}[(Mx+m)^2F_\mu (x)]`$ $`C_1`$ $`=4\|𝐩_1\|^2(p^0E_1+MmM\nu ){\displaystyle \frac{a}{b^3}}2\|𝐩_1\|^3p^z{\displaystyle \frac{1}{b^2}}\left(1{\displaystyle \frac{3a^2}{b^2}}\right)(2M\nu q^z\|𝐩_1\|+Q^2p^z\|𝐩_1\|){\displaystyle \frac{1}{b^2}}`$ (80) $`x1`$ $`C_2=`$ $`8\|𝐩_1\|^2(p_0E_1+MmM\nu ){\displaystyle \frac{1}{b^2}}12\|𝐩_1\|^3p^z{\displaystyle \frac{a}{b^3}}`$ (82) $`x1`$ and we have defined $$F_\mu (x)=m^2(1x)+\mu ^2xM^2x(1x).$$ (83) Furthermore, the phase-space factor approaches a constant, $`{\displaystyle \frac{\|𝐩_1\|}{W}}{\displaystyle \frac{1}{2}}.`$ (84) In the expressions given above, the limiting form as Q$`{}_{}{}^{2}\mathrm{}`$ with x = Q<sup>2</sup>/(2 M $`\nu `$) fixed is indicated following the arrow. We have used a number of kinematical relations that can be found in the paper of Batiz and Gross. Note that ln((a$`+`$b)/(a$``$b)) arises in the structure function. Because a$`+`$b $``$ F<sub>μ</sub>(x)/(1 $``$ x) is independent of Q, but a $``$ b $``$Q<sup>2</sup>/x is not, there is a $`\mathrm{ln}`$(Q<sup>2</sup>) term in W<sup>xx</sup>. This is cancelled when the Pauli-Villars subtraction is made, i.e., when the parton distribution is calculated as the discontinuity of the subtracted bubble graph. Clearly, scaling in 3+1 dimensions depends upon the subtraction, which was not the case in the 1+1 dimensional analysis of Ref. . ¿From Eqs. (71),(61) and (62), we find the following parton distribution for the bubble graph with a fermion of mass m, a boson of mass $`\mu `$, and a Pauli-Villars subtraction of mass $`\mathrm{\Lambda }_1`$. $`f_b(x;m,\mu ,\mathrm{\Lambda }_1)={\displaystyle \frac{\|gZ_2\|(1x)}{16\pi ^2}}\{(Mx+m)^2({\displaystyle \frac{1}{F_\mu (x)}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}(x)}})\mathrm{ln}({\displaystyle \frac{F_\mu (x)}{F_{\mathrm{\Lambda }_1}(x)}})\}`$ (85) For the case in which additional subtractions are made, the parton distribution becomes a linear combination of terms as in Eq.( 13), i.e., $`f(x)=f_b(x;m,\mu ,\mathrm{\Lambda }_1)+\beta f_b(x;m,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)`$ (86) $`\alpha \left[f_b(x;m_1,\mu ,\mathrm{\Lambda }_1)+\beta f_b(x;m_1,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)\right]`$ (87) $`(1\alpha )\left[f_b(x;m_2,\mu ,\mathrm{\Lambda }_1)+\beta f_b(x;m_2,\mathrm{\Lambda }_1,\mathrm{\Lambda }_2)\right]`$ (88) Figure 11 shows the inelastic structure function 2MxW<sub>1</sub>(x,Q<sup>2</sup>) and its limit as Q$`\mathrm{}`$, xf(x), for the same parameters that allow a reasonable description of the nucleon form factors. For finite Q, the energy transfer $`\nu =Q^2/(2Mx)`$ also is finite. It must be greater than the total mass of the constituent quark and diquark for each combination that enters Eq. (88) in order to avoid spurious threshold effects. This restricts x not to be too close to 1. For finite Q, we show W$`{}_{1}{}^{}(x,Q)`$ only for $`\nu `$ greater than about 200 MeV above threshold. The solid line in Fig. 11 shows the asymptotic limit Q$`\mathrm{}`$, which is xf(x). Already by Q$``$ 2 GeV/c, the inelastic structure function is close to its asymptotic limit for our quark-diquark model of a nucleon. Our result for $`xf(x)`$ is more peaked than that obtained recently by Mineo, Bentz and Yazaki from consideration of a three-quark model, and both results lack sufficient strength near x = 0 in comparison with experimentally determined parton distributions. Reference considers what one should expect for xf(x) at low Q. Using QCD evolution to relate high and low Q, the nucleon’s parton distribution is found to be more peaked at low Q, and not unlike the xf(x) that we find, with a peak near x = m/M, where m is the lightest fermion mass. However, our results at Q<sup>2</sup> = 4 (GeV/c)<sup>2</sup> are more peaked than the results of Ref. at Q<sup>2</sup> = 0.16 (GeV/c)<sup>2</sup>. For x $`>`$ .6, xf(x) in Figure 11 is negative owing to the fact that the fermionic subtractions are not hermitian. This is a deficiency of the model used. Our results for the parton distribution are influenced by the choice of subtractions that have been incorporated in order to obtain a good description of the nucleon’s form factors. We have given preference to obtaining a realistic form factor in selecting parameters. Owing to the normalization factor $`\|gZ_2\|`$, the parton distribution automatically is normalized according to $$_0^1𝑑xf(x)=1.$$ (89) See Ref. and Appendix A in this regard. However, only the fermion constituent has charge and the momentum sum rule is, $$_0^1𝑑xxf(x)=.304.$$ (90) Thus, 70% of the momentum is carried by the diquark in this model. To summarize this section, the model for a nucleon as a bound state of a quark and diquark is consistent with scaling in deep inelastic scattering. The normalized parton distribution provides a rough normalization for the contact-like terms in the large Q limit. ## VI Calculations for meson-exchange current amplitude Virtual photopion production amplitude $`A^\mu `$ that has been defined in Eq. (46) and discussed in Sec. 4 contributes to the electromagnetic current in electron-deuteron scattering. For this case, the emitted pion is absorbed on a second nucleon and, in general, there is a loop integration involving the pion momentum and the deuteron wave functions of initial and final states. For electron-deuteron scattering, the loop integration receives important contributions from the quasifree kinematics indicated in Figure 12. Each nucleon in the deuteron, only one of which is shown, has initial momentum $`p_i=\frac{1}{2}`$P $`\frac{1}{4}`$q. Figure 12 shows the nucleon which absorbs a photon of momentum q and emits a pion of momentum $`r=\frac{1}{2}`$q, ending up with momentum $`p_f=\frac{1}{2}`$P $`+\frac{1}{4}`$q. The second nucleon, not shown, also has initial momentum $`p_i=\frac{1}{2}`$P $`\frac{1}{4}`$q. When the pion is absorbed on the second nucleon, its final momentum also becomes $`p_f`$. This process begins and ends with the two nucleons at zero relative momentum. It is favored because the deuteron wave function is largest at zero relative momentum. Pion-in-flight terms vanish for the selected kinematics, since $`2r_\mu q_\mu =0`$. They are omitted from our calculations. A calculation using deuteron wave functions is planned for a future work. For now, we focus on the quasi-free photopion amplitude and simply vary the momentum of the space-like virtual photon: q = (0, 0, 0, Q). Although there are sixteen helicity amplitudes $`A_{\lambda _f,\lambda _i}^\mu `$, for the quasifree kinematics with collinear momenta that we consider only three amplitudes are significant. An isospin-nonflip amplitude, a, occurs in the time-component of the photopion amplitude as follows, $`A_{\lambda _f,\lambda _i}^0=a\{\widehat{e},\stackrel{}{\tau }\widehat{\varphi }\}\chi _{\lambda _f}^{}\sigma _z\chi _{\lambda _i},`$ (91) where the isospin factor involves an anticommutator. Two isospin flip amplitudes, b and c, occur in the space-vector parts of the photopion amplitude as follows, $`A_{\lambda _f,\lambda _i}^\pm =b[\widehat{e},\stackrel{}{\tau }\widehat{\varphi }]\chi _{\lambda _f}^{}\sigma _\pm \chi _{\lambda _i},`$ (92) $`A_{\lambda _f,\lambda _i}^3=c[\widehat{e},\stackrel{}{\tau }\widehat{\varphi }]\chi _{\lambda _f}^{}\sigma _z\chi _{\lambda _i},`$ (93) where the isospin factors involve a commutator. Although we calculate only the photopion amplitude, its role as a meson-exchange current in electron-deuteron scattering is of interest. In that case, the isospin wave functions $`\widehat{\varphi }`$ for the pion are replaced by the isospin operators $`\stackrel{}{\tau }_2`$ for the second nucleon. The isospin-nonflip amplitude, a, is the only one that contributes as a meson-exchange current in elastic electron-deuteron scattering. However, for breakup of the deuteron, both isospin nonflip and isospin-flip amplitudes contribute. ### A Isospin nonflip amplitude: a Figure 13 shows the absolute value of Born and contact-like contributions to amplitude $`a`$ in comparison with simple estimates of these contributions suggested in Sec. 4: $`\mathrm{V}_{0,5}+\mathrm{V}_{5,0}\mathrm{C}_1\mathrm{QF}_{\mathrm{dipole}}(\mathrm{q}^2)\mathrm{G}^+((\mathrm{p}_i+\mathrm{q})\mathrm{F}_{\mathrm{dipole}}((\mathrm{r}^2)`$ and $`\mathrm{C}_{0,5}+\mathrm{C}_{5,0}`$ C<sub>2</sub>Q F((q-r)<sup>2</sup>) S(q), where C<sub>1</sub> and C<sub>2</sub> are constants and s(q) = $`\kappa ^2`$/(Q<sup>2</sup> \+ $`\kappa ^2)`$. We find that $`\kappa ^2`$ = .20 (GeV/c)<sup>2</sup> describes the Q dependence caused by the extra quark propagator in the contact-like terms. A factor Q is included in the estimates because there is such a factor in the amplitude for kinematical reasons. The point of this comparison is to show that the contact-like contributions are decreasing with Q faster than a form factor, as expected. They nevertheless can be dominant when Q $`>`$ 1 GeV/c because the Born terms decrease even faster. There is a zero in the Born amplitude that is not reproduced in the estimate. However, the estimate is quite good for individual covariant amplitues that go into the helicity matrix element (see Appendix C for the definition of covariant amplitudes). In Fig. 14, we show $`\|a\|`$ for the Born (long dash line) and full amplitudes (solid line) for the composite nucleon. Also shown (dash line) is the elementary amplitude $`\|a\|`$ that is based upon pseudoscalar pion coupling. Finally, we show (dotted line) a nonrelativistic amplitude that is based on pseudoscalar pion coupling and the standard positive-energy propagator: $`\mathrm{\Lambda }^+(𝐩)/(p^0\sqrt{M^2+𝐩^2})`$. Thus, the nonrelativistic amplitude differs from the elementary one by omission of the Z-graph part. Form factors used in the elementary particle and nonrelativistic analyses are based on the quark-diquark model except that on-shell ++ form factors are used. For small Q, the Born contributions dominate for all cases because there is a pole in the intermediate nucleon propagator. In the vicinity of Q $`=`$ 1.2 GeV/c, two amplitudes that involve an intermediate propagator for a nucleon, i.e., the the Born and nonrelativistic amplitudes, pass through zero. They are negative at higher Q and their magnitude is 40% to 10% of the full amplitude at Q $`=`$ 3 GeV/c. The elementary amplitude based on pseudoscalar pion coupling has a zero near Q $``$ 2.8 GeV/c. It provides the best approximation to the results of the composite model but is much smaller in magnitude for Q $``$ 1.5 GeV/c. The amplitude for the composite nucleon model is dominated at large Q by the contact-like terms, which do not change sign. Although the nonrelativistic amplitude tends at large Q to a magnitude similar to that of the full amplitude, it has the opposite sign and is not a useful approximation to the full result. Figure 15 shows the contact-like amplitude of the composite model, (solid line). The part of the Born amplitude that comes from excited states and Z-graphs is shown by the long dash line. It has been calculated by evaluating Eqs. (37) and (36) with the positive-energy propagator, $`\mathrm{\Lambda }^+(𝐩)/(p^0\sqrt{M^2+𝐩^2})`$ and then subtracting that result from the Born amplitude based on the full propagator of the composite model. Next we show by the dash dot line the sum of parts of the composite nucleon amplitude that do not come from the Born terms with the positive propagator, i.e., the sum of contact-like parts, excited-state parts and Z-graph parts. Thus, the dot-dash line shows the sum of the amplitudes used in the solid and long-dash lines. The dash line shows the Z-graph part of the elementary amplitude based on pseudoscalar pion coupling. The excited-states-plus-Z-graph part of the composite-nucleon Born amplitude (long dash line in Fig. 15) is larger than the Z-graph part of the elementary Born amplitude (dashed line) at low Q, but it decreases rapidly with Q. Because a less point-like structure for the composite nucleon would be expected to provide even smaller contributions at large Q from excited states and Z-graphs, the smallness of the Born contributions of the composite model at large Q is expected to hold more generally. The Q-dependence of the pseudoscalar Z-graph contribution (dashed line) is notable for its similarity to that of the contact-like contribution of the composite model (solid line). Because the pion vertex of the composite model is about 75% pseudovector, we consider next the same set of comparisons using an elementary amplitude in which the pion coupling is pseudovector. In this case, we also include the contact term that is implied by gauging the derivative of the pion field, $`C_\mu ^{Elem(PV)}`$. Figure 16 shows that the elementary amplitude based on pseudovector pion coupling (dashed line) provides a poorer approximation to the composite nucleon result (solid line). This is because it has a zero near 1 GeV/c and has the wrong sign at large Q. Figure 17 shows the difference between the pseudovector elementary amplitude and the nonrelativistic amplitude that uses pseudoscalar pion coupling by the dashed line. We find that the use of pseudoscalar pion coupling in the elementary amplitude provides a better approximation to the amplitude of the composite model. This is because it has a large Z-graph contribution that approximates the contact-like contribution of the composite model. As mentioned, the pseudoscalar and pseudovector pion couplings produce different results only because of the magnetic couplings of the photon. If the photon were to couple only via the charge current, $`\gamma ^\mu F_1(Q)`$, the pseudoscalar and pseudovector elementary amplitudes that we consider would be equal. However, the magnetic part of the charge current, $`\sigma ^{\mu \nu }q_\nu F_2(Q)`$, changes this. For pseudoscalar coupling, the Z-graph amplitude shown by the dashed line in Figure 15 is proportional to $`F_1(Q)+F_2(Q)=G_M(Q)`$, whereas for pseudovector coupling, the dashed line in Figure 17 is proportional to $`F_1F_2(Q)Q^2/(4M^2)=G_E(Q)`$. The effect of the magnetic parts explains the different results in these graphs. Low-energy theorems that apply to the photopion amplitude at low Q arise from chiral invariance. Because the composite nucleon model has essentially pseudovector pion coupling, which is consistent with chiral invariance, one might expect the pseudovector elementary amplitude to provide a better approximation to the composite nucleon results. This expectation fails at large Q because of the important contributions of contact-like terms. ### B Isospin-flip amplitudes: b and c Isospin-flip amplitudes $`\|b\|`$ and $`\|c\|`$ are shown in Figures (18) and (19). These amplitudes do not exhibit a pole at $`Q=0`$ like the one in the isospin-nonflip amplitude, a. Each vertex in the b and c amplitudes has a factor Q, thus cancelling the $`1/Q^2`$ from the propagator. Consequently, the isospin flip amplitudes are much smaller than a at small Q, but they can be comparable at large Q. Elementary amplitudes based upon pseudoscalar pion coupling and pseudovector pion coupling are equivalent for the b and c amplitudes, and thus we show only the pseudoscalar elementary amplitude. This equivalence results because the hadronic contact terms involving the $`F_2`$ electromagnetic form factor give a vanishing contribution for the kinematics that we consider. Figure (18) shows that Born and nonrelativistic results for $`\|b\|`$ are very close to one another at small Q. However, the Born amplitude is significantly smaller at larger Q. The full composite model result is close to that of the elementary amplitude over the entire range of Q. Born and nonrelativistic results both omit Z-graphs, whereas the full and elementary results both include Z-graphs. It is apparent that the Z-graphs make a significant contribution at Q $`=`$ 0, lowering the full and elementary results in comparsion with the Born and nonrelativistic ones. Contact-like terms of the composite nucleon cause the difference between full and long-dash lines: these are significant but not dominant in the way they are for the isospin-nonflip amplitude, a. Figure (19) shows that the full composite model provides a much larger result for $`\|c\|`$ than is obtained from the Born, elementary or nonrelativistic amplitudes. Thus, $`\|c\|`$ and $`\|a\|`$ amplitudes show dominance of contact-like contributions at large Q, but $`\|b\|`$ does not. It is clear that one would like to have better control of the normalization of contact-like terms in order to determine the transition point where they may become dominant contributions to electron scattering from nuclei. However, the present model suggests that this could be near 1 GeV/c for the isoscalar MEC appropriate to elastic electron-deuteron scattering, based upon the strong dominance of contact-like contributions to $`\|a\|`$. For the isospin-flip MEC contributions, which are relevant to electrodisintegration of the deuteron, each of the amplitudes a, b and c contributes. A more complete calculation is required to see if the contact-like terms may dominate the MEC at large Q. Use of pseudoscalar pion coupling improves the agreement with results of the composite model significantly for $`\|a\|`$, and is equivalent to pseudovector pion coupling for $`\|b\|`$ and $`\|c\|`$. The fact that pseudoscalar pion coupling seems to work fairly well is not because it provides a description of the underlying physics, which requires consideration of scattering from the quarks at large Q. ## VII Conclusion A simple model of a composite nucleon is developed in which a fermion and a boson, representing quark and diquark constituents of the nucleon, form a bound state owing to a contact interaction. Photon and pion couplings to the quark provide vertex functions for the photon and pion interactions with the composite nucleon. By introducing and exploiting cutoff parameters of the Pauli-Villars type, realistic electromagnetic form factors are obtained for the proton. When a pseudoscalar pion-quark coupling is used, the pion-nucleon coupling is 75% pseudovector. The small quark mass produces a vertex behavior close to that expected from chiral invariance. A virtual photopion amplitude is considered in which there are two types of contributions: hadronic contributions where the photon and pion interactions have an intervening propagator of the nucleon, or its excited states, and contact-like contributions where the photon and pion interactions occur within a single vertex. Relative normalization of the two types of contribution is controlled by Ward-Takahashi identities at low momentum transfer. At high momentum transfer, scaling behavior is obtained for the composite nucleon already by Q $``$ 2 GeV/c. This provides a rough normalization of the contact-like parts because the parton distribution is normalized (see Eq. (89)). However, our model of a composite nucleon as a bound state of a quark and diquark yields a parton distribution that is peaked near x $`=`$ m/M, the ratio of quark to nucleon mass, whereas the data suggest much less peaking and more strength at low x values than the model gives. Calculations for the virtual photopion amplitude are performed using kinematics appropriate to its occurrence as a meson-exchange current in electron-deuteron scattering. The results show that the contact-like terms dominate the meson-exchange current for Q $`>`$ 1 GeV/c for the case of elastic electron-deuteron scattering. As Q increases, the dominance of the contact-like terms over the Born terms of the composite nucleon can become very large, suggesting that hadronic processes become unimportant when this occurs. Our results indicate that contact-like terms still have substantial Q dependence when they become dominant. For the inelastic electron deuteron scattering, both isospin-nonflip and isospin-flip parts of the photopion amplitude can contribute. Two of the three contributing amplitudes are dominated at large Q by contact-like terms and the other is not. A more complete calculation using deuteron wave functions is needed in order to understand the role of contact-like contributions in deuteron breakup. Off-shell effects in the hadronic vertex functions are found to be significant in the composite model. They cause a significant suppression of Born contributions to the virtual photopion amplitude for Q $``$ 1 GeV/c. This result is model-dependent, but it suggests that use of on-shell form factors could be a poor approximation for momenta that are significantly off the mass shell. An elementary amplitude based upon pseudovector pion coupling fails to provide a useful approximation to the full result of the composite model for the isospin-nonflip amplitude. This can be improved somewhat by using pseudoscalar pion coupling in the elementary amplitude. The increased Z-graph contribution gives a better approximation to the contact-like terms of the composite nucleon, but not to the underlying physics. Compositeness requires contact-like terms in second-order interactions. They have a direct connection to off-forward parton distributions and can dominate the scattering at large Q as they contain the leading partonic scattering process. Hadronic form factors and off-shell effects tend to quench the Born scattering processes that involve intermediate hadronic states. For the considered nucleon model, we find that scattering from the quark constituent can be significant at modest Q values such as Q $`>`$ 1 to 2 GeV/c. Once partonic scattering becomes dominant, it is expected to remain dominant for higher Q. Where the transition to dominance of the partonic interactions actually takes place is a matter of great interest. The model calculation of this paper suggests that this is determined by the size of contact-like contributions, or equivalently, by the size of the off-forward parton distributions. It may occur in some processes at momentum transfer as low as 1 GeV/c and seems to be likely by 2 GeV/c for the considered isoscalar meson-exchange current. Support for this work from the U.S. Department of Energy is gratefully acknowledged through DOE grant DE-FG02-93ER-40762 at the University of Maryland and DOE contract DE-AC05-84ER40150, under which the Southeastern Universities Research Association (SURA) operates the Thomas Jefferson National Accelerator Facility. S. J. W. gratefully acknowledges support of SURA under its Sabbatical Fellowship Program. ## A Self-energy and loop integrals Details that have gone into calculations but are omitted from the text are collected in this appendix. The fermion-boson self energy graph, defined in Eq. (4), vertex functions defined in Eqs. (23) and (29) and the contact-like terms defined in Eqs. (42) and (40) require evaluations of Feynman integrals and subsequent reductions of the Dirac matrices to standard forms. Integrations over loop momentum k are performed by standard methods: n propagator factors are combined by means of integrals over Feynman parameters $`\alpha _1,\alpha _2,\mathrm{},\alpha _{n1}`$, into a single denominator function of the form $`[(k\mathrm{})^2F+i\eta ]^n`$, where the shift vector $`\mathrm{}^\mu `$ and the function F depend upon the external momenta and Feynman parameters. Numerator functions involve one power of the loop momentum, $`k^\mu `$ for each fermion propagator. Two divergent k-integrations arise and these are evaluated by using subtractions. The required formulas are, $`ig{\displaystyle \frac{d^4k}{(2\pi )^4}\left(\frac{1}{[(k\mathrm{})^2F_\mu ]^2}\frac{1}{[(k\mathrm{})^2F_{\mathrm{\Lambda }_1}]^2}\right)\{1,k^\mu \}}={\displaystyle \frac{g}{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{F_\mu }{F_{\mathrm{\Lambda }_1}}}\right)\{1,\mathrm{}^\mu \},`$ (A1) $`2ig{\displaystyle \frac{d^4k}{(2\pi )^4}\left(\frac{1}{[(k\mathrm{})^2F_\mu ]^3}\frac{1}{[(k\mathrm{})^2F_{\mathrm{\Lambda }_1}]^3}\right)k^\mu k^\nu }=`$ (A2) $`{\displaystyle \frac{g}{16\pi ^2}}\mathrm{ln}\left({\displaystyle \frac{F_\mu }{F_{\mathrm{\Lambda }_1}}}\right)\left({\displaystyle \frac{1}{2}}g^{\mu \nu }\right)+{\displaystyle \frac{g}{16\pi ^2}}\left({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}}\right)\mathrm{}^\mu \mathrm{}^\nu .`$ (A3) In all other cases, the k-integrations can be performed before subtractions by using the formulas ( for $`n3`$), $`\{<1>,<k^\mu >,<k^\mu k^\nu >,<k^\mu k^\nu k^\sigma >\}`$ (A4) $`ig(n1)!\times {\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}{\displaystyle \frac{1}{[(k\mathrm{})^2F+i\eta ]^n}}\{1,k^\mu ,k^\mu k^\nu ,k^\mu k^\nu k^\sigma \}=`$ (A5) $`{\displaystyle \frac{(n3)!g}{(1)^{(n+1)}16\pi ^2}}\{{\displaystyle \frac{1}{F^{n2}}},{\displaystyle \frac{\mathrm{}^\mu }{F^{n2}}},{\displaystyle \frac{\mathrm{}^\mu \mathrm{}^\nu }{F^{n2}}}{\displaystyle \frac{g^{\mu \nu }}{2(n3)F^{n3}}},{\displaystyle \frac{\mathrm{}^\mu \mathrm{}^\nu \mathrm{}^\sigma }{F^{n2}}}{\displaystyle \frac{n^{\mu \nu \sigma }}{2(n3)F^{n3}}}\}`$ (A6) where $$n^{\mu \nu \sigma }g^{\mu \nu }\mathrm{}^\sigma +g^{\nu \sigma }\mathrm{}^\mu +g^{\sigma \mu }\mathrm{}^\nu .$$ (A7) Considering the self energy of an elementary fermion-boson bubble graph, we have $`\mathrm{\Sigma }_b(p;m,\mu ,\mathrm{\Lambda }_1)=ig{\displaystyle \frac{d^4k}{(2\pi )^4}\frac{(p/k/+m}{(pk)^2m^2+i\eta }\left(\frac{1}{k^2\mu ^2+i\eta }\frac{1}{k^2\mathrm{\Lambda }_1^2+i\eta }\right)}.`$ (A8) Using the Feynman parameterization $$\frac{1}{ab}=_0^1𝑑\alpha \frac{1}{[\alpha a+(1\alpha )b]^2},$$ (A9) to combine denominators, we have in this case the shift vector $`\mathrm{}=\alpha p`$, and denominator functions $`F_\mu `$ and $`F_{\mathrm{\Lambda }_1}`$, where the general form is $$F_\mathrm{\Lambda }=\alpha \mathrm{\Lambda }^2+(1\alpha )m^2\alpha (1\alpha )p^2.$$ (A10) Integrating over loop momentum produces the two scalar parts defined in Eq. (15), as follows, $$\{A(p^2),B(p^2)\}=\frac{g}{16\pi ^2}_0^1𝑑\alpha \mathrm{ln}\left(\frac{F_\mu }{F_{\mathrm{\Lambda }_1}}\right)\{\alpha ,m\}.$$ (A11) Using these formulas and the condition $`MA(M^2)+B(M^2)=1`$, one may determine the coupling constant g such that there is a bound state of mass M, where M $`<`$ m + $`\mu `$. The corresponding formulas for $`A^{}(p^2)`$ and $`B^{}(p^2)`$ are obtained by differentiating with respect to $`p^2`$, $`\{A^{}(p^2),B^{}(p^2)\}={\displaystyle \frac{g}{16\pi ^2}}{\displaystyle _0^1}𝑑\alpha \alpha (1\alpha )\left({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}}\right)\{\alpha ,m\}.`$ (A12) Wave-function renormalization constant $`Z_2`$ has contributions from the elementary bubble graph that may be expressed, using Eqs. (17) and (A12), as follows, $`Z_2^1={\displaystyle \frac{g}{16\pi ^2}}{\displaystyle _0^1}d\alpha \{2M(M\alpha +m)\alpha (1\alpha )({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}})\alpha \mathrm{ln}({\displaystyle \frac{F_\mu }{F_{\mathrm{\Lambda }_1}}})\}`$ (A13) and for $`Z_2`$, $`p^2=M^2`$ in $`F_\mu `$ and $`F_{\mathrm{\Lambda }_1}`$. Using the identity, $`{\displaystyle _0^1}d\alpha \{(1\alpha )(m^2M^2\alpha ^2)({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}})+(2\alpha 1)\mathrm{ln}\left({\displaystyle \frac{F_\mu }{F_{\mathrm{\Lambda }_1}}}\right)\}=0,`$ (A14) an equivalent expression for $`Z_2^1`$ is, $`Z_2^1`$ $`={\displaystyle \frac{g}{16\pi ^2}}{\displaystyle _0^1}d\alpha (1\alpha )\{(M\alpha +m)^2({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}})\mathrm{ln}({\displaystyle \frac{F_\mu }{F_{\mathrm{\Lambda }_1}}})\}`$ (A15) The integral here is the same as for the contribution of the elementary bubble graph to the normalization of the parton distribution, showing that the factor $`Z_2`$ guarantees the normalization as in Eq. (89). Equation (A14) can be verified by integrating by parts the term involving $`2\alpha 1`$. ## B Three-point functions Three-point functions required for photon and pion vertices are defined by Eqs. (23) and (29). They are calculated numerically from formulas involving integrations over two Feynman parameters, $`\alpha _1`$ and $`\alpha _2`$. The denominator function F<sub>Λ</sub> that results from combining the propagators in the three-point function, assuming a generic mass $`\mathrm{\Lambda }`$ for the boson, is $`F_\mathrm{\Lambda }=`$ $`\alpha _1\mathrm{\Lambda }^2+\alpha _2(m^2p_i^2)+(1\alpha _1\alpha _2)(m^2p_f^2)+\mathrm{}^2,`$ (B1) and the shift vector is $$\mathrm{}=\alpha _1p_i+(1\alpha _1\alpha _2)p_f.$$ (B2) Moments of loop momentum need to be expanded in terms of the independent external momenta, which we choose to be $`p_i`$ and $`p_f`$. For this expansion, we define, $`<1>=C_0,`$ (B3) $`<k^\mu >=C_{11}p_i^\mu +C_{12}p_f^\mu ,`$ (B4) $`<k^\mu k^\nu >=`$ $`C_{21}p_i^\mu p_i^\nu +C_{22}p_f^\mu p_f^\nu +C_{23}(p_i^\mu p_f^\nu +p_i^\nu p_f^\mu )+C_{24}g^{\mu \nu }.`$ (B5) where the coefficients are calculated from $`\{C_0,C_{11},C_{12},C_{21},C_{22},C_{23}\}={\displaystyle \frac{g}{16\pi ^2}}{\displaystyle _0^1}𝑑\alpha _1{\displaystyle _0^{1\alpha _1}}𝑑\alpha _2\left({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}}\right)`$ (B6) $`\{1,(1\alpha _1),(1\alpha _1\alpha _2),(1\alpha _1)^2,(1\alpha _1\alpha _2)^2,(1\alpha _1)(1\alpha _1\alpha _2)\}`$ (B7) and the final coefficient is $`C_{24}={\displaystyle \frac{g}{32\pi ^2}}{\displaystyle _0^1}𝑑\alpha _1{\displaystyle _0^{1\alpha _1}}𝑑\alpha _2\mathrm{ln}\left({\displaystyle \frac{F_\mu }{F_{\mathrm{\Lambda }_1}}}\right)`$ (B8) Finally, the Dirac matrices from numerators of two fermion propagators are simplified to standard forms with the assistance of projection operators $`L^{\rho _i}(p_i)`$ and $`L^{\rho _f}(p_f)`$. Once a factor $`p/_i`$ is commuted, if necessary, to act on $`L^{\rho _i}(p_i)`$, it becomes $`\rho _iW_i`$, where $`W_i=\sqrt{p_i^2}`$. Similarly commuting $`p/_f`$ as necessary to act on $`L^{\rho _f}(p_f)`$ results in $`\rho _fW_f`$. The final expressions for form factors are: $`F_1`$ $`{}_{}{}^{\rho _f,\rho _i}(p_f,p_i)=(\rho _fW_f+m)(\rho _iW_i+m)C_0+(\rho _fW_f+m)\rho _iW_iC_{11}`$ (B10) $`+\rho _fW_f(\rho _iW_i+m)C_{11}+\rho _iW_i\rho _fW_fC_{21}(p_ip_f)^2(C_{22}C_{23})2C_{24}`$ $`F_2`$ $`{}_{}{}^{\rho _f,\rho _i}(p_f,p_i)=(\rho _fW_f+m)C_{12}+(\rho _iW_i+m)C_{12}`$ (B12) $`(\rho _iW_i+m)C_{11}\rho _iW_iC_{21}+(\rho _fW_f+\rho _iW_i)C_{23}`$ $`F_3`$ $`{}_{}{}^{\rho _f,\rho _i}(p_f,p_i)=(\rho _fW_f+m)C_{12}+(\rho _iW_i+m)\times (C_{12}C_{11})\rho _iW_iC_{21}`$ (B15) $`+2(\rho _fW_f\rho _iW_i)(C_{22}C_{23})+(\rho _fW_f+\rho _iW_i)C_{23}`$ $`F_5^{\rho _f,\rho _i}(p_f,p_i)=(\rho _fW_f+m)(\rho _iW_i+m)C_0+(\rho _fW_f+m)\rho _iW_iC_{11}`$ (B16) $`(\rho _fW_f+m)(\rho _fW_f+\rho _iW_i)C_{12}\rho _iW_i(\rho _iW_i+m)C_{11}`$ (B17) $`+(\rho _fW_f+\rho _iW_i)(\rho _iW_i+m)C_{12}W_i^2C_{21}(W_f^2W_i^2)C_{23}`$ (B18) $`(p_ip_f)^2[C_{22}C_{23}]4C_{24}`$ (B19) Dependences of the form factors on $`\rho _i`$, $`\rho _f`$ and off-shell momenta are made explicit in these formulas. ## C Expansion in terms of covariants and matrix elements of Born terms It is convenient to expand amplitudes $`V^{\mu 5}`$ and $`V^{5\mu }`$ in terms of kinematical covariants and associated scalar amplitudes. For the kinematical covariants, we use helicity matrix elements of a set of eight Dirac operators, as follows, $$k_1^\mu =\overline{u}_{\lambda _f}(p_f)\gamma _5\gamma ^\mu u_{\lambda _i}(p_i)$$ (C1) $$k_2^\mu =p_i^\mu \overline{u}_{\lambda _f}(p_f)\gamma _5q/u_{\lambda _i}(p_i)$$ (C2) $$k_3^\mu =p_f^\mu \overline{u}_{\lambda _f}(p_f)\gamma _5q/u_{\lambda _i}(p_i)$$ (C3) $$k_4^\mu =q^\mu \overline{u}_{\lambda _f}(p_f)\gamma _5q/u_{\lambda _i}(p_i)$$ (C4) $$k_5^\mu =p_i^\mu \overline{u}_{\lambda _f}(p_f)\gamma _5u_{\lambda _i}(p_i)$$ (C5) $$k_6^\mu =p_f^\mu \overline{u}_{\lambda _f}(p_f)\gamma _5u_{\lambda _i}(p_i)$$ (C6) $$k_7^\mu =q^\mu \overline{u}_{\lambda _f}(p_f)\gamma _5u_{\lambda _i}(p_i)$$ (C7) $$k_8^\mu =\frac{1}{2}\overline{u}_{\lambda _f}(p_f)\gamma _5[\gamma ^\mu ,\gamma ^\nu ]q_\nu u_{\lambda _i}(p_i)$$ (C8) where $`\lambda _i`$ and $`\lambda _f`$ denote the helicities of initial and final states. The direct Born graph of Eq. (37) is expanded as follows, $$V_{\lambda _f\lambda _i}^{\mu ,5}=\underset{n=1}{\overset{8}{}}V_{Dn}k_n^\mu $$ (C9) where the scalar coefficients are $`V_{D1}^{\rho ,+}={\displaystyle \underset{\rho }{}}F_5^{+,\rho }{\displaystyle \frac{1}{D^\rho }}\left[(W_{p_i+q}\rho M)F_1^{\rho ,+}+\rho F_2^{\rho ,+}((p_i+q)^2p_i^2)\right]`$ (C10) $`V_{D2}^{\rho ,+}={\displaystyle \underset{\rho }{}}F_5^{+,\rho }{\displaystyle \frac{1}{D^\rho }}[2\rho F_2^{\rho ,+}]`$ (C11) $`V_{D3}^{\rho ,+}=0`$ (C12) $`V_{D4}^{\rho ,+}={\displaystyle \underset{\rho }{}}F_5^{+,\rho }{\displaystyle \frac{1}{D^\rho }}[\rho (F_3^{\rho ,+}F_2^{\rho ,+})]`$ (C13) $`V_{D5}^{\rho ,+}={\displaystyle \underset{\rho }{}}F_5^{+,\rho }{\displaystyle \frac{1}{D^\rho }}[2\rho F_1^{\rho ,+}]`$ (C14) $`V_{D6}^{\rho ,+}=0`$ (C15) $`V_{D7}^{\rho ,+}=`$ $`{\displaystyle \underset{\rho }{}}F_5^{+,\rho }{\displaystyle \frac{1}{D^\rho }}[W_{p_i+q}F_3^{\rho ,+}+\rho MF_3^{\rho ,+}+\rho F_1^{\rho ,+}]`$ (C16) $`V_{D8}^{\rho ,+}=`$ $`{\displaystyle \underset{\rho }{}}F_5^{+,\rho }{\displaystyle \frac{1}{D^\rho }}[W_{p_i+q}F_2^{\rho ,+}\rho F_1^{\rho ,+}\rho MF_2^{\rho ,+}]`$ (C17) In the $`V_{Dn}`$ expressions, $`D^\rho =2W_{p_i+q}Z_2[1B((p_i+q)^2)\rho W_{p_i+q}A((p_i+q)^2)]`$. Similarly, the cross Born graph of Eq. (36) is expanded as follows, $$V_{\lambda _f\lambda _i}^{5,\mu }=\underset{\rho }{}\underset{n=1}{\overset{8}{}}V_{Xn}k_n$$ (C18) where the scalar coefficients are $`V_{X1}^{+,\rho }={\displaystyle \underset{\rho }{}}\left[(W_{p_fq}\rho M)F_1^{+,\rho }+\rho F_2^{+,\rho }(p_f^2(p_fq)^2)\right]{\displaystyle \frac{1}{D^\rho }}F_5^{\rho ,+}`$ (C19) $`V_{X2}^{+,\rho }=0`$ (C20) $`V_{X3}^{+,\rho }={\displaystyle \underset{\rho }{}}2\rho F_2^{+,\rho }{\displaystyle \frac{1}{D^\rho }}F_5^{\rho ,+}`$ (C21) $`V_{X4}^{+,\rho }={\displaystyle \underset{\rho }{}}\rho (F_2^{+,\rho }+F3^{+,\rho }){\displaystyle \frac{1}{D^\rho }}F_5^{\rho ,+}`$ (C22) $`V_{X5}^{+,\rho }=0`$ (C23) $`V_{X6}^{+,\rho }={\displaystyle \underset{\rho }{}}2\rho F_1^{+,\rho }{\displaystyle \frac{1}{D^\rho }}F_5^{\rho ,+}`$ (C24) $`V_{X7}^{+,\rho }=`$ $`{\displaystyle \underset{\rho }{}}\left[\rho F_1^{+,\rho }+(W_{p_fq}+\rho M)F_3^{+,\rho }\right]{\displaystyle \frac{1}{D^\rho }}F_5^{\rho ,+}`$ (C25) $`V_{X8}^{+,\rho }=`$ $`{\displaystyle \underset{\rho }{}}\left[\rho F_1^{+,\rho }(W_{p_fq}+\rho M)F_2^{+,\rho }\right]{\displaystyle \frac{1}{D^\rho }}F_5^{\rho ,+}`$ (C26) In the $`V_{Xn}`$ expressions, $`D^\rho (p)=2W_{p_fq}Z_2[1B((p_fq)^2)\rho W_{p_fq}A((p_fq)^2)]`$. ## D Contact-like terms Contact-like terms involve three fermion propagators and Dirac matrices $`\gamma _\mu `$ and $`\gamma _5`$. For $`C_{5,\mu }`$, we have $`C^{5\mu }=\overline{u}(p_f)[ig{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}S(p_fk;m)\gamma ^5S(p_i+qk;m)\gamma ^\mu S(p_ik;m)`$ (D1) $`D(k;\mu ,\mathrm{\Lambda }_1)]u(p_i).`$ (D2) We denote the numerator of this expression as $`N^{5\mu }(k)\overline{u}(p_f)[(p/_fk/+m)\gamma ^5(p/_i+q/k/+m)\gamma ^\mu (p/_ik/+m)]u(p_i).`$ (D3) The Feynman parameterization used is $$\frac{1}{d_1d_2d_3d_4}=3![d\alpha ]\frac{1}{[\alpha _1d_1+\alpha _2d_2+\alpha _3d_3+\alpha _4d_4]^4},$$ (D4) where $$[d\alpha ]_0^1𝑑\alpha _1_0^{1\alpha _1}𝑑\alpha _2_0^{1\alpha _1\alpha _2}𝑑\alpha _3,$$ (D5) and $$\alpha _4=1\alpha _1\alpha _2\alpha _3.$$ (D6) This leads to a shift vector $$\mathrm{}=\alpha _4p_f+\alpha _3(p_i+q)+\alpha _2p_i,$$ (D7) and a denominator function, for boson mass $`\mathrm{\Lambda }`$, $`F_\mathrm{\Lambda }=`$ $`\alpha _4[m^2p_f^2]+\alpha _3[m^2(p_i+q)^2]+\alpha _2[m^2p_i^2]+\alpha _1\mathrm{\Lambda }^2+\mathrm{}^2.`$ (D8) The required integration is $`C^{5\mu }=`$ $`{\displaystyle \frac{6g}{16\pi ^2}}{\displaystyle }[d\alpha ]{\displaystyle }{\displaystyle \frac{d^4k}{i\pi ^2}}N^{5\mu }(k)({\displaystyle \frac{1}{[(k\mathrm{})^2F_\mu ]^4}}{\displaystyle \frac{1}{[(k\mathrm{})^2F_{\mathrm{\Lambda }_1}]^4}}).`$ (D9) ¿From the general rules stated above, one sees that a $`k^\mu `$ in the numerator in general is replaced by $`\mathrm{}^\mu `$ after integration over $`k`$, but there are additional contributions from combinations of $`k^\mu k^\nu `$ that involve $`g^{\mu \nu }`$. Therefore we write $`k^\mu =\mathrm{}^\mu +(k\mathrm{})^\mu `$, and expand in powers of $`k\mathrm{}`$. Terms that are odd in $`k\mathrm{}`$ do not contribute because of symmetry. The parts that do contribute are $$N^{5\mu }(k)=N^{5\mu }(\mathrm{})+\mathrm{\Delta }N^{5\mu }(k\mathrm{}),$$ (D10) where $`\mathrm{\Delta }N^{5\mu }(k)=\overline{u}(p_f)[(p/_f\mathrm{}/+m)\gamma ^5k/\gamma ^\mu k/+k/\gamma ^5(p/_i\mathrm{}/+q/+m)\gamma ^\mu k/+k/\gamma ^5k/\gamma ^\mu (p/_i\mathrm{}/+m)]u(p_i).`$ (D11) Integration over k produces $`C_{5\mu }={\displaystyle \frac{g}{16\pi ^2}}{\displaystyle }[d\alpha ]\{N^{5\mu }(\mathrm{})({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}}){\displaystyle \frac{1}{2}}\mathrm{\Delta }N^{5\mu }\times ({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}})\},`$ (D12) where $`\mathrm{\Delta }N^{5\mu }=\overline{u}(p_f)[(p/_f\mathrm{}/+m)\gamma ^5\gamma ^\alpha \gamma ^\mu \gamma ^\beta g_{\alpha \beta }+\gamma ^\alpha \gamma ^5\times (p/_i+q/\mathrm{}/+m)\gamma ^\mu \gamma ^\beta g_{\alpha \beta }`$ (D13) $`+\gamma ^\alpha \gamma ^5\gamma ^\beta g_{\alpha \beta }(p/_i\mathrm{}/+m)]u(p_i).`$ (D14) The resulting expressions are reduced and expanded in terms of scalar amplitudes times kinematical covariants defined in Appendix C as follows. $$C_{\lambda _f\lambda _i}^{5\mu }=\underset{n=1}{\overset{8}{}}C_{Dn}k_n^\mu ,$$ (D15) where $$C_{D1}=[d\alpha ]\left[A_3MA_4MA_6(\alpha _2+\alpha _3+\alpha _4)\right]$$ (D16) $$C_{D2}=[d\alpha ]\left[(\alpha _2+\alpha _3)(A_5+\alpha _3A_7)\right]$$ (D17) $$C_{D3}=[d\alpha ]\left[\alpha _4(A_5+\alpha _3A_7)\right]$$ (D18) $$C_{D4}=[d\alpha ]\left[\alpha _3(A_5+\alpha _3A_7)\right]$$ (D19) $`C_{D5}=`$ $`{\displaystyle [d\alpha ]\left[A_1+(\alpha _2+\alpha _3)(A_2+A_6+A_8)+A_4+(\alpha _2+\alpha _3)(\alpha _4\alpha _2\alpha _3)MA_7\right]}`$ (D20) $`C_{D6}=`$ $`{\displaystyle [d\alpha ]\left[\alpha _4(A_2A_6+\stackrel{~}{A}_8)+\alpha _4(\alpha _4\alpha _2\alpha _3)\times MA_7\right]}`$ (D21) $`C_{D7}=`$ $`{\displaystyle [d\alpha ]\left[A_1+\alpha _3(A_2+A_8)+\alpha _3(\alpha _4\alpha _2\alpha _3)\times MA_7\right]}`$ (D22) $$C_{D8}=[d\alpha ]\left[A_4\alpha _3A_6\right]$$ (D23) and where $$A_1=\left[(M+m)^2\mathrm{}^2\right]\left(\frac{1}{F_\mu ^2}\frac{1}{F_{\mathrm{\Lambda }_1}^2}\right)+2\left(\frac{1}{F_\mu }\frac{1}{F_{\mathrm{\Lambda }_1}}\right)$$ (D24) $`A_2=\left[\mathrm{}^2(M+m)^22M(M+m)\right]\left({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}}\right)3\left({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}}\right)`$ (D25) $`A_3=\left[(M+m)[2\mathrm{}^22\mathrm{}p_{\mathrm{int}}m(M+m)]m\mathrm{}^2\right]\left({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}}\right)`$ (D26) $`\left[3M+2m\right]\left({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}}\right)`$ (D27) $$A_4=\left[(M+m)^2\mathrm{}^2\right]\left(\frac{1}{F_\mu ^2}\frac{1}{F_{\mathrm{\Lambda }_1}^2}\right)$$ (D28) $$A_5=\left[2(M+m)2M(\alpha _2+\alpha _4)\right]\left(\frac{1}{F_\mu ^2}\frac{1}{F_{\mathrm{\Lambda }_1}^2}\right)$$ (D29) $`A_6=\left[m^2M^2\mathrm{}^2+2\mathrm{}p_{\mathrm{int}}\right]\left({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}}\right)+3\left({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}}\right)`$ (D30) $$A_7=\left[2M\right]\left(\frac{1}{F_\mu ^2}\frac{1}{F_{\mathrm{\Lambda }_1}^2}\right)$$ (D31) $`A_8=[2M^2(\alpha _2+\alpha _3+\alpha _4)2p_i(p_i+q)(\alpha _2+\alpha _3)+2p_f(p_i+q)\alpha _4`$ (D32) $`+2p_iq\alpha _3\left]\right({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}})`$ (D33) Here $`p_{\mathrm{int}}=p_i+q`$ is the intermediate momentum, and $`\mathrm{}`$, $`F_\mu `$ and $`F_{\mathrm{\Lambda }_1}`$ are expressed as appropriate for $`C^{5\mu }`$. A very similar analysis is carried out for the crossed contact-like term, $`C^{\mu 5}`$. $`C^{\mu 5}=\overline{u}(p_f)[ig{\displaystyle }{\displaystyle \frac{d^4k}{(2\pi )^4}}S(p_fk;m)\gamma ^\mu S(p_fqk;m)\gamma ^5`$ (D34) $`S(p_ik;m)D(k;\mu ,\mathrm{\Lambda }_1)]u(p_i).`$ (D35) We denote the numerator of this expression as $`N^{\mu 5}\overline{u}(p_f)[(p/_fk/+m)\gamma ^\mu (p/_fq/k/+m)\gamma ^5`$ (D36) $`\times (p/_ik/+m)]u(p_i).`$ (D37) Proceeding as before leads to a shift vector $$\mathrm{}=\alpha _4p_f+\alpha _3(p_fq)+\alpha _2p_i,$$ (D38) and a denominator function, for boson mass $`\mathrm{\Lambda }`$, $`F_\mathrm{\Lambda }=`$ $`\alpha _4[m^2p_f^2]+\alpha _3[m^2(p_fq)^2]+\alpha _2[m^2p_i^2]+\alpha _1\mathrm{\Lambda }^2+\mathrm{}^2.`$ (D39) Numerator factors that produce nonzero results are expressed in a similar way as above, $`N^{\mu 5}(k)=N^{\mu 5}(\mathrm{})+\mathrm{\Delta }N^{\mu 5}(k\mathrm{}),`$ (D40) where $`\mathrm{\Delta }N^{\mu 5}(k)=\overline{u}(p_f)[(p/_f\mathrm{}/+m)\gamma ^\mu k/\gamma ^5k/+k/\gamma ^\mu (p/_f\mathrm{}/q/+m)\gamma ^5k/`$ (D41) $`+k/\gamma ^\mu k/\gamma ^5(p/_i\mathrm{}/+m)]u(p_i).`$ (D42) Integration over k produces $`C_{\mu 5}=`$ $`{\displaystyle \frac{g}{16\pi ^2}}{\displaystyle }[d\alpha ]\{N^{\mu 5}(\mathrm{})({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}}){\displaystyle \frac{1}{2}}\mathrm{\Delta }N^{\mu 5}({\displaystyle \frac{1}{F_\mu }}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}}})\},`$ (D43) where $`\mathrm{\Delta }N^{\mu 5}=\overline{u}(p_f)[(p/_f\mathrm{}/+m)\gamma ^\mu \gamma ^\alpha \gamma ^5\gamma ^\beta g_{\alpha \beta }+\gamma ^\alpha \gamma ^\mu \times (p/_fq/\mathrm{}/+m)\gamma ^5\gamma ^\beta g_{\alpha \beta }`$ (D44) $`+\gamma ^\alpha \gamma ^\mu \gamma ^\beta g_{\alpha \beta }(p/_i\mathrm{}/+m)]u(p_i).`$ (D45) The resulting expressions are expanded in terms of scalar amplitudes times the kinematical covariants, $$C^{\mu 5}=\underset{n=1}{\overset{8}{}}C_{Xn}k_n^\mu ,$$ (D46) where $$+\alpha _3+\alpha _4)]C_{X1}=C_{D1}$$ (D47) $$C_{X2}=[d\alpha ]\left[\alpha _2(A_5+\alpha _3A_7)\right]$$ (D48) $$C_{X3}=[d\alpha ]\left[(\alpha _3+\alpha _4)(A_5+\alpha _3A_7)\right]$$ (D49) $$C_{X4}=C_{D4}$$ (D50) $`C_{X5}=`$ $`{\displaystyle [d\alpha ]\left[\alpha _2(A_2A_6+\stackrel{~}{A}_8)+\alpha _2(\alpha _2\alpha _3\alpha _4)MA_7\right]}`$ (D51) $`C_{X6}=`$ $`{\displaystyle [d\alpha ]\left[A_1+(\alpha _3+\alpha _4)(A_2+A_6+\stackrel{~}{A}_8)+A_4+(\alpha _3+\alpha _4)(\alpha _2\alpha _3\alpha _4)MA_7\right]}`$ (D52) $`C_{X7}=`$ $`{\displaystyle [d\alpha ]\left[A_1\alpha _3(A_2+\stackrel{~}{A}_8)\alpha _3(\alpha _2\alpha _3\alpha _4)MA_7\right]}`$ (D53) $$C_{X8}=C_{D8}$$ (D54) and where the functions $`A_1`$ to $`A_7`$ take the same form as before, except that $`p_{\mathrm{int}}=p_fq`$, and the appropriate $`\mathrm{}`$, $`F_\mu `$ and $`F_{\mathrm{\Lambda }_1}`$ for $`C^{\mu 5}`$ must be used. Also, equalities such as $`C_{X1}=C_{D1}`$ mean that $`C_{X1}`$ takes the same form as $`C_{D1}`$, but of course must be evaluated with the appropriate $`\mathrm{}`$, and so on. Function $`\stackrel{~}{A}_8`$ takes a different form from $`A_8`$, as follows, $`\stackrel{~}{A}_8=[2M^2(\alpha _2+\alpha _3+\alpha _4)+2p_i(p_fq)\alpha _22p_f(p_fq)(\alpha _3+\alpha _4)`$ (D55) $`2p_fq\alpha _3\left]\right({\displaystyle \frac{1}{F_\mu ^2}}{\displaystyle \frac{1}{F_{\mathrm{\Lambda }_1}^2}}).`$ (D56) Calculations have been performed in two ways. One uses the expressions given above and the other uses expressions that have been developed by use of the symbolic manipulation program SCHOONSCHIP in order to reduce the Dirac matrices to the desired forms and FORMF to calculate the moments of the one-loop graphs . Two independent computer codes were written and checked against one another to verify that the algebra and the numerics was done correctly.
warning/0002/hep-ex0002017.html	ar5iv	text	# Measurement of 𝜔 meson parameters in 𝜋⁺⁢𝜋⁻⁢𝜋⁰ decay mode with CMD-2 Work is supported in part by grants RFBR-98-02-17851, RFBR-99-02-17053, RFBR-99-02-17119 ## 1 Introduction High precision measurements of the $`\omega `$ meson parameters provide valuable information for testing various theoretical models describing interactions of light quarks. This paper presents a precise determination of the mass, total width and leptonic width of the $`\omega `$, based on its dominant decay mode, $`\omega \pi ^+\pi ^{}\pi ^0`$. The data sample was collected with the CMD-2 detector in 1994-1995 while scanning the center of mass energy range 2$`E_{beam}`$ from 760 to 810 MeV at the high luminosity collider VEPP-2M . The resonant depolarization method was used for the precise calibration of the beam energy at each point. The integrated luminosity of 141 nb<sup>-1</sup> corresponds to $`7\times 10^4`$ $`\omega `$ meson decays. ## 2 CMD-2 detector The CMD-2 detector has been described in detail elsewhere . It is a general purpose detector consisting of a drift chamber (DC) and proportional Z-chamber (ZC), both inside a thin (0.38 $`X_0`$) superconducting solenoid with a field of 1 T. The barrel calorimeter placed outside the solenoid consists of 892 CsI crystals of 6$`\times `$6$`\times `$15 cm<sup>3</sup> size. It covers polar angles from 0.8 to 2.3 radian. The energy resolution for photons is about 8% in the energy range from 100 to 700 MeV. The trigger signal is generated either by the charged trigger based on DC and ZC hits or by the neutral trigger which takes into account the number and relative position of clusters detected in the CsI calorimeter as well as the total energy deposition. These two independent triggers have been used to study the trigger efficiency. ## 3 Analysis Events with two tracks originating from the same vertex, each with a polar angle $`0.85<\theta <\pi 0.85`$ within the fiducial volume of the detector, were selected for further analysis. To minimize a systematic error of the detection efficiency, only DC information has been used for the selection of $`\omega \pi ^+\pi ^{}\pi ^0`$ events. Most of the background comes from the processes with the hard photon emission: $$e^+e^{}e^+e^{}\gamma ,\pi ^+\pi ^{}\gamma ,\mu ^+\mu ^{}\gamma .$$ These processes have the same signature as the reaction $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$, except for the very different acollinearity angle ($`\mathrm{\Delta }\varphi =\pi \|\phi _1\phi _2\|`$) distribution peaked near $`\mathrm{\Delta }\varphi =0`$. Thus, the rejection of events with a small $`\mathrm{\Delta }\varphi `$ drastically reduces the background, at the same time decreasing the number of $`\pi ^+\pi ^{}\pi ^0`$ events. The value of $`\mathrm{\Delta }\varphi =0.25`$ was used as a reasonable compromise (see Fig. 1-a). Additional background suppression was achieved using the "$`\pi ^+\pi ^{}`$ missing mass" parameter $`M_X`$ which is calculated assuming charged particles to be pions and taking into account energy-momentum conservation. For real $`\pi ^+\pi ^{}\pi ^0`$ events the distribution of the missing mass squared has a peak in the region of $`M_{\pi ^0}^2`$ in contrast to the background processes which have a peak around zero for $`e^+e^{}\pi ^+\pi ^{}(\gamma ),\mu ^+\mu ^{}(\gamma )`$ or in the negative region for $`e^+e^{}e^+e^{}(\gamma )`$. Further rejection of events of the process $`e^+e^{}e^+e^{}\gamma `$ which has the largest cross section among the backround processes, is based on the maximum energy deposition of two charged particles in the calorimeter $`E_{CsI}^{max}`$. The corresponding cuts shown by two lines in Fig. 1-b reject its lower right corner mostly populated by events of this background source. The number of $`\pi ^+\pi ^{}\pi ^0`$ events was obtained in two different ways. The first method was to fit the $`M_X`$ distributions with the sum of Gaussian functions describing $`\pi ^+\pi ^{}\pi ^0`$ and background events. In the second method the cosmic and beam background was rejected by fitting the distribution of the $`z`$-coordinate of the vertex with the sum of a Gaussian function and a constant background. In the last case the remaining $`e^+e^{}e^+e^{}(\gamma ),\pi ^+\pi ^{}(\gamma ),\mu ^+\mu ^{}(\gamma )`$ events were simulated and subtracted from the total number of events at each point according to the corresponding integrated luminosity. Both approaches gave the same result within statistical errors. At each energy the $`\pi ^+\pi ^{}\pi ^0`$ production cross section was calculated according to the formula: $`\sigma (e^+e^{}\pi ^+\pi ^{}\pi ^0)={\displaystyle \frac{N_{\pi ^+\pi ^{}\pi ^0}}{L\epsilon _{trig}\epsilon _{MC}\epsilon _{M_X^2}(1+\delta _{rad})}},`$ where $`N_{\pi ^+\pi ^{}\pi ^0}`$ is the number of events; $`L`$ is the integrated luminosity determined from large angle Bhabha events with the help of the procedure described in ; $`\delta _{rad}`$ is the radiative correction calculated according to with an accuracy better than 0.5%; and $`\epsilon _{trig},\epsilon _{MC},\epsilon _{M_X^2}`$ are respectively the trigger efficiency, the geometrical efficiency (acceptance) multiplied by the reconstruction efficiency, and the efficiency of the cut shown in Fig. 1-b. The acceptance is the probability to detect two pions from the $`\omega `$ decay within a given solid angle. It was calculated by Monte Carlo simulation taking into account radiative photons emitted by initial electrons. The efficiencies $`\epsilon _{MC}`$ and $`\epsilon _{M_X^2}`$ were calculated by Monte Carlo simulation. Their systematic errors were estimated with the help of special "test" events obtained as a result of the constrained fit based on the information from the ZC and CsI calorimeter only. About 40% of the $`\pi ^+\pi ^{}\pi ^0`$ events have two clusters in the CsI calorimeter resulting from a neutral pion decay. Using the polar and azimuthal angles of these clusters as well as the hits of charged tracks in the ZC, one can reconstruct the $`\omega \pi ^+\pi ^{}\pi ^0`$ event without DC information. “Test” events with the neutral trigger were also used to determine the charged trigger efficiency. Typical values of the efficiencies and corrections are presented in Table 1 for $`2E_{beam}=782.0`$ MeV (the $`\omega `$ meson peak). The integrated luminosity, radiative correction, number of selected $`\pi ^+\pi ^{}\pi ^0`$ events and cross section for $`e^+e^{}\omega \pi ^+\pi ^{}\pi ^0`$ at each energy are presented in Table 2. ## 4 $`\omega `$ meson parameters The experimental data were fitted with a function which includes the interference of the $`\omega `$ and $`\varphi `$ mesons and non-resonant background: $$\sigma _{3\pi }(s)=\frac{F_{3\pi }(s)}{s^{3/2}}\|A_\omega +e^{i\alpha }A_\varphi +A_{bg}\|^2,$$ (1) $$A_V=\frac{m_V^2\mathrm{\Gamma }_V\sqrt{\sigma _Vm_V/F_{3\pi }(m_V^2)}}{sm_V^2+i\sqrt{s}\mathrm{\Gamma }_V(s)},$$ $$A_{bg}=m_\omega ^{3/2}\sqrt{\sigma _{bg}/F_{3\pi }(m_\omega ^2)},$$ $$\mathrm{\Gamma }_\omega (s)=\mathrm{\Gamma }_\omega \left(Br_{\pi ^+\pi ^{}}\frac{m_\omega ^2F_{2\pi }(s)}{sF_{2\pi }(m_\omega ^2)}+Br_{\pi ^0\gamma }\frac{F_{\pi ^0\gamma }(s)}{F_{\pi ^0\gamma }(m_\omega ^2)}+Br_{3\pi }\frac{\sqrt{s}F_{3\pi }(s)}{m_\omega F_{3\pi }(m_\omega ^2)}\right),$$ $$F_{\pi ^0\gamma }(s)=(\sqrt{s}(1m_{\pi ^0}^2/s))^3,F_{2\pi }(s)=(s/4m_\pi ^2)^{3/2},$$ where $`m_V,\mathrm{\Gamma }_V,\sigma _V`$ are mass, width and peak cross section ($`s=m_V^2`$) for the vector meson $`\omega `$ or $`\varphi `$; $`\alpha `$ is a relative phase of $`\omega \varphi `$ mixing taken to be $`(155\pm 15)^{}`$ according to . $`F_{3\pi }(s)`$ is a smooth function which describes the dynamics of $`V\rho \pi \pi ^+\pi ^{}\pi ^0`$ decay including the phase space . $`\mathrm{\Gamma }_\varphi (s)`$ has been parametrized similarly to $`\omega `$ using the corresponding branching ratios and phase space factors . The cross section values were fit by the function (1). The $`\omega `$ meson mass, width, peak cross section and background cross section were optimized, while the $`\varphi `$ meson parameters were fixed at their world average values . The energy dependence of the cross section is shown in Fig. 3 (experimental points and the optimal fitting curve). The following $`\omega `$ meson parameters were obtained from the fit: $`\sigma _0=(1457\pm 23)`$ nb, $`M_\omega =(782.71\pm 0.07)`$ MeV/c<sup>2</sup>, $`\mathrm{\Gamma }_\omega =(8.68\pm 0.23)`$ MeV, $`\sigma _{bg}=(12\pm 5)`$ nb. The systematic error of $`\sigma _0`$ is about 1.3% and comes from the following sources: reconstruction efficiency 0.5% ; trigger efficiency 0.1% ; radiative corrections for the process $`e^+e^{}\pi ^+\pi ^{}\pi ^0`$ 0.5% ; decays in flight 0.1% ; pion nuclear interaction 0.2% ; solid angle uncertainty 0.3% ; luminosity determination 1.0% . The systematic error of the mass was found to be about 40 keV dominated by the stability of the beam energy. The systematic error of the width was found to be about 100 keV dominated by the scatter of results of various fits corresponding to different selection criteria. ## 5 Discussion The measurements of $`\sigma _0(\omega \pi ^+\pi ^{}\pi ^0)`$ have been performed by a number of groups from Orsay and Novosibirsk with the results presented in Table 3. One can see that the result of this work $`\sigma _0(\omega \pi ^+\pi ^{}\pi ^0)=(1457\pm 30)`$ nb does not contradict these measurements and is the most precise. The cross section in the peak obtained in our experiment is related to the product $`\mathrm{\Gamma }_{e^+e^{}}`$Br$`(\omega \pi ^+\pi ^{}\pi ^0)`$. To obtain this value, the fit with this product as a free parameter has been performed with the following result: $`\mathrm{\Gamma }_{e^+e^{}}\text{Br}(\omega \pi ^+\pi ^{}\pi ^0)=(0.528\pm 0.012\pm 0.007)\text{keV},`$ which is the most precise direct measurement of this quantity. Using $`\mathrm{\Gamma }_{e^+e^{}}`$ from other experiments, one can obtain $`Br(\omega \pi ^+\pi ^{}\pi ^0)`$. For example, for $`\mathrm{\Gamma }_{e^+e^{}}=(0.60\pm 0.02)`$ keV from , $`Br(\omega \pi ^+\pi ^{}\pi ^0)=0.880\pm 0.020\pm 0.032`$ can be obtained. Alternatively, taking $`Br(\omega \pi ^+\pi ^{}\pi ^0)`$ from other measurements, $`\mathrm{\Gamma }_{e^+e^{}}`$ can be calculated. For $`Br(\omega \pi ^+\pi ^{}\pi ^0)=0.888\pm 0.007`$ (from ), we obtain for the leptonic width $`\mathrm{\Gamma }_{e^+e^{}}=(0.595\pm 0.014\pm 0.009)`$ keV or for the leptonic branching ratio $`\mathrm{\Gamma }_{e^+e^{}}/\mathrm{\Gamma }_\omega =(6.85\pm 0.11\pm 0.11)10^5`$. In Fig. 3 the result of this work (CMD95) is compared to the previous measurements of the $`\omega `$ meson mass. The left shaded bar corresponds to the current world average. Its value is dominated by the CMD87 experiment which was also performed at the VEPP-2M collider with the CMD detector . The reported precision of CMD87 was much better than in all other experiments. Our new measurement gives the $`\omega `$ meson mass value 930 keV higher (more than seven standard deviations) than in CMD87. Since both measurements were performed at VEPP-2M and used the resonant depolarization method (RDM), thorough comparison of two experiments has been carried out. We now assume that the difference between the two results is due to the fact that resonant depolarization method (RDM) measurements in CMD87 were regularly performed at some side band resonance. Such a resonance could arise from a parasitic modulation of the depolarizer frequency since the RF device used for the RDM had the power about $`10^4`$$`10^5`$ times higher than required. Thus, the absolute calibration of the beam energy gave wrong results. Unfortunately, after a lapse of more than ten years, it is impossible to reproduce the CMD87 environment and prove the above hypothesis. However, we know that because of various technical problems inadequate attention was given at that time to the possibility of the low modulation leading to a side band resonance. Since the time of the CMD87 experiment, the VEPP-2M collider and the RDM hardware have been upgraded. The applied RF power is of the order of a few mW excluding any “parasitic” depolarization. Furthermore, the frequency spectrum of the depolarizer was investigated before the RDM measurements and we believe that in the present experiment the sources of the systematic error in RDM considered above have been completely removed. There were also some other differences between the RDM measurements in both experiments. In CMD87 the beam was polarized in the VEPP-2M ring itself at the beam energy of about 700 MeV. This led to variations of the collider parameters before each RDM measurement including the change of the betatron frequencies $`\nu _x,\nu _z`$ in order to pass through intrinsic spin resonances. The imperfection resonance at the “magic energy” $`E_{beam}=440.65`$ MeV was crossed adiabatically by decompensating the longitudinal magnetic field of the detector (so called “partial siberian snake” mode ). Another parameter affecting the beam energy was the collider temperature which changed by approximately 10 C between the polarization at the high energy and a subsequent RDM measurement. In the present experiment the beam was polarized in the new booster ring at the high energy and after that injected into VEPP-2M at the energy of the experiment. Thus, the parameters of the collider itself were not changed and RDM measurements were performed under the same conditions as data taking. The beam energy stability during data taking has been thoroughly analysed. This analysis was based both on the deviations of about 60 RDM measurements at different energies from the predicted values and on direct measurements of the beam energy stability by the tracking system of the CMD-2 detector . The RDM measurements were consistent with each other in the energy range covering the $`\omega `$ and $`\varphi `$ mesons and showed the long-term beam energy instability of the order of 50 keV. The latter was used as an uncertainty of the beam energy at each point for the calculation of the $`\omega `$ meson mass systematic error. The value of the $`\omega `$ meson mass obtained in this work is close to the world average before the CMD87 experiment $`M_\omega =(782.55\pm 0.17)`$ MeV/$`c^2`$ (the right shaded bar in Fig. 3) and is the most precise today. The results of this work on the total width of the $`\omega `$ meson (Fig. 4-a) as well as on the leptonic branching ratio $`\mathrm{\Gamma }_{e^+e^{}}/\mathrm{\Gamma }_\omega `$ (Fig. 4-b) are in good agreement with those from previous experiments. ## 6 Conclusion Using the CMD-2 data sample of about 11 200 $`\omega \pi ^+\pi ^{}\pi ^0`$ events, the following values of the $`\omega `$ meson parameters have been obtained: $`\sigma _0`$ $`=`$ $`(1457\pm 23\pm 19)\text{nb},`$ $`M_\omega `$ $`=`$ $`(782.71\pm 0.07\pm 0.04)\text{MeV}/c^2,`$ $`\mathrm{\Gamma }_\omega `$ $`=`$ $`(8.68\pm 0.23\pm 0.10)\text{MeV},`$ $`\mathrm{\Gamma }_{e^+e^{}}Br(\omega \pi ^+\pi ^{}\pi ^0)`$ $`=`$ $`(0.528\pm 0.012\pm 0.007)10^3\text{MeV}.`$ These results, except for the total width, are more precise than those from previous experiments. The mass value differs significantly from the previous most precise measurement which was performed by a group including many authors of this work. Due to the present more thorough study of systematic errors, our mass measurement supersedes that of Ref. . ## 7 Acknowledgements The authors are grateful to the staff of VEPP-2M for the excellent performance of the collider and to all engineers and technicians who participated in the design, commissioning and operation of CMD-2.
warning/0002/hep-th0002235.html	ar5iv	text	# Aspects of Holography and Rotating AdS Black Holes ## 1 Introduction The AdS/CFT correspondence has provided many insights into the connection between gauge theories and gravity. (See for an extensive review of this vast subject.) In some sense the correspondence provides a duality map between large N Yang-Mills and IIB supergravity in an $`AdS_5\times S^5`$ background. The origin of this duality relation is from considering the physics of D3 branes from two perspectives. The world volume description, as can be derived from considering the D-brane as a submanifold on which strings can end, gives a Yang-Mills theory. The soliton description, whereby the the D-brane is seen as a solution in IIB supergravity provides, after taking appropriate limits, the alternative $`AdS`$ supergravity point of view. One may use this duality relation to calculate quantities in the Yang-Mills theory at large ‘t Hooft coupling. In particular one may calculate thermodynamic quantities in the Yang-Mills theory from considering black holes in the AdS space-time . Comparing these thermodynamic quantities with those calculated directly from the Yang-Mills theory at weak coupling, one sees that they correspond up to a factor of four thirds. (It should be stressed that there is no reason to expect them to match as typically entropies change as the coupling is altered). For nonrotating AdS black holes , the thermodynamics has been described by thermal conformal field theory . Recently a five-dimensional rotating black hole embedded in anti-de Sitter space has been discovered . Since rotation introduces an extra dimensionful parameter, the conformal field theory entropy is not so tightly constrained by the combination of extensivity and dimensional analysis; thus a comparison between thermodynamic quantities is much more nontrivial. Our purpose in this paper, as first reported in , is to compare the thermodynamics of the new rotating black hole with that of the dual conformal field theory in four dimensions. If one then assumes the AdS/CFT correspondence to be correct then the result may be interpreted as the comparison of thermodynamic quantites in the field theory at strong and weak coupling. This correspondence has also been investigated in the case of rotating D branes by . We begin by demonstrating the holographic nature of the duality for nonrotating black holes: the thermodynamics of a nonrotating black hole in anti-de Sitter space emerges from a thermal conformal field theory whose thermodynamic variables are read off from the boundary of the black hole spacetime. In the high temperature limit, the field theory calculation gives the correct entropy of the Hawking-Page black hole up to a factor of 4/3. We then describe the new rotating Kerr-AdS black hole solution and show how its thermodynamic properties can be recovered from the dual field theory, in the high temperature limit. In that limit, the entropy, energy and angular momentum, compared with the statistical mechanics of the field theory, all agree with their gravitational counterparts, again up to a common factor of 4/3. ## 2 AdS/CFT Correspondence for Nonrotating Holes The five-dimensional Einstein-Hilbert action with a cosmological constant is given by $$I=\frac{1}{16\pi G_5}d^5x\sqrt{g}\left(R+12l^2\right),$$ (1) where $`G_5`$ is the five-dimensional Newton constant, $`R`$ is the Ricci scalar, the cosmological constant is $`\mathrm{\Lambda }=6l^2`$, and we have neglected a surface term at infinity. Anti-de Sitter solutions derived from this action can be embedded in ten-dimensional IIB supergravity such that the supergravity background is of the form $`AdS_5\times S^5`$. The AdS/CFT correspondence then states that there is a dual conformal field theory in four dimensions from which one can extract the physics. The line element of a “Schwarzschild” black hole in anti-de Sitter space in five spacetime dimensions can be written as $`ds^2=\left(1{\displaystyle \frac{2MG_5}{r^2}}+r^2l^2\right)dt^2`$ $`+\left(1{\displaystyle \frac{2MG_5}{r^2}}+r^2l^2\right)^1dr^2+r^2d\mathrm{\Omega }_3^2.`$ (2) This solution has a horizon at $`r=r_+`$ where $$r_+^2=\frac{1}{2l^2}\left(1+\sqrt{1+8MG_5l^2}\right).$$ (3) The substitution $`\tau =it`$ makes the metric positive definite and, by the usual removal of the conical singularity at $`r_+`$, yields a periodicity in $`\tau `$ of $$\beta =\frac{2\pi r_+}{1+2r_+^2l^2},$$ (4) which is identified with the inverse temperature of the black hole. The entropy is given by $$S=\frac{A}{4G_5}=\frac{\pi ^2r_+^3}{2G_5},$$ (5) where $`A`$ is the “area” (that is 3-volume) of the horizon. We shall take the dual conformal field theory to be $`𝒩=4`$, U(N) super-Yang-Mills theory. But since it is only possible to do calculations in the weak coupling regime, we shall consider only the free field limit of Yang-Mills theory. Then, in the high-energy regime which dominates the state counting, the spectrum of free fields on a sphere is essentially that of blackbody radiation in flat space, with $`8N^2`$ bosonic and $`8N^2`$ fermionic degrees of freedom. The entropy is therefore $$S_{\mathrm{CFT}}=\frac{2}{3}\pi ^2N^2V_{\mathrm{CFT}}T_{\mathrm{CFT}}^3.$$ (6) We would like to evaluate this “holographically”, i.e. by substituting physical data taken from the boundary of the black hole spacetime. At fixed $`rr_0r_+`$, the boundary line element tends to $$ds^2r_0^2\left[l^2dt^2+d\mathrm{\Omega }_3^2\right].$$ (7) The physical temperature at the boundary is consequently red-shifted to $$T_{\mathrm{CFT}}=\frac{T_{BH}}{\sqrt{g_{tt}}}=\frac{T_{BH}}{lr_0},$$ (8) while the volume is $$V_{\mathrm{CFT}}=2\pi ^2r_0^3.$$ (9) To obtain an expression for $`N`$, we invoke the AdS/CFT correspondence. Originating in the near horizon geometry of the D3-brane solution in IIB supergravity, the correspondence , relates $`N`$ to the radius of $`S^5`$ and the cosmological constant: $$R_{S^5}^2=\sqrt{4\pi g_s\alpha _{}^{}{}_{}{}^{2}N}=\frac{1}{l^2}.$$ (10) Then, since $$(2\pi )^7g_s^2\alpha _{}^{}{}_{}{}^{4}=16\pi G_{10}=16\frac{\pi ^4}{l^5}G_5,$$ (11) we have $$N^2=\frac{\pi }{2l^3G_5}.$$ (12) Substituting the expressions for $`N`$, $`V_{\mathrm{CFT}}`$ and $`T_{\mathrm{CFT}}`$ into Eq. (6), we obtain $$S_{\mathrm{CFT}}=\frac{1}{12}\frac{\pi ^2}{l^6G_5}\left(\frac{1+2r_+^2l^2}{r_+}\right)^3,$$ (13) which, in the limit $`r_+l1`$, reduces to $$S_{\mathrm{CFT}}=\frac{2}{3}\frac{\pi ^2r_+^3}{G_5}=\frac{4}{3}S_{\mathrm{BH}},$$ (14) in agreement with the black hole result, Eq. (5), but for a numerical factor of 4/3. Similarly, the red-shifted energy of the conformal field theory matches the black hole mass, modulo a coefficient. The mass above the anti-de Sitter background is $$M^{}=\frac{3\pi }{4}M.$$ (15) This is the AdS equivalent of the ADM mass, or energy-at-infinity. The corresponding expression in the field theory is $$U_{\mathrm{CFT}}^{\mathrm{}}=\sqrt{g_{tt}}\frac{\pi ^2}{2}N^2V_{\mathrm{CFT}}T_{\mathrm{CFT}}^4=\frac{\pi }{2}r_+^4l^2=\frac{4}{3}M^{},$$ (16) where $`U_{\mathrm{CFT}}^{\mathrm{}}`$ is the conformal field theory energy red-shifted to infinity, and we have again taken the $`r_+l1`$ limit. The $`4/3`$ discrepancy in Eqs. (5) and (16) is construed to be an artifact of having calculated the gauge theory entropy in the free field limit rather than in the strong coupling limit required by the correspondence; intuitively, one expects the free energy to decrease when the coupling increases. The 4/3 factor was first noticed in the context of D3-brane thermodynamics . Our approach differs in that we take the idea of holography at face value, by explicitly reading physical data from the boundary of spacetime; nonetheless, Eq. (12) refers to an underlying brane solution. At this level, the correspondence only goes through in the $`r_+l>>1`$ limit. Note that, in terms of the conformal field theory $`r_+l=T_{\mathrm{CFT}}r_0`$. The limit we have taken means that $`T_{\mathrm{CFT}}1/r_0`$, allowing us to neglect finite-size effects in the field theory which we have implicitly done in calculating the entropy in (6). ## 3 Five-Dimensional Rotating AdS Black Holes The general rotating black hole in five dimensions has two independent angular momenta. Here we consider the case of a rotating black hole with one angular momentum in an ambient AdS space. The line element is $`ds^2`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }}{\rho ^2}}\left(dt{\displaystyle \frac{a\mathrm{sin}^2\theta }{\mathrm{\Xi }_a}}d\varphi \right)^2`$ (17) $`+`$ $`{\displaystyle \frac{\mathrm{\Delta }_\theta \mathrm{sin}^2\theta }{\rho ^2}}\left(adt{\displaystyle \frac{\left(r^2+a^2\right)}{\mathrm{\Xi }}}d\varphi \right)^2`$ $`+{\displaystyle \frac{\rho ^2}{\mathrm{\Delta }}}dr^2+{\displaystyle \frac{\rho ^2}{\mathrm{\Delta }_\theta }}d\theta ^2+r^2\mathrm{cos}^2\theta d\psi ^2,`$ where $`0\varphi ,\psi 2\pi `$ and $`0\theta \pi /2`$, and $`\mathrm{\Delta }`$ $`=`$ $`\left(r^2+a^2\right)\left(1+r^2l^2\right)2MG_5`$ $`\mathrm{\Delta }_\theta `$ $`=`$ $`1a^2l^2\mathrm{cos}^2\theta `$ $`\rho ^2`$ $`=`$ $`r^2+a^2\mathrm{cos}^2\theta `$ $`\mathrm{\Xi }`$ $`=`$ $`1a^2l^2.`$ (18) This solution is an anti-de Sitter space with curvature given by $$R_{ab}=4l^2g_{ab}.$$ (19) The horizon is at $$r_+^2=\frac{1}{2l^2}\left((1+a^2l^2)+\sqrt{(1+a^2l^2)^2+8MG_5l^2}\right),$$ (20) which can be inverted to give $$MG_5=\frac{1}{2}(r_+^2+a^2)(1+r_+^2l^2).$$ (21) The entropy is one-fourth the “area” of the horizon: $$S=\frac{1}{2G_5}\frac{\pi ^2\left(r_+^2+a^2\right)r_+}{\left(1a^2l^2\right)}.$$ (22) The entropy diverges in two different limits: $`r_+\mathrm{}`$ and $`a^2l^21`$. The first of these descibes an infinite temperature and infinite radius black hole, while the second corresponds to “critical angular velocity”, at which the Einstein universe at infinity has to rotate at the speed of light. The inverse Hawking temperature is $$\beta =\frac{2\pi \left(r_+^2+a^2\right)}{r_+\left(1+a^2l^2+2r_+^2l^2\right)}.$$ (23) The mass above the anti-de Sitter background is now $$M^{}=\frac{3\pi }{4\mathrm{\Xi }}M,$$ (24) the angular velocity at the horizon is $$\mathrm{\Omega }_H=\frac{a\mathrm{\Xi }}{r_+^2+a^2},$$ (25) and the angular momentum is defined as $$J_\varphi =\frac{1}{16\pi }_Sϵ_{abcde}^d\psi ^edS^{abc}=\frac{\pi Ma}{2\mathrm{\Xi }^2},$$ (26) where $`\psi ^a=\left(\frac{}{\varphi }\right)^a`$ is the Killing vector conjugate to the angular momentum in the $`\varphi `$ direction, and S is the boundary of a hypersurface normal to $`\left(\frac{}{t}\right)^a`$, with $`dS^{abc}`$ being the volume element on S. Following methods discussed in , one can derive a finite action for this solution from the regularized spacetime volume after an appropriate matching of hypersurfaces at large $`r`$. The result is $$I=\frac{\pi ^2\left(r_+^2+a^2\right)^2(1r_+^2l^2)}{4G_5\mathrm{\Xi }r_+(1+a^2l^2+2r_+^2l^2)}.$$ (27) As noted in , the action changes sign at $`r_+l=1`$, signalling the presence of a phase transition in the conformal field theory. For $`r_+l>1`$, the theory is in an unconfined phase and has a free energy proportional to $`N^2`$. One can also check that the action satisfies the thermodynamic relation $$S=\beta (M^{}J_\varphi \mathrm{\Omega }_H)I.$$ (28) It is interesting to note that, by formally dividing both the free energy, $`F=I/\beta `$, and the mass by an arbitrary volume, one obtains an equation of state: $$p=\frac{1}{3}\frac{r_+^2l^21}{r_+^2l^2+1}\rho ,$$ (29) where $`p=F/V`$ is the pressure, and $`\rho `$ is the energy density. In the limit $`r_+l1`$ that we have been taking, this equation becomes $$p=\frac{1}{3}\rho ,$$ (30) as is appropriate for the equation of state of a conformal theory. This suggests that if a conformal field theory is to reproduce the thermodynamic properties of this gravitational solution, it has to be in such a limit. ## 4 The dual CFT description The gauge theory dual to supergravity on $`AdS_5\times S^5`$ is $`𝒩=4`$ super Yang-Mills with gauge group $`U(N)`$ where $`N`$ tends to infinity . The action is $`S={\displaystyle }d^4x\sqrt{g}\mathrm{Tr}({\displaystyle \frac{1}{4g^2}}F^2+{\displaystyle \frac{1}{2}}\left(D\mathrm{\Phi }\right)^2`$ $`+{\displaystyle \frac{1}{12}}R\mathrm{\Phi }^2+\overline{\psi }/D\psi ).`$ (31) All fields take values in the adjoint representation of U(N). The six scalars, $`\mathrm{\Phi }`$, transform under $`SO(6)`$ R-symmetry, while the four Weyl fermions, $`\psi `$, transform under $`SU(4)`$, the spin cover of $`SO(6)`$. The scalars are conformally coupled; otherwise, all fields are massless. We shall again take the free field limit. The angular momentum operators can be computed from the relevant components of the stress energy tensor in spherical coordinates. This approach is to be contrasted with in which generators of R-rotations are used corresponding to spinning D-branes. The free energy of the gauge theory is given by $`F_{\mathrm{CFT}}=+T_{\mathrm{CFT}}{\displaystyle \underset{i}{}}\eta _i{\displaystyle _0^{\mathrm{}}}𝑑l_i{\displaystyle 𝑑m_i^\varphi 𝑑m_i^\psi }`$ $`\mathrm{ln}\left(1\eta _ie^{\beta (\omega _im_i^\varphi \mathrm{\Omega }_\varphi )}\right),`$ (32) where $`i`$ labels the particle species, $`\eta =+1`$ for bosons and -1 for fermions, $`l_i`$ is the quantum number associated with the total orbital angular momentum of the ith particle, and $`m_i^{\varphi (\psi )}`$ is its angular momentum component in the $`\varphi (\psi )`$ direction. Here $`\mathrm{\Omega }`$ plays the role of a “voltage” while the “chemical potential” $`m^\varphi \mathrm{\Omega }`$ serves to constrain the total angular momentum of the system. The free energy is easiest to evaluate in a corotating frame, which corresponds to the constant-time foliation choice of hypersurfaces orthogonal to $`t^a`$. Since, at constant $`r=r_0`$, the boundary has the metric $`ds^2`$ $`=`$ $`r_0^2[l^2dt^2+{\displaystyle \frac{2al^2\mathrm{sin}^2\theta }{\mathrm{\Xi }}}dtd\varphi +{\displaystyle \frac{\mathrm{sin}^2\theta }{\mathrm{\Xi }}}d\varphi ^2`$ (33) $`+`$ $`{\displaystyle \frac{d\theta ^2}{\mathrm{\Delta }_\theta }}+\mathrm{cos}^2\theta d\psi ^2],`$ the constant-time slices of the corotating frame have a spatial volume of $$V=\frac{2\pi ^2r_0^3}{1a^2l^2}.$$ (34) The spectrum of a conformally coupled field on $`S^3`$ is essentially given by $$\omega _l\frac{l}{r_0},$$ (35) where $`l`$ is the quantum number for total orbital angular momentum. Eq. (32) can now be evaluated by making use of the identities $`{\displaystyle _0^{\mathrm{}}}𝑑xx^n\mathrm{ln}\left(1e^{x+c}\right)`$ $`=`$ $`\mathrm{\Gamma }(n+1)\mathrm{Li}_{n+2}(e^c)`$ $`=`$ $`\mathrm{\Gamma }(n+1){\displaystyle \underset{k=1}{\overset{k=\mathrm{}}{}}}{\displaystyle \frac{e^{kc}}{k^{n+2}}},`$ $`{\displaystyle 𝑑xx\mathrm{Li}_2(e^{ax+c})}`$ $`=`$ $`{\displaystyle \frac{1}{a^2}}[ax\mathrm{Li}_3(e^{ax+c})`$ (36) $`+`$ $`\mathrm{Li}_4(e^{ax+c})],`$ where $`\mathrm{Li}_n`$ is the nth polylogarithmic function, defined by the sum above. The result is $$F_{\mathrm{CFT}}=\frac{\pi ^4}{24}\frac{r_0}{\frac{1}{r_0^2}\mathrm{\Omega }^2}(8N^2)T_{\mathrm{CFT}}^4,$$ (37) yielding an entropy of $$S_{\mathrm{CFT}}=\frac{2}{3}\frac{\pi ^5}{l^3G_5}\frac{r_0^3}{1\mathrm{\Omega }^2r_0^2}T_{\mathrm{CFT}}^3.$$ (38) The physical temperature that enters the conformal field theory is $$T_{\mathrm{CFT}}=\frac{1}{lr_0}T_{\mathrm{BH}}.$$ (39) Similarly, the angular velocity is scaled to $$\mathrm{\Omega }_{\mathrm{CFT}}=\frac{al^2}{lr_0}.$$ (40) Substituting Eqs. (39) and (40) into Eq. (38) and taking the high temperature limit as before, we have $$S_{\mathrm{CFT}}=\frac{2}{3G_5}\frac{\pi ^2r_+^3}{(1a^2l^2)}=\frac{4}{3}S_{\mathrm{BH}}.$$ (41) The inclusion of rotation evidently does not affect the ratio of the black hole and field theory entropies. In the corotating frame, the free energy is simply of the form $`N^2VT^4`$, with the volume given by Eq. (34). However, with respect to a nonrotating AdS space, the free energy takes a more complicated form since now the volume is simply $`2\pi ^2r_0^3`$. By keeping this volume and the temperature fixed, one may calculate the angular momentum of the system with respect to the nonrotating background: $$J_\varphi ^{\mathrm{CFT}}=\frac{F}{\mathrm{\Omega }}\|_{V,T_{\mathrm{CFT}}}=\frac{ar_+^4\pi \left(1+a^2l^2+2r_+^2l^2\right)^4}{48l^6\mathrm{\Xi }^2\left(r_+^2+a^2\right)^4}.$$ (42) In the usual $`r_+l1`$ limit, we obtain $$J_\varphi ^{\mathrm{CFT}}=\frac{2\pi Ma}{3\mathrm{\Xi }^2}=\frac{4}{3}J_\varphi ^{\mathrm{BH}},$$ (43) so that the gauge theory angular momentum is proportional to the black hole angular momentum, Eq. (26), with a factor of 4/3. The black hole mass formula, Eq. (24), refers to the energy above the nonrotating anti-de Sitter background. We should therefore compare this quantity with the red-shifted energy in the conformal field theory. Here a slight subtlety enters. Since the statistical mechanical calculation gives the energy in the corotating frame, we must add the center-of-mass rotational energy before comparing with the black hole mass. Then we find that $$U_{\mathrm{CFT}}^{\mathrm{}}=\sqrt{g_{tt}}\left(U_{\mathrm{corotating}}+J_{\mathrm{CFT}}\mathrm{\Omega }_{\mathrm{CFT}}\right)=\frac{4}{3}M^{},$$ (44) with $`M^{}`$ given by Eq. (24), evaluated at high temperature. Using $`U_{\mathrm{CFT}}^{\mathrm{}}=\sqrt{g_{tt}}U_{\mathrm{CFT}}`$ and previous expressions for thermodynamic quantities, one may check that the first law of thermodynamics is satisfied. ## 5 Discussion There are several interesting aspects of these results. The first is that the same relative factor that appears in the entropy appears in the angular momentum and the energy. A priori, one has no reason to believe that the functional form of the free energy will be such as to guarantee this result (see, for example,). The second is that the relative factor of 4/3 in the entropy is unaffected by rotation. Indeed, one could expand the entropy of the rotating system in powers and inverse powers of the ’t Hooft coupling. The correspondence implies that $$S_{\mathrm{CFT}}=\underset{m}{}a_m\lambda ^m=\underset{n}{}b_n\left(\frac{1}{\sqrt{\lambda }}\right)^n=S_{\mathrm{BH}}.$$ (45) We may approximate the series on the gauge theory side as $`a_0`$ and on the gravity side as $`b_0`$. Then, generically, we would expect these coefficients to be functions of the dimensionless rotational parameter $`\mathrm{\Xi }`$ so that $`a_0(\mathrm{\Xi })=f(\mathrm{\Xi })b_0(\mathrm{\Xi })`$ with $`f(\mathrm{\Xi }=1)=4/3`$. Our somewhat unexpected result is that $`f(\mathrm{\Xi })=4/3`$ has, in fact, no dependence on $`\mathrm{\Xi }`$. Ofcourse, in some sense the most fascinating aspect is that the ratio between strong and weak coupling is such a simple rational number. A futher understanding of this property of SYM would be very desirable. Similar properties have been investigated with the incorporation of finite size effects by . ## 6 Acknowledgments We would like to thank José Barbón, Roberto Emparan, and Kostas Skenderis for helpful discussions. D. B. is supported by European Commission TMR programme ERBFMRX-CT96-0045. M. P. is supported by the Netherlands Organization for Scientific Research (NWO).
warning/0002/astro-ph0002312.html	ar5iv	text	# Gamma-Ray Bursts via the Neutrino Emission from Heated Neutron Stars ## 1 Introduction Understanding the origin of gamma-ray bursts (GRBs) has been a perplexing problem since they were first detected almost three decades ago (Klebesadel et al., 1973). The fact that GRBs are distributed isotropically (Meegan et al., 1992) suggests a cosmological origin. Furthermore, arcminute burst locations from BeppoSax have revealed that at least some $`\gamma `$-ray bursts involve weak X-ray, optical, or radio transients, and are of cosmological origin (Groot et al., 1997). The Mg I absorption and \[O II\] emission lines along the line of sight from the GRB970508 optical transient, for example, indicate a redshift $`Z0.835`$ (Galama et al., 1997). The implied distance means that this burst must have released of order $`{}_{>}{}^{}10_{}^{51}`$ ergs in $`\gamma `$-rays on a time scale $``$ seconds. This energy requirement has been rendered even more demanding by other events such as GRB971214 (Kulkarni et al., 1998) which appears to be centered on a galaxy at redshift 3.42. This implies that the energy of a $`4\pi `$ burst would have to be as much as $`3\times 10^{53}`$ ergs, comparable to the visible light output of $`10^9`$ galaxies. Based upon the accumulated evidence one can can now conclude that the following four features probably characterize the source environment: 1) If the total burst energies are in the range of $`10^{51}10^{52}`$ ergs, then a beaming factor of 10 to 100 is necessary; 2) The multiple peak temporal structure of most bursts probably requires either multiple colliding shocks (Rees & Mészáros, 1994; Kobayashi et al., 1998) or a single shock impinging upon a clumpy interstellar medium (Mészáros & Rees, 1993; Dermer & Mitman, 1999); 3) The observed afterglows imply some surrounding material on a scale of light hours; and 4) the presence of \[O II\] emission lines suggests that the bursts occur in a young, metal-enriched stellar population. Some proposed sources for the production of GRBs include accretion onto supermassive black holes, AGN’s, relativistic stellar collisions, hypernovae, and binary neutron star coalescence. Each of these possibilities, however, remain speculative until realistic models can be constructed for their evolution. In this paper we construct a model for GRBs produced by energetic neutrino emission from a heated neutron star. Our specific model for the emission derives from the relativistic compression and heating of neutron stars near their last stable orbit, however any scenario by which energetic neutrino emission above a neutron star can endure for several seconds (e.g. tidal heating, MHD induced heating, accretion shocks, etc) might also power the gamma-ray burst paradigm described herein. Our model is as follows. A compressionally heated neutron star emits thermal neutrino pairs which, in turn, annihilate to produce a hot electron-positron pair plasma. We model the expansion of the plasma with a spherically symmetric relativistic hydrodynamics computer program. This simplification is justified at this stage of the calculations since the rotational velocity of the stars is about one third of the sound speed in the $`e^+e^{}`$ pair plasma. We then analyze and compare the contributions of photons from $`e^+e^{}`$ pair annihilation as well as from an external synchrotron shock as the plasma plows into the interstellar medium. We show that the characteristic features of GRBs, i.e. total energy, duration and gamma-ray spectrum, can be accounted for in the context of this model. ## 2 Compression in Close Neutron Star Binaries It has been speculated for some time that inspiraling neutron stars could provide a power source for cosmological gamma-ray bursts. However, previous Newtonian and post Newtonian studies (Janka & Ruffert, 1996; Ruffert & Janka, 1998, 1999) of the final merger of two neutron stars have found that the neutrino emission time scales are so short that it would be difficult to drive a gamma-ray burst from this source. It is clear that a mechanism is required for extending the duration of energetic neutrino emission. A number of possibilities could be envisioned, for example, neutrino emission powered by accretion shocks, MHD or tidal interactions between the neutron stars, etc. The present study, however, has been primarily motivated by numerical studies of the strong field relativistic hydrodynamics of close neutron-star binaries in three spatial dimensions. These studies (Wilson & Mathews, 1995; Wilson et al., 1996; Mathews & Wilson, 1997; Mathews et al., 1998; Mathews & Wilson, 2000) suggest that neutron stars in a close binary can experience relativistic compression and heating over a period of seconds. During the compression phase released gravitational binding energy can be converted into internal energy. Subsequently, up to $`10^{53}`$ ergs in thermally produced neutrinos can be emitted before the stars collapse (Mathews & Wilson, 1997, 2000). Here we briefly summarize the physical basis of this model and numerically explore its consequences for the development of an $`e^+e^{}`$ plasma and associated GRB. In (Mathews & Wilson, 1997, 2000) properties of equal-mass neutron-star binaries were computed as a function of mass and EOS (Equation of State). From these studies it was deduced that compression, heating and collapse could occur at times from a few seconds to tens of seconds before binary merger. Our calculation of the rates of released binding energy and neutron star cooling suggests that interior temperatures as hot as 70 MeV are achieved. This leads to a high neutrino luminosity which peaks at $`L_\nu 10^{53}`$ ergs sec<sup>-1</sup>. This much neutrino luminosity would partially convert to an $`e^+e^{}`$ pair plasma above the stars as is also observed above the nascent neutron star in supernova simulations (Wilson & Mayle, 1993). We should point out, however, that many papers have been published claiming the compression is nonexistent. In Mathews et al. (1998) we presented a rebuttal to the critics. Subsequently, however, Flanagan (1999) pointed out a spurious term in our formula for the momentum constraint. We (Mathews & Wilson, 2000) have corrected the momentum constraint equation and redone a sequence of calculations for a binary neutron star system with various angular momenta. A compression effect still exists which is able to release $`10^{52}`$ \- $`10^{53}`$ ergs of gravitational binding energy. The compression does not occur for corotating stars with a polytropic equation of state. For irrotational binary stars our compression effect is consistent with the results of at least two other groups (Bonazzola et al., 1999a, b; Marronetti et al., 1999; Uryu & Eriguchi, 1999) using different numerical methods to compute the relativistic hydrostatic equilibrium (Bonazzola et al., 1997). However, these calculations were done with a polytropic equation of state and found only very small compression; much less than 1%. For the polytropic equation of state Mathews & Wilson (2000) also found a compression less than 1%. In simulations with the realistic, somewhat soft, EOS described below we found clear evidence (Mathews & Wilson, 2000) that significant compression, heating and collapse still occurs for sufficiently close orbits. The reason for this EOS dependence is straightforward. Table 1 shows the realistic EOS used in the present work and the (Mathews & Wilson, 2000) studies. The key difference between the polytropic and realistic EOSs is that the adiabatic index $`\mathrm{\Gamma }`$ is not constant but decreases at low density for a realistic EOS. This causes the outer regions of the star to be more compact, and therefore, less affected by tidal stabilization than for polytropes. At the same time, the maximum central density tends to be larger for a neutron star of a fixed baryon mass. Therefore the relativistic effects are more dramatic when a realistic EOS is employed. The hydrodynamic calculations that demonstrate compression have been made with the stars constrained to remain at zero temperature (i.e. efficient radiators). As the compression rate increases, however, it is expected that the rate of released binding energy will exceed the ability of the star to radiate and internal heating will result. Large scale off-center vortices are observed to form (Mathews et al., 1998) within the stars with a characteristic circulation time scale of $`0.005`$ sec. The maximum velocities are nearly sonic. Among other things, this circulatory motion should help dissipate the compressional motion into thermal energy by shocks thereby heating the interior of the stars. We have run several sets of calculations with realistic neutron-star equations of state. We first considered stars like our earlier bench mark cases (Mathews et al., 1998) with a baryon mass of 1.548 $`M_{}`$ corresponding to a typical (cf. Appendix A) gravitational mass of M<sub>G</sub> = 1.39 M and a central density of $`\rho _c=1.34\times 10^{15}`$ g cm<sup>-3</sup> in isolation. These stars are based upon the “realistic” EOS of table 1 for which the maximum critical mass is $`M_c=1.575`$ M. As summarized in Appendix A, this maximum mass is typical of the somewhat soft EOS’s in which relativistic particles and/or condensates have been included. Parameters for this EOS were motivated by the necessity of such a soft EOS to obtain the correct neutrino signal in simulations of SN 1987A (Wilson & Mayle 1993). As noted in Table 5 of the Appendix this maximum mass is consistent with the measured masses of all binary pulsar systems for which the orbits have been well determined. As noted above, the stars calculated using this realistic EOS show significant compression and released binding energy before inspiral but do not individually collapse. The released gravitational energy from this calculation is summarized in Table 2. Even without the collapse instability enough internal heating occurs to produce a significant gamma-ray burst. We also found (Mathews & Wilson, 2000) that the individual collapse of stars would occur if the stars are increased in mass from M$`{}_{G}{}^{}=1.39`$ to 1.44 M (M$`{}_{B}{}^{}=1.61`$ M) for this EOS. Collapse of this star system is observed to occur for very close separation ($`d=2.4R`$) near (but before) the final stable orbit. Thus, collapse is a reasonable possibility for typical masses and a moderately soft EOS. For example, collapse would always occur prior to inspiral for stars in the typically observed mass range modeled with the EOS of Bethe & Brown (1995). For a critical mass of 1.54, even stars of initial mass of 1.35 collapse before reaching the innermost stable orbit. Based upon the above results, we model the thermal energy deposition due to neutron star compression as follows: we expect that the fluid motion within the stars will quickly convert released gravitational binding energy into thermal energy in the interior of the stars. Thus, we estimate that the rate of thermal energy deposition is comparable to the rate of released binding energy due to compression. The amount of released binding energy scales with the orbital four velocity (Mathews & Wilson, 1997, 2000). An estimate of the rate of increase of the orbital four velocity can be obtained (Mathews & Wilson, 1997) from the gravitational radiation timescale. Then, from the relation between released binding energy and increasing four velocity (Mathews & Wilson, 1997, 2000), the energy deposition rate into the stars can be deduced in approximate analytic form (Mathews & Wilson, 1997). We write, $$\dot{E}_{th}=\frac{(32/5)(Mf)^{5/3}fE_{th}^0}{[1(64/5)(Mf)^{5/3}ft]^{3/2}},$$ (1) where $`f`$ is the orbital angular frequency and $`E_{th}^0`$ is the total thermal energy deposited into the stars. In the hydrodynamic pair plasma discussions below we consider a range of deposited thermal energy of $`E_{th}^0=10^{51}`$, $`10^{52}`$, $`10^{53}`$ ergs, consistent with the hydrodynamics simulations. We use the convention that $`t<0`$ and $`t=0`$ is the end of energy deposition when the neutron stars either have collapsed into two black holes or have reached the last stable orbit and collapsed to a single black hole. At the time that a typical neutron star binary system is near the last stable obit, the orbital frequency is $``$ a few$`\times 10^3`$ sec<sup>-1</sup>. Hence, by Equation (1), the energy deposition rate would be $$\dot{E}_{th}10^2\times E_{th}^0\mathrm{erg}\mathrm{sec}^1.$$ (2) Thus, for $`E_{th}^0=2\times 10^{52}`$ ergs, $`\dot{E}_{th}2\times 10^{54}`$ ergs sec<sup>-1</sup>. The magnitude of the neutrino luminosity is very critical since the subsequent fireball is formed by neutrino-antineutrino annihilation. In order to model the thermal energy emitted by neutrinos before either stellar or orbital collapse we have constructed a computer model which treats the diffusion of energy in a static neutron star. This energy transport occurs by a combination of neutrino diffusion plus energy diffusion via a convective velocity-dependent diffusion coefficient. For the neutrino energy diffusion we write: $$\frac{dE_\nu }{dt}=\stackrel{}{}D_\nu \stackrel{}{}E_\nu $$ (3) where the coefficient for neutrino diffusion is just the form, $$D_\nu ^{rad}=\frac{c}{3\rho \kappa _\nu }+\frac{RV_c}{3},$$ (4) where a simple estimate for the neutrino opacity $`\kappa _\nu `$ is used, $`\kappa _\nu 9\times 10^{43}T_{MeV}^2`$ cm<sup>2</sup> g<sup>-1</sup> based upon the cross section for neutrino nuclear absorption and scattering. Characteristic convective velocities $`V_c`$ were deduced from simulations using our three-dimensional binary neutron star code (Mathews & Wilson, 2000). We calculated an angle averaged radial component of the fluid velocity in the frame of the star, $$V_c=\frac{1}{4\pi }\|V_r\|d(\mathrm{cos}\theta )𝑑\varphi .$$ (5) For our studies, these velocities were fit with an ansatz of the form $$V_c\frac{32}{105}V_{c,ave}\frac{r}{R}\sqrt{1r/R},$$ (6) where $`r`$ is the radial position inside a star of radius $`R`$ and $`V_{c,ave}`$ is the volume averaged $`V_c`$. This gives a good fit to the numerical results and has the correct form in that the velocity goes to zero at the surface and also at $`r=0`$. Energy was deposited in accordance with equation 1 and the calculations terminated at time=0. In Figure 1 the fraction of energy released, the peak luminosity, and $`\overline{L}`$ the average luminosity weighted by $`L^{5/4}`$ (see next section) are shown. The Energy input was $`2\times 10^{52}`$ ergs and the orbital frequency was 4000 rad sec<sup>-1</sup> (see Mathews & Wilson, 2000). The average convective velocity was found to be $`0.003c`$ by analyzing the hydrodynamical calculations Mathews & Wilson (2000) of neutron star binaries. From Figure 1 we see that a high emission efficiency and luminosity are obtained from this convective velocity. These produced lower thermal energies but about the same fraction of the energy emitted, $`68\%`$, but the $`\overline{L}`$ is reduced. For $`E_{th}^0=0.5`$ ($`1.0`$) $`\times 10^{52}`$ ergs $`\overline{L}=0.75`$ ($`1.5`$) $`\times 10^{53}`$ ergs sec<sup>-1</sup>. From these calculations we estimate that the conversion of compressional energy to fireball energy is probably $`{}_{}{}^{>}20\%`$. ## 3 Neutrino Annihilation and Pair Creation In the previous section we have outlined a mechanism by which neutrino luminosities of $`10^{53}`$ erg sec<sup>-1</sup> may arise from binary neutron stars approaching their final orbits. Here we argue that the efficiency for converting these neutrinos into pair plasma is probably quite high. Neutrinos emerging from the stars will deposit energy outside the stars predominantly by $`\nu \overline{\nu }`$ annihilation to form electron pairs. A secondary mechanism for energy deposition is the scattering of neutrinos from the $`e^+e^{}`$ pairs. Strong gravitational fields near the stars will bend the neutrino trajectories. This greatly enhances the annihilation and scattering rates (Salmonson & Wilson, 1999). Figure 2 taken from Salmonson & Wilson (1999) shows the relativistic enhancement factor, $`(R/M)`$, of the rate of annihilation by gravitational bending versus the radius to mass ratio (in units $`G=c=1`$). For our employed neutron-star equations of state the radius to mass ratio is typically between $`R/M3`$ and 4 just before stellar collapse. Thus, the enhancement factor ranges from $``$ 8 to 28. Defining the efficiency of energy deposition as the ratio of energy deposition to neutrino luminosity, then from Equation 24 of Salmonson & Wilson (1999) we obtain, $$\frac{\dot{Q}}{L_\nu }0.03(R/M)L_{53}^{5/4}.$$ (7) Thus, the efficiency of annihilation ranges from $``$0.1 to $`0.84\times L_{53}^{5/4}`$. For the upper range of luminosity the efficiency is quite large. To better analyze the annihilation process we have adapted the Mayle-Wilson (Wilson & Mayle, 1993) supernova model to this problem. We emphasize that the Mayle-Wilson model is fully general relativistic. To investigate this problem, a hot neutron star of the appropriate $`R/M`$ was constructed and the internal temperature adjusted to achieve the correct neutrino luminosities. The Courant condition requires that the time steps be quite small ($`10^9`$ second), and zonal masses as low as $`10^{13}M_{}`$ are required just outside of the neutrinosphere. Hence, the calculations could only be evolved for a short time. The entropy per baryon $`s/k`$ of the $`e^+e^{}`$ pair plasma is the critical quantity for gamma-ray production. It can be written, $$s/k=\frac{4m_bc^2(ae^3)^{1/4}}{3k\rho },$$ (8) where $`\rho `$ is the baryon density and $`e`$ is the total energy density. The entropy per baryon was found to be in the range of $`10^510^6`$ for the high luminosities. For a luminosity of $`10^{53}`$ ergs sec<sup>-1</sup>, an efficiency of energy transfer from the neutrinos to the $`e^+e^{}`$ pair plasma due to annihilation and electron scattering was found to be about 50 %. This efficiency of neutrino annihilation determines the total energy of the pair plasma and the entropy. This provides the initial conditions for the subsequent fireball expansion. ## 4 Pair Plasma Expansion Having determined the initial conditions of the hot $`e^+e^{}`$ pair plasma near the surface of a neutron star, we wish to follow its evolution and characterize the observable gamma-ray emission. To study this we have developed a spherically symmetric, general relativistic hydrodynamic computer code to track the flow of baryons, $`e^+e^{}`$ pairs, and photons. For the present discussion we consider the plasma deposited at the surface of a $`1.45M_{}`$ neutron star with a radius of 10 km. The fluid is modeled by evolving the following spherically symmetric general relativistic hydrodynamic equations: $$\frac{D}{t}=\frac{\alpha }{r^2}\frac{}{r}(\frac{r^2}{\alpha }DV^r)+\dot{D}_{in}$$ (9) $$\frac{E}{t}=\frac{\alpha }{r^2}\frac{}{r}(\frac{r^2}{\alpha }EV^r)P\left[\frac{W}{t}+\frac{\alpha }{r^2}\frac{}{r}(\frac{r^2}{\alpha }WV^r)\right]+\dot{E}_{in}$$ (10) $$\frac{S_r}{t}=\frac{\alpha }{r^2}\frac{}{r}(\frac{r^2}{\alpha }S_rV^r)\alpha \frac{P}{r}\alpha \frac{M}{r^2}\left(\frac{D+\mathrm{\Gamma }E}{W}\right)\left[\left(\frac{W}{\alpha }\right)^2+\frac{(U^r)^2}{\alpha ^4}\right]$$ (11) where $`D=\rho W`$ and $`E=ϵ\rho W`$ are the Lorentz contracted coordinate densities of baryonic and thermal mass energy ($`e^+e^{}`$ and photons) respectively. The quantities $`\dot{D}_{in}`$ and $`\dot{E}_{in}`$ refer to the injected plasma from neutrino pair annihilation, and $`S_r`$ is the radial coordinate momentum density. $`U_r`$ is the radial component of the covariant 4-velocity. $`W\alpha U^t`$ is the generalized Lorentz factor, $`V^r`$ is the radial coordinate three velocity, and $`\mathrm{\Gamma }`$ is an equation of state index. These quantities are defined by $`\alpha \sqrt{1{\displaystyle \frac{2M}{r}}};U_r`$ $``$ $`{\displaystyle \frac{S_r}{D+\mathrm{\Gamma }E}};W\sqrt{1+U^rU_r}`$ $`V^r{\displaystyle \frac{U^r}{W}}`$ ; $`\mathrm{\Gamma }1+{\displaystyle \frac{PW}{E}}`$ To evolve the $`e^+e^{}`$ pair plasma, we define a pair equation. The observed pair annihilation rate must be corrected for relativistic effects; specifically, time dilation slows the apparent pair annihilation process for a fast moving fluid with respect to an observer. Thus, we construct a continuity equation analogous to Equation (9) and add a term to account for annihilation and pair-production reactions: $$\frac{N_{pairs}}{t}=\frac{\alpha }{r^2}\frac{}{r}(\frac{r^2}{\alpha }N_{pairs}V^r)+\overline{\sigma v}((N_{pairs}^0(T))^2N_{pairs}^2)/W^2.$$ (13) Here, $`N_{pairs}`$ is the coordinate pair number density, and $`\overline{\sigma v}`$ is the Maxwellian averaged mean pair annihilation rate per particle. Although $`\overline{\sigma v}`$ depends on $`T`$, it varies little in the temperature range of interest, and thus, can be taken as constant: $`\overline{\sigma v}=2.5\times 10^{25}`$ cm<sup>3</sup> sec<sup>-1</sup>. $`N_{pairs}^0(T)=n_{pairs}^0(T)W`$, where $`n_{pairs}^0(T)`$ is the local proper equilibrium $`e^+e^{}`$ pair density at temperature T given by the appropriate Fermi integral with a chemical potential of zero. Zero chemical potential is a good approximation when $`N_{pairs}^0(T)`$ of Equation (13) is important. The total proper energy equation, including photons and $`e^+e^{}`$ pairs (baryon thermal energy is negligible), is $$e_{tot}=aT^4+e_{pairs}$$ (14) where coordinate energy in Equation (10) is related to proper energy by $`E=e_{tot}W`$ and $`e_{pairs}`$ is the appropriate zero chemical potential Fermi integral normalized to give the proper $`e^+e^{}`$ pair density $`n_{pair}=N_{pairs}/W`$ as determined by Equation (13). The entropy per baryon (Equation 8) of the wind is crucial to the behavior of the burst. An entropy that is too high will create a burst which is much hotter than those observed, while an entropy that is too low will extinguish the burst with baryons. We find that entropies of the order $`10^7`$ to $`10^8`$ are ideal for producing an isotropic burst directly from the expanding pair-photon plasma. In the calculations shown below (Sections 5 & 6) we cover a range of possible entropies per baryon from $`10^6`$ to $`10^8`$. Other possible sources of high entropy-per-baryon plasmas include the formation of magnetized black holes (Ruffini et al., 1998, 1999) and the high-energy collisions ($`\gamma 2`$) of stars in collapsing globular clusters, which we are studying in a separate work. We will deal with two paradigms for $`\gamma `$ray production. First, we treat the high entropy case ($`s/k>10^7`$) where the emission is from the fireball. Secondly, we present a low entropy case in which gamma emission arises from the collision of the fireball with the local interstellar medium. In the first case, the hydrodynamic equations are evolved as the plasma expands. Once the system becomes transparent to Thomson scattering, ($`N_{pair}(r)\sigma _T𝑑r_{}^<1`$ where $`\sigma _T`$ is the Thomson cross-section) we assume the photons are free-streaming, the calculation is stopped and the photon gas is analyzed to determine the photon signal. ## 5 Analysis of the Spectrum and Light curve We find that the photons and $`e^+e^{}`$ pairs appear to decouple at virtually the same time throughout the entire photon-$`e^+e^{}`$ pair plasma (when the cloud has reached a radius $`10^{12}10^{13}`$ cm and the temperature is typically a few 10’s of eV). As such, the photons will be well approximated as thermal and so we neglect any radiation transport effects. Thus, we take decoupling to be instantaneous and to occur when the plasma becomes optically thin to Thomson scattering. Furthermore, we find that virtually none of the energy deposited in the $`e^+e^{}`$ pair plasma remains in the pairs ($`.001`$%). Thus, the conversion of $`e^+e^{}`$ pair energy to photons and baryons is very efficient. From this simulation we derive two observables, the time integrated energy spectrum $`N(ϵ)`$ and the total energy received as a function of observer time $`\epsilon (t)`$. ### 5.1 The Spectrum As mentioned above, we assume that the $`e^+e^{}`$ pairs and photons are equilibrated to the same $`T`$ when they decouple. Thus, the photons in the fluid frame (denoted with a prime: ) make up a Planck distribution of the form $$u_{}^{}{}_{ϵ^{}}{}^{}(T^{})\frac{ϵ_{}^{}{}_{}{}^{3}}{exp(ϵ^{}/T^{})1},$$ (15) but $`u_ϵ/ϵ^3`$ is a relativistic invariant (Rybicki & Lightman, 1975). This implies $`ϵ/T`$ is also a relativistic invariant. So a Planck distribution in an emitter’s rest-frame with temperature $`T^{}`$ will appear Planckian to a moving observer, but with boosted temperature $`T=T^{}/(\gamma (1v\mathrm{cos}\theta ))`$ where $`v\mathrm{cos}\theta `$ is the component of fluid velocity (c=1) directed toward the observer. Thus, $$u_ϵ(\theta ,v,T^{})\frac{ϵ^3}{exp(\gamma (1v\mathrm{cos}\theta )\frac{ϵ}{T^{}})1}$$ (16) gives the observed spectrum of a blackbody with rest-frame temperature $`T^{}`$ moving at velocity $`v`$ and angle $`\theta `$ with respect to the observer. In the present case we wish to calculate the spectrum from a spherical, relativistically expanding shell as seen by a distant observer. Since we know $`v`$, $`T^{}`$ and the radius $`R`$ of the shell, we integrate over volume (i.e., shell, angle) with respect to the observer. We thus obtain the observed photon energy spectrum $`N_ϵ=(u_ϵ/ϵ)d^3x`$, from a relativistically expanding spherical shell with radius $`R`$, thickness $`dR`$, velocity $`v`$, Lorentz factor $`\gamma `$ and fluid-frame temperature $`T^{}`$, to be (in photons/eV/steradian) $$N_ϵ(v,T^{},R)=(5.23\times 10^{11})4\pi R^2dR\frac{ϵT^{}}{v\gamma }\mathrm{log}\left[\frac{1exp[\gamma ϵ(1+v)/T^{}]}{1exp[\gamma ϵ(1v)/T^{}]}\right](\mathrm{eV}^1sr^1),$$ (17) where $`R`$ is in cm. Note, that this spectrum has a maximum at $`ϵ_{max}1.39\gamma T^{}eV`$ for $`\gamma 1`$. We may then sum this spectrum over all shells (the zones in our computer code) of the fireball to get the total spectrum. Figure 3 shows an example of such a spectrum up to 500 keV. Since we assume a priori that the photons are thermal, our spectrum has a high frequency exponential tail, but the resultant total spectrum is clearly not thermal in the high energies. ### 5.2 The Light Curve To construct the observed light curve $`\epsilon (t)`$ we again decompose the spherical plasma into concentric shells and consider two effects: First, is the relative arrival time of the first light from each shell: light from outer shells will be observed before light from inner shells; Second, is the shape of the light curve from a single shell. Emission from moving pair plasma is beamed along the direction of travel within an angle $`\theta 1/\gamma `$. The surface of simultaneity of a relativistically expanding spherical shell, as seen by an observer, is an ellipsoid (Fenimore et al., 1996). The observer time of intersection of an expanding ellipse with a fixed shell of radius R as a function of $`\theta `$ (i.e. the time at which emission from this intersection circle is received) is: $$t=\frac{R}{v}(1v\mathrm{cos}\theta )\frac{R}{2\gamma ^2c},$$ (18) for $`\theta 1,\gamma 1`$. Integrating our boosted Planck distribution of photons (Equation 16) over frequency, we find that a relativistically expanding shell of radius R will have a time profile (energy/time/steradian) $$\epsilon (\tau ,v,T^{},R)=\frac{a}{2}\left(\frac{T^{}}{\gamma \tau }\right)^4cRdR1/\tau ^4,$$ (19) for $`\tau >1`$ and where $`\tau \frac{vt}{R}`$. Emission starts at $`\tau _i=(1v/c)`$ and ends at $`\tau _f=(1+v/c)`$. The final light curve is constructed by summing the signal from all shells. The total thickness of the expanding plasma is $`cJ/\dot{J}`$ because it expands at near the speed of light and $`J/\dot{J}`$ is the timescale of compression and coalescence which sets the emission timescale. Typically $`R10^{12}`$ cm and $`J/\dot{J}`$ a few seconds, so $`cJ/\dot{J}R`$ and the emitting plasma is a thin shell. The duration of the burst is determined by the duration of emission because the observed timescale of emission from the plasma shell is very short, $`R/2\gamma ^2c0.01`$ seconds (Equation 18) for $`\gamma 100`$, compared to $`J/\dot{J}`$. ## 6 Results of Pair Plasma Emission We have run a variety of models over a range of entropies per baryon and total energies. The results are summarized in Figures 4, 5 & 6. We see that more powerful bursts are derived from higher entropies per baryon and higher total energies. In particular, entropies per baryon of a few $`\times 10^7`$ allow a burst with a spectral peak $`100`$ keV and efficiencies $`E_\gamma /E_{tot}10`$%. This is consistent with, although at the upper end of the range of, the entropies calculated for the $`e^+e^{}`$ plasma deposited above the neutron stars. Much further work needs to be done to better characterize the nature of the stellar compression and energy transport within the stars. Also, more elaborate simulations must be done to resolve the plasma flow in three dimensions and to consider the effects of magnetic fields. In Table 3 we see the final Lorentz factor for a range of expanding fireballs. This data will be used when we look at the collision of the fireball into an external medium. ## 7 External Shock Emission In previous sections the emission from an expanding fireball was studied. We found that the resulting emission spectrum and total energy strongly depends upon the energy of the plasma deposited near the surface of the neutron stars; entropies of $`10^6`$ resulted in weak emission with most of the original energy manifesting itself as kinetic energy of the baryons. Thus, for the low entropy per baryon fireballs ($`s10^510^6`$) produced by NSBs it is necessary to examine the emission due to the interaction of the relativistically expanding baryon wind with the interstellar medium (ISM). After becoming optically thin and decoupling with the photons, the matter component of the fireball continues to expand and interact with the ISM via collisionless shocks. As the ISM is swept up, the matter decelerates. We model this process as an inelastic collision between the expanding fireball and the ISM as in, for example, Piran (1999). We assume that the absorbed internal energy is immediately radiated away. From this we construct a simple picture of the emission due to the matter component of the fireball “snowplowing” into the ISM of baryon number density $`n`$. For a shell of a given rest mass $`M`$ expanding at Lorentz factor $`\gamma `$, the conservation of momentum leads to the following constraint equation: $$\frac{d\gamma }{\gamma ^21}=\frac{dM}{M},$$ (20) which has the solution $$\frac{M}{M_0}=\sqrt{\frac{(\gamma _01)(\gamma +1)}{(\gamma _0+1)(\gamma 1)}}.$$ (21) Now we can put this in terms of radius by noting $$M=M_0+\frac{4\pi }{3}nm_pc^2R^3.$$ (22) Thus, $$R(\gamma )=R_0\left(\frac{M}{M_0}1\right)^{1/3}R_0\left(\frac{1}{\gamma }\frac{1}{\gamma _0}\right)^{1/3}\text{for }\gamma \text{}\gamma _01,$$ (23) where $$R_0\sqrt[3]{\frac{3M_0}{4\pi nm_pc^2}},$$ (24) is the radius at which $`M=2M_0`$. This is the characteristic radius at which the shock decelerates. We assume that the local thermal energy radiated away after a thin shell of ISM mass $`dM`$ is swept up by the shock is $$dE^{}=(\gamma 1)dM.$$ (25) The observer time elapsed for the mass to expand a distance $`dR`$ is $$dt_{obs}=\frac{dR}{2\gamma ^2c}.$$ (26) Equations (23,26) can be solved in the relativistic limit to give $$t(\gamma ,\gamma _0)\frac{R_0}{28c}\left(\frac{9}{\gamma _0^2}+\frac{3}{\gamma \gamma _0}+\frac{2}{\gamma ^2}\right)\left(\frac{1}{\gamma }\frac{1}{\gamma _0}\right)^{1/3}\text{for }\gamma \text{}\gamma _01.$$ (27) The implied observer luminosity, from Equations (26, 25), is $$L=\frac{dE}{dt_{obs}}=\frac{\gamma dE^{}}{dt_{obs}}8\pi R^2\gamma ^4nm_pc^3$$ (28) for $`\gamma 1`$. Using Equation (27), a relativistic ($`\gamma 1`$) solution for observed luminosity, in ergs sec<sup>-1</sup>, over several epochs is $$L(t)\{\begin{array}{cc}2.68\times 10^{50}n\gamma _{300}^8t^2(16.27\times 10^3\gamma _{300}^8nE_{52}^1t^3)^{10/3}\hfill & t<t_1\hfill \\ & \\ 7.88\times 10^{51}n^{1/3}E_{52}^{2/3}\gamma _{300}^{8/3}(0.32\frac{t}{t_{max}}0.15)^{2/3}(1.150.32\frac{t}{t_{max}})^{10/3}\hfill & t_1<t<t_2\hfill \\ & \\ 5.3\times 10^{51}n\gamma _{300}^4\left[\frac{E_{52}}{n\gamma _{300}}\right]^{4/7}t^{2/7}\left(1\sqrt{f(t)}\right)^4\left(f(t)+\sqrt{f(t)}\right)^{2/3}\hfill & t>t_2\hfill \end{array}$$ (29) where constant parameters are, $`E_{52}E/10^{52}`$ ergs, $`\gamma _{300}\gamma _0/300`$ and $`n`$ is in baryons cm<sup>-3</sup>. For $`t>t_2`$: $$f(t)11.05\left(\frac{t_{max}}{t}\right)^{3/7}$$ (30) and $$t_{max}3.5\sqrt[3]{\frac{E_{52}}{n\gamma _{300}^8}}\text{seconds}$$ (31) is the observer time at maximum luminosity $`L_{max}`$: $$L(t_{max})=L_{max}=1.3\times 10^{51}n^{1/3}\gamma _{300}^{8/3}E_{52}^{2/3}\text{ergs/sec}.$$ (32) The times at which the solutions for each epoch are spliced together are roughly $`t_1`$ $`0.6t_{max}`$ (33) $`t_2`$ $`1.5t_{max}.`$ (34) Figure 7 shows the light curve for a $`10^{52}`$ erg fireball expanding at $`\gamma =300`$ for a range of ISM densities. This corresponds to an initial energy deposition above the neutron stars with an entropy per baryon of $`s=10^5`$ as seen in Table 3. The expansion can be divided into a free-expansion phase and a deceleration phase: $$L(t)\{\begin{array}{cc}t^2\hfill & \text{free expansion phase }(t<t_{max})\hfill \\ t^{10/7}\hfill & \text{deceleration phase }(t>t_{max})\hfill \end{array}$$ (35) Figure 8 shows a linear plot of the light curve for ISM density $`n=1.0`$ baryons cm<sup>-3</sup>. The “fast-rise, exponential-decay” or “FRED”-like shape is evident and is in good qualitative agreement with “smooth” GRBs. ### 7.1 Synchrotron Shock Spectrum Now we wish to model the spectrum of light emitted as the fireball expands into the ISM. To do this we assume an external synchrotron shock model (Shemi & Piran, 1990; Rees & Mészáros, 1992; Mészáros & Rees, 1993; Rees & Mészáros, 1994). This analysis is analogous to afterglow models in the radiative limit. Thus, the spectrum will have the form (Sari et al., 1998) $$L_\nu =\{\begin{array}{cc}(\nu /\nu _c)^{1/3}L_{\nu ,max}\hfill & \nu <\nu _c\hfill \\ (\nu /\nu _c)^{1/2}L_{\nu ,max}\hfill & \nu _m>\nu >\nu _c\hfill \\ (\nu _m/\nu _c)^{1/2}(\nu /\nu _m)^{p/2}L_{\nu ,max}\hfill & \nu >\nu _m.\hfill \end{array}$$ (36) Photon frequency $`\nu _m`$ is the frequency corresponding to the minimum energy of the electron distribution above which the electrons are assumed to have a power law functional form $`n(\gamma )\gamma ^p`$. In the numerical examples that follow, we take the spectral index to be $`p=2.5`$. This is consistent with that calculated for ultrarelativistic shocks (Bednarz & Ostrowski, 1998). The “cooling frequency” $`\nu _c`$ corresponds to the energy below which the electrons cannot cool on a hydrodynamic timescale. The peak of the luminosity spectrum is $$L_{\nu ,max}\left(\frac{p2}{2p2}\right)\frac{L}{\sqrt{\nu _m\nu _c}},$$ (37) assuming $`\nu _m\nu _c`$, which is valid throughout the burst duration. There are two free parameters in this model. $`ϵ_e`$ is the fraction of the kinetic energy of the baryons that is deposited into the electrons by the shock. $`ϵ_B`$ is the ratio of the magnetic field energy density to the kinetic energy density of the baryons. In these simulations we take each of these values to be $`1/4`$. The evolution of the characteristic frequency $`\nu _m`$ is described by $$\nu _m1.4\times 10^4ϵ_e^2ϵ_B^{1/2}\sqrt{n}\gamma _{300}^4\text{keV}\{\begin{array}{cc}(1\left(\frac{2ct}{R_0}\right)^3\gamma _0^7)^4\hfill & t<t_1\hfill \\ & \\ (1.150.32t/t_{max})^4\hfill & t_1<t<t_2\hfill \\ & \\ \left(1\sqrt{1\frac{2}{\gamma _0}(\frac{14ct}{R_0})^{3/7}}\right)^4\hfill & t>t_2.\hfill \end{array}$$ (38) The behavior of the cooling frequency $`\nu _c`$ is more difficult to characterize since it depends on the hydrodynamical timescale of the fluid. Fortunately, however, $`\nu _c`$ is much smaller than $`\nu _m`$. Therefore, its exact behavior is not important for this analysis. Thus, we assume $`\nu _c`$ to be constant at early times and follow its asymptotic power-law at later times: $$\nu _c2.7\times 10^3ϵ_B^{3/2}E_{52}^{4/7}\gamma _{300}^{4/7}n^{13/14}\text{keV}\{\begin{array}{cc}t_{max}^{2/7}\hfill & tt_{max}\hfill \\ t^{2/7}\hfill & t>t_{max}.\hfill \end{array}$$ (39) The spectrum of the burst at peak luminosity $`L_{max}`$ is shown in Figure 9. For $`n=1.0`$ baryons cm<sup>-3</sup>, most of the energy is emitted at photon energies $`100`$ keV. Using Equations (29,38,39) for $`L`$, $`\nu _m`$ and $`\nu _c`$ respectively, we can determine the spectrum (Equation 36). The fluence spectrum of the burst is obtained by integrating the evolving luminosity spectrum (Equation 36) in time. This is shown in Figure 10. This figure again shows that most of the burst energy is in photons of several hundred keV energy. Now we can ask what the efficiency is of gamma-ray production by the shock compared to other wavelengths. At any time, the fraction of luminosity above a given minimum frequency $`\nu _{min}`$ is $$\epsilon _{ff}=\{\begin{array}{cc}\frac{\left[\left(\frac{2p2}{p2}\right)\nu _m^{1/2}2\nu _{min}^{1/2}\right]}{\left[\left(\frac{2p2}{p2}\right)\nu _m^{1/2}\frac{5}{4}\nu _c^{1/2}\right]}\hfill & \nu _c<\nu _{min}<\nu _m\hfill \\ & \\ \frac{\left(\frac{2}{p2}\right)\left(\frac{\nu _m^{p1}}{\nu _{min}^{p2}}\right)^{1/2}}{\left[\left(\frac{2p2}{p2}\right)\nu _m^{1/2}\frac{5}{4}\nu _c^{1/2}\right]}\hfill & \nu _{min}>\nu _c.\hfill \end{array}$$ (40) Thus, we can calculate the duration, $`t_{90}`$, of the luminosity at energies above this minimum energy. This is done in Table 4 for the various fireballs shown in Table 3. There is a competition between factors limiting the duration; lower energy fireballs simply have fewer high energy photons, and thus, shorter duration; while higher energy fireballs expand and evolve faster and thus have shorter duration. Fireballs with energy of order $`10^{52}`$ ergs and entropies per baryon of order $`10^5`$ yield a value for $`t_{90}`$ which is consistent with observation. The overall efficiency of the production of photons above a frequency $`\nu _{min}`$ is $$\begin{array}{cc}\hfill \epsilon _{fftot}& \frac{_0^{\mathrm{}}_{\nu _{min}}^{\mathrm{}}L_\nu 𝑑\nu 𝑑t}{_0^{\mathrm{}}_0^{\mathrm{}}L_\nu 𝑑\nu 𝑑t}\hfill \\ & 1\left(\frac{\nu _{min}}{\nu _{max}}\right)^{1/6}\text{for }\nu _{min}\nu _{max},\hfill \end{array}$$ (41) where $`\nu _{max}`$ is the value of $`\nu _m`$ (Equation 38) at $`t=t_{max}`$. For $`\nu _{min}=10`$ keV we have an overall efficiency of about $`\epsilon _{fftot}75`$ %. Thus, the radiative external shock GRB is quite efficient at producing gamma-rays if our assumptions are reasonable. ## 8 Conclusions In this paper we have argued that heated neutron stars (perhaps by compression of close neutron-star binaries) are viable candidates for the production of large, high entropy per baryon, $`e^+e^{}`$ pair plasma fireballs, and thus, for the creation of gamma-ray bursts. We find that fireballs of total energy $`E10^{51}`$ to $`3\times 10^{52}`$ ergs and an entropy per baryon of $`s/k10^510^6`$ are possible. Values for the entropy as high as $`10^7`$ may be realized during the peak $`\nu \overline{\nu }`$ luminosity (of $`10^{53}`$ ergs sec<sup>-1</sup>). Emergent gamma-rays yield a quasi-thermal spectrum peaked at $`100`$ keV with an efficiency of conversion from pair plasma to photons of $`30`$%. The lower entropy component of the fireball will initiate a shock which propagates into the ISM, generating an external shock GRB. The calculation utilizing the supernova computer program (Section 3) to describe the neutrino and matter transport, produces a baryonic wind that contains $`90`$% neutrons. The decay of the neutrons to protons occurs on the same time scale as that for which the protons are decelerated by the intersteller medium. This delayed conversion of neutrons to protons will broaden the gamma-ray signal by a factor of a few. In addition, the decay electrons will strongly increase the entropy of the expanding plasma at the late times. In future work we will quantify the role of neutrons and explore the possibility of fireball photons inverse-Compton up-scattering off of the accelerated electron distribution of the external shock. This corresponds to the emission scenario put forward by Liang et al. (1997). As of yet this model is spherically symmetric. Thus, it can only generate bursts with a smooth light-curve structure. However, we expect a large variety of GRB morphologies with varied time structure due to: 1) three dimensional resolution of the plasma flow; 2) plasma instabilities due to increased heating of the deposited plasma with time; and 3) variation in the ratio of star mass in the NSB, effecting the relative compression and heating rate of each star. In future work we will numerically model, in three dimensions, the flow of the $`e^+e^{}`$ pair plasma in the midst of the orbiting neutron stars. We have written a three-dimensional general relativistic hydrodynamic code to study this three-dimensional behavior. In particular, we wish to study the possible formation of jets along the orbital axis due to the collision of plasma blowing away from each star. Also, we have done simulations which suggest that the internal magnetic field of the neutron stars may be high. Thus, the inclusion of magnetohydrodynamic plasma effects including Alfvén instabilities and reconnections may ultimately be necessary. The authors wish to thank the late Jean-Alain Marck for his inciteful and encouraging remarks on an earlier version of this manuscript. This work was performed under the auspices of the U.S. Department of Energy by University of California Lawrence Livermore National Laboratory under contract W-7405-ENG-48. J.R.W. was partly supported by NSF grant PHY-9401636. Work at University of Notre Dame supported in part by DOE grant DE-FG02-95ER40934, NSF grant PHY-97-22086, and by NASA CGRO grant NAG5-3818. ## Appendix A Appendix: Neutron Star EOS A key requirement of a gamma-ray burst paradigm based upon collapsing neutron stars in binaries is that the equation of state be relatively soft so that significant compression and heating can occur before inspiral. Therefore, for completeness in this Appendix we review arguments for and against a “soft” neutron star EOS. The neutron star EOS must extend from normal iron nuclei on the surface to as much as 15 times nuclear matter density in the interior. At the same time, one must consider that neutron stars in weak-interaction equilibrium are highly isospin asymmetric. They may also carry net strangeness. Therefore, only pieces of the neutron-star equation of state, e.g. the nuclear compressibility, are accessible in laboratory experiments. The value for the nuclear compressibility $`K_s`$ can be derived from the nuclear monopole resonance (Blaizot, 1980). The present value ($`K_s=`$ 230 MeV) is consistent with a modestly soft nuclear equation of state. Nuclear heavy-ion collision data can also be used to shed some insight, particularly for the heated neutron-star equation of state. For example, McAbee & Wilson (1994) studied heavy ion collisions of <sup>139</sup>La on <sup>139</sup>La as a means to constrain the supernova EOS. The electron fraction for $`{}_{57}{}^{}{}_{}{}^{139}`$La ($`Y_e=0.41`$) overlaps that of supernovae which range from $`Y_e=0.05`$ to 0.50. They showed that the pion contribution to the EOS could be constrained by the observed pion multiplicities from central collisions. The formation and evolution of pions was computed in the context of Landau-Migdal theory to model the effective energy and momenta of the pions. A key aspect of hydrodynamic simulations of the heavy-ion data was the determination of the Landau parameter $`g^{}`$. Their determination of the pion contribution to the equation of state implies a relatively soft equation of state after pion condensation such that a maximum neutron-star mass of $`M1.64`$ M is inferred. There have been dozens of nuclear equations of state introduced over the years. Summaries of some of them can be found in Schaab & Weigel (1999) and Arnett & Bowers (1977). As far as the maximum mass of a neutron star is concerned, most theoretical equations of state fall into two groups, those which only describe the mean nuclear field even at high density and those which allow for various condensates, e.g. pions, kaons, hyperons, and even quark-gluon plasma. Table 5 summarizes the basic neutron star properties based upon most available nuclear equations of state (Lattimer, 1998). Equations of state which are based upon the mean nuclear field tend to be “stiff” at high density. Therefore, they reach lower interior densities for the same baryonic mass and tend to allow a higher maximum neutron-star mass $`m_{max}1.82.2`$ M. Such equations of state also tend to become acausal at the high densities associated with the maximum neutron-star mass. On the other hand, the relativistic equations of state are generally causal at high density. They also tend to be somewhat “soft”, therefore allowing a higher central density for a given baryon mass and generally implying a maximum neutron-star mass in the range $`m_{max}1.31.7`$ M. We note, however, that recent 3-body corrections to a relativistic EOS (Akmal et al., 1999) tend to stiffen an otherwise soft relativistic EOS. For the most part, constraints on the neutron star equation of state must ultimately come from observations of neutron stars themselves. Over the years attempts have been made with limited success to constrain the equation of state based upon the maximum observed rotation frequency (e.g. Friedman et al., 1986) or the thermal response to neutron star glitches (e.g. Page, 1998). In recent years, however, new observational constraints on the structure and properties of neutron stars are becoming available (Lattimer, 1998). Observations of quasi-periodic oscillations (QPOs) (Strohmayer et al., 1996; Van Der Klis et al., 1996, 1997), pulsar light curves (Yancopoulos et al., 1994; Thorsett & Chakrabarty, 1999), and glitches (Link et al., 1999), studies of soft-gamma repeaters (Kouveliotou et al., 1998; Gotthelf et al., 1999); and even the identification of an isolated non-pulsing neutron stars (Walter et al., 1996; Haberl et al., 1997) have all led to the hope that significant constraints on the mass-radius relation and maximum mass of neutron stars may be soon coming. ### A.1 Pulsars Two possible constraints come from measured pulsar systems. The most precisely measured property of any pulsar system is its spin frequency. The frequencies of the fastest pulsars (PSR B1937+21 at 641.9 Hz and B1957+20 at 622.1 Hz) already constrain the equation of state under the assumption that these pulsars are near their maximum spin frequency (Friedman et al., 1986). In particular, the equation of state cannot be too stiff, though maximum masses as large as 3 M are still allowed. A much more stringent constraint may come from the numerous determinations of neutron-star masses in pulsar binaries. There are now about 50 known pulsars in binary systems. Of these 50, approximately 15 of them have significantly constrained masses. These are summarized in Table 6. The measured masses are all consistent with low neutron-star masses in the range $`m1.35\pm 0.10`$ M (Thorsett & Chakrabarty, 1999). Even though these masses are low, this does not necessarily mean that the maximum neutron-star mass is in this range. If one adopts these masses as approaching the maximum neutron-star mass, then the softer equations of state are preferred. However, this narrow mass range may be the result of the mechanism of neutron-star formation in supernovae and not an indication of the maximum neutron-star mass. In a recent paper, Link et al. (1999) have proposed that glitches observed in the Vela pulsar and six other pulsars may place some constraint on the nuclear EOS. In particular, if the glitches originate from the liquid of the inner crust, and if the mass of the Vela pulsar is 1.35 consistent with Table 2, then the radius of the Vela pulsar must be $`R_{}^>8.9`$ km. This result is consistent with either a soft or stiff equations of state. A better theoretical determination of the pressure at the crust-core interface might lead to a more stringent constraint. ### A.2 QPO’s The identification of kilohertz QPO’s with the last stable orbit around a neutron star also could significantly constrain the neutron-star equation of state (e.g. Schaab & Weigel, 1999). For example, demanding that the 1.2 khz QPO from source KS 1731-260 be the last stable orbit requires a neutron-star mass of 1.8 M, On the other hand, other interpretations are possible for the origin of QPO’s. For example, they could be a harmonic of a lower frequency outer orbit, or they might result from effects closer to the neutron-star surface. Among proposals for the source of the QPO phenomenon are: boundary layer oscillations (Collins et al., 1998); radial oscillations and diffusive propagation in the transition region between the neutron star and the last Keplerian orbit (Titarchuk & Osherovich, 1999); Lense-Thirring precession for fluid particles near the last stable orbit (Miller et al., 1998; Merloni et al., 1999; Morsink & Stella, 1999); and nonequitorial resonant oscillations of magnetic fluid blobs (Vietri & Stella, 1998). ### A.3 Supernova constraints The lack of a radio pulsar in SN1987A, along with nucleosynthesis constraints on the observed change of helium abundance with metallicity has led to the suggestion (Brown & Bethe, 1994; Bethe & Brown, 1995) that the maximum neutron-star mass must be $`{}_{}{}^{<}1.56`$ M. In this picture, the development of a kaon condensate tends to greatly soften the EOS after $`12`$ sec. Thus, even though neutrinos were emitted, the core subsequently collapses to a black hole. One constraint comes from the neutrino signal itself observed to arise from supernova SN1987A. The fact that the neutrinos arrived over an interval of at least twelve seconds implies a significant cooling and neutrino diffusion time from the core. This favors a soft equation of state in which the core is more compact and at higher temperature in the supernova models. For example, the simulations of Wilson & Mayle (1993) require a maximum neutron-star mass of $`{}_{}{}^{<}1.6`$ M. ### A.4 Isolated Neutron Star A most promising constraint on the neutron-star EOS may come from the determination of the radius for the isolated nonpulsing neutron star RX J185635-3754, first detected by ROSAT (Walter et al., 1996). The inferred (redshifted) surface temperature from the X-ray emission is about 35 eV. Atmospheric models of this emission then imply (Lattimer, 1998; An et al., 1998; Wang et al., 1999) that for a distance between 31 and 41 pc, a radius between $`5.75<R/\mathrm{km}<11.4`$ and a mass of $`1.3<M<1.8`$, is most consistent with the observed emission. This is suggestive of a soft equation of state. However, this constraint requires that the distance be less than 41 kpc. On the other hand, Wang et al. (1999) find that the cooling properties of the soft X-ray source RX J0720.4-3125 are most consistent with a moderately stiff or stiff EOS provided that the age of this star is less than 10<sup>5</sup> yr. Proper motion studies with HST are currently underway to determine a reliable distance to RX J185635-3754. These studies will provide a key constraint on the nuclear equation of state.
warning/0002/math-ph0002027.html	ar5iv	text	# Dominos and the Gaussian free field. ## 1 Introduction A domino tiling of a polyomino $`P`$ in $`^2`$ is a tiling of $`P`$ with $`2\times 1`$ and $`1\times 2`$ rectangles. For a polyomino $`P`$ let $`\mu =\mu (P)`$ denote the uniform measure on the set of all domino tilings of $`P`$. Let $`U^2`$ be a Jordan domain with smooth boundary. We study uniform random domino tilings of polyominos $`P_ϵ`$ in $`ϵ^2`$ which approximate $`U`$ (and using dominos which are $`2ϵ\times ϵ`$ and $`ϵ\times 2ϵ`$ rectangles). A domino tiling of a polyomino $`P_ϵ`$ in $`ϵ^2`$ can be thought of as a random map from $`ϵ^2P_ϵ`$ to $``$ in the following way. Let $`V_ϵ=ϵ^2P_ϵ`$ be the set of lattice points in the polyomino $`P_ϵ`$. Let $`h:V_ϵ`$ be a function which has the property that around every lattice square of $`P_ϵ`$ the four values of $`h`$ are $`4`$ consecutive integers $`h_0,h_0+1,h_0+2,h_0+3`$, with the values on any two adjacent boundary vertices of $`P_ϵ`$ differing by $`1`$. The set of such functions $`h`$ (up to additive constants and a global sign change) is in bijection with the set of domino tilings of $`P_ϵ`$: dominos cross exactly those edges whose $`h`$-difference is $`3`$. The function $`h`$ associated to a tiling is called its height function . See Figure 1. Note that the height function takes values in $``$, not in $`ϵ`$. Our aim is to prove that in the limit as $`ϵ0`$ the height function on a random tiling of $`P_ϵ`$ tends to a random (generalized) function which has a succinct description in terms of the eigenbasis of the Laplacian operator on $`U`$. ###### Theorem 1.1 Let $`U`$ be a Jordan domain with smooth boundary in $`^2`$. For each $`ϵ>0`$ sufficiently small let $`P_ϵ`$ be a Temperleyan polyomino approximating $`U`$ as described below. Let $`h_ϵ`$ be the height of a random domino tiling of $`P_ϵ`$ and $`\overline{h}_ϵ`$ be its mean value. Then as $`ϵ`$ tends to $`0`$, $`h_ϵ\overline{h}_ϵ`$ tends weakly in distribution to $`4/\sqrt{\pi }`$ times the “massless free field” $`F`$ on $`U`$, in the sense that for any smooth function $`\varphi `$ on $`U`$, the random variable $`_{xV_ϵ}\varphi (x)(h_ϵ(x)\overline{h_ϵ(x)})`$ tends in distribution to $`\frac{4}{\sqrt{\pi }}_U\varphi F`$. For the definition of Temperleyan polyominos see below. The massless free field $`F`$ on $`U`$ is a random variable taking values in the space of distributions <sup>1</sup><sup>1</sup>1Henceforth we will refer to these objects as “generalized functions” to avoid confusion. which are continuous linear functionals on the space of $`C^1`$ functions on $`U`$ (with a $`C^1`$-norm). For background on the massless free field see . It can be defined as follows: let $`\{f_i\}_{i1}`$ be an $`L^2`$-orthonormal eigenbasis for the Laplacian $`\mathrm{\Delta }=\frac{^2}{x^2}+\frac{^2}{y^2}`$ on $`U`$ with Dirichlet boundary conditions (that is, $`f_i0`$ on $`U`$). Let $`\lambda _i`$ be the eigenvalue of $`f_i`$. Then $$F=\underset{i1}{}\frac{c_if_i}{(\lambda _i)^{1/2}},$$ (1) where the $`c_i`$ are i.i.d. Gaussian random variables of mean $`0`$ and variance $`1`$. Here this expression is interpreted as the generalized function $`F`$ satisfying, for any $`C^1`$ function $`\varphi `$, $$_U\varphi F=\underset{i1}{}\frac{c_i}{(\lambda _i)^{1/2}}_U\varphi f_i,$$ series which converges almost surely. The expression (1) does not define a function since the series diverges almost everywhere. Remarks. 1. The above theorem describes the limiting value of $`h_ϵ\overline{h}_ϵ`$. The limiting average value $`\overline{h}=lim\overline{h}_ϵ`$ was computed in : it is a harmonic function whose boundary values are given by $`\frac{2}{\pi }`$ times the angle of turning of the boundary tangent counterclockwise from a fixed basepoint. (Regarding the choice of basepoint, see the definition of “Temperleyan” polyomino below.) 2. Theorem 1.1 has a well-known one-dimensional analog: Let $`X`$ be the sets of random maps $`h`$ from $`0,\frac{1}{n},\frac{2}{n}\mathrm{},1`$ to $``$ satisfying $`h(0)=h(1)=0`$ and $`\|h(\frac{i+1}{n})h(\frac{i}{n})\|=1`$. A random element of $`X`$, when divided by $`\sqrt{n}`$, converges to a random function known as the “Brownian bridge” . In the eigenbasis of the one-dimensional Laplacian $`\frac{^2}{x^2}`$ the coefficients of the Brownian bridge are again independent Gaussians. One difference between the one-dimensional case and Theorem 1.1, however, is that the height function $`h`$ of Theorem 1.1 is unnormalized. It is therefore all the more surprising that the integer-valued function $`h`$ of Theorem 1.1 converges to a continuous-valued object. 3. An important open problem is to compute the distribution of the height function on a non-simply connected domain, even an annulus. In particular for an annulus the distribution of the height difference between the two boundary components (in the limit $`ϵ0`$) is unknown although it was shown in to depend only on the conformal modulus of the annulus. 4. Temperley gave a bijection between the uniform spanning tree process on subgraphs of $`^2`$ and domino tilings. The function $`h`$ of Theorem 1.1 corresponds under this bijection to the “winding number” of the branches of a spanning tree , as first conjectured by I. Benjamini. As it is an open question to show that a scaling limit exists for the uniform spanning tree process , one might hope that the reconstruction of the tree from its winding numbers, which is possible for $`ϵ>0`$, also works in the limit $`ϵ=0`$. So far this remains an open problem. 5. The result of Theorem 1.1 depends strongly on the choice of boundary conditions for the approximating polyominos $`P_ϵ`$. For even slight generalizations of these boundary conditions our methods will not work: see for a discussion of this issue. 6. When the region $`U`$ is a rectangle $`U=[0,a]\times [0,b]`$, the orthonormal eigenvectors of $`\mathrm{\Delta }`$ with Dirichlet boundary conditions are $`\frac{4}{ab}\mathrm{sin}\frac{\pi jx}{a}\mathrm{sin}\frac{\pi ky}{b}`$, where $`j,k`$ are positive integers. So in this case the massless free field has independent Fourier coefficients. 7. Most of the work to prove Theorem 1.1 was done in , where we proved Proposition 2.2, below. If we consider the massless free field $`F`$ to be a continuous linear functional on the space of smooth $`2`$-forms on $`U`$ (rather than on the space of smooth functions on $`U`$) then $`F`$ is conformally invariant, in the following sense. ###### Proposition 1.2 Let $`\omega `$ be a smooth $`2`$-form on $`U`$ and let $`f:VU`$ be a conformal bijection. Let $`F_U,F_V`$ be the massless free fields on $`U`$ and $`V`$ respectively. Let $`X=_UF_U(z)\omega (z)`$ and $`Y=_VF_V(z)f^{}\omega (z),`$ where $`f^{}\omega `$ is the pullback of $`\omega `$ to $`V`$. Then the random variables $`X`$ and $`Y`$ are equal in distribution. For the proof see section 4. ## 2 Background and preliminaries ### 2.1 Temperleyan polyominos and approximation Define the $`(i,j)`$-lattice square in $`^2`$ to be the lattice square whose lower left corner is $`(i,j)`$. A lattice square is said to be even if the coordinates of its lower left corner are even. A polyomino is a union of lattice squares which is bounded by a simple closed lattice curve. A polyomino is even if all of its corner squares are even, where by corner squares we mean those lattice squares adjacent to the corners and containing the interior angle bisector at the corner. In particular note that an edge of an even polyomino $`P^{}`$ has odd length if its two extremities are both concave or both convex corners; if the extremities consist of one concave and one convex corner the edge length is even. Let $`P`$ be a polyomino obtained from an even polyomino $`P^{}`$ by removing one lattice square $`b`$ adjacent to its boundary, where $`b`$ is of the same parity as the corners of $`P^{}`$. Such a polyomino is called Temperleyan, and the removed square $`b`$ is called its root. In Figure 1, the polyomino is Temperleyan with root the lower left (removed) square. All Temperleyan polyominos have domino tilings \[5, section 7\]. The term Temperleyan comes from the bijection due to Temperley between the set of spanning trees of a rectangle in $`^2`$ and the set of domino tilings of a rectangular region with a corner removed . This bijection was generalized in and further in . Let $`U`$ be a smooth Jordan domain with a marked point $`bU`$. For each $`ϵ>0`$ let $`P_ϵ`$ be a Temperleyan polyomino in $`ϵ^2`$ approximating $`U`$ as follows. The boundary of $`P_ϵ`$ lies within $`O(ϵ)`$ of $`U`$, and the counterclockwise boundary path of $`P_ϵ`$ points locally into the same half-space as the (directed) tangent to $`U`$ which it is near. Furthermore the root $`b_ϵ`$ of $`P_ϵ`$ should be within $`O(ϵ)`$ of $`b`$. ### 2.2 Green’s functions Let $`U`$ be a Jordan domain with basepoint $`bU`$. The Green’s function with Dirichlet boundary conditions, or simply Dirichlet Green’s function, $`g_D(z_1,z_2)`$, is defined to be the unique function (of $`z_2`$) satisfying $`\mathrm{\Delta }g_D(z_1,z_2)=\delta _{z_1}(z_2)`$ (the Dirac delta), and which is zero when $`z_2U`$, where the Laplacian is with respect to the second variable. This function is well-defined and when $`z_2`$ is near $`z_1`$ has the form $`g_D(z_1,z_2)=\frac{1}{2\pi }\mathrm{log}\|z_2z_1\|+O(1)`$. The Dirichlet Green’s function has the following simple expression in the basis of eigenfunctions of the Laplacian on $`U`$: ###### Lemma 2.1 $$g_D(z_1,z_2)=\underset{i1}{}\frac{f_i(z_1)f_i(z_2)}{\lambda _i}.$$ Proof. Since the eigenbasis $`\{f_i\}`$ of $`\mathrm{\Delta }`$ is an orthonormal basis for $`L^2(U)`$, it suffices to show that for each $`i`$, $`f_i(z_2),g_D(z_1,z_2)=\frac{f_i(z_1)}{\lambda _i}`$. But $`f_i(z_2),g_D(z_1,z_2)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _i}}\lambda _if_i(z_2),g_D(z_1,z_2)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _i}}\mathrm{\Delta }f_i(z_2),g_D(z_1,z_2)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _i}}f_i(z_2),\mathrm{\Delta }g_D(z_1,z_2)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _i}}f_i(z_2),\delta _{z_1}(z_2)`$ $`=`$ $`{\displaystyle \frac{1}{\lambda _i}}f_i(z_1).`$ $`\mathrm{}`$ We will also need to define the Neumann Green’s function $`g_N`$. On a bounded domain, the Green’s function with Neumann boundary conditions does not exist. However we can define (see below) the function which is the “difference of two Neumann Green’s functions”: for points $`z_1,z_1^{}U`$ define the difference of Neumann Green’s functions $`g_N(z_1,z_2)g_N(z_1^{},z_2)`$ to be the function of $`z_2`$ satisfying $`\mathrm{\Delta }(g_N(z_1,z_2)g_N(z_1^{},z_2))=\delta _{z_1}(z_2)\delta _{z_1^{}}(z_2)`$, where the Laplacian is with respect to the second variable, and which satisfies $`\frac{}{\widehat{n}}(g_N(z_1,z_2)g_N(z_1^{},z_2))=0,`$ that is, the derivative normal to the boundary with respect to the second variable is zero. This function is well-defined up to an additive constant, and we set the constant so that $`g(z_1,b)g(z_1^{},b)=0`$. Near $`z_1`$ we have $`g_N(z_1,z_2)g_N(z_1^{},z_2)=\frac{1}{2\pi }\mathrm{log}\|z_2z_1\|+c_1(z_1,z_1^{})+O(z_2z_1).`$ Let $`\widehat{g}_N(z_1,z_2)\widehat{g}_N(z_1^{},z_2)`$ be the harmonic conjugate (with respect to the second variable) of $`g_N(z_1,z_2)g_N(z_1^{},z_2)`$. This function is multiply-valued, increasing by $`1`$ when $`z_2`$ turns counterclockwise around $`z_1`$ and by $`1`$ when $`z_2`$ turns counterclockwise around $`z_1^{}`$. The function $`\stackrel{~}{g}_N(z_1,z_2)\stackrel{~}{g}_N(z_1^{},z_2):=g_N(z_1,z_2)g_N(z_1^{},z_2)+i(\widehat{g}_N(z_1,z_2)\widehat{g}_N(z_1^{},z_2))`$ is analytic in $`z_2`$ (except at $`z_1`$ and $`z_1^{}`$) and is the analytic Neumann Green’s function. It is also multiply-valued. We can define the exterior derivative of $`\stackrel{~}{g}_N`$ with respect to the first variable as $`d\stackrel{~}{g}_N(z_1,z_2)=F_0(z_1,z_2)dx_1+iF_1(x,z_2)dy_1`$, where $`z_1=x_1+iy_1`$ and $`F_0,F_1`$ are defined by taking limits (for $`\delta `$ real) $`F_0(z_1,z_2)=lim_{\delta 0}\frac{\stackrel{~}{g}_N(z_1+\delta ,z_2)\stackrel{~}{g}_N(z_1,z_2)}{\delta }`$ and $`iF_1(z_1,z_2)=lim_{\delta 0}\frac{\stackrel{~}{g}_N(z_1+i\delta ,z_2)\stackrel{~}{g}_N(z_1,z_2)}{\delta }`$. These functions $`F_0,F_1`$ are well-defined (single-valued) and vanish at $`z_2=b`$. As examples of these functions, on the upper half-plane $``$ with $`b=\mathrm{}`$ we have $$g_D(z_1,z_2)=\frac{1}{2\pi }\mathrm{log}\left\|\frac{z_2z_1}{z_2\overline{z}_1}\right\|,$$ (2) and $$\stackrel{~}{g}_N(z_1,z_2)\stackrel{~}{g}_N(z_1^{},z_2)=\frac{1}{2\pi }\mathrm{log}\frac{(z_2z_1)(z_2\overline{z}_1)}{(z_2z_1^{})(z_2\overline{z_1^{}})}.$$ (3) Note that when $`z_2`$ the imaginary part of (3) is constant (in fact vanishes); this implies that the real part has Neumann boundary conditions. For a more general Jordan domain $`V`$, let $`f`$ be a Riemann map from $`V`$ to the upper half-plane sending $`b`$ (the base point of $`V`$) to $`\mathrm{}`$. Then the Dirichlet Green’s function on $`V`$ is $`g_D^V(z_1,z_2)=g_D^{}(f(z_1),f(z_2))`$, and the analytic Neumann Green’s function is defined similarly $`\stackrel{~}{g}_N^V(z_1,z_2)\stackrel{~}{g}_N^V(z_1^{},z_2)=\stackrel{~}{g}_N^{}(f(z_1),f(z_2))\stackrel{~}{g}_N^{}(f(z_1^{}),f(z_2))`$. One can in fact take this to be the definition of the Green’s functions on $`V`$. ### 2.3 Moment formula For a region $`U`$ with basepoint $`bU`$, define the functions $`F_+(z_1,z_2)`$ and $`F_{}(z_1,z_2)`$ by $$4d\stackrel{~}{g}_N(z_1,z_2)=F_+(z_1,z_2)dz_1+F_{}(z_1,z_2)d\overline{z_1},$$ (4) where $`d`$ is exterior differentiation with respect to the first variable. Then $`F_+(z_1,b)=0=F_{}(z_1,b)`$ (in terms of the functions $`F_0,F_1`$ of the previous section we have $`F_\pm =2(F_0\pm F_1)`$). Let $`h_0(x)=h(x)\overline{h}(x)`$. ###### Proposition 2.2 () Under the hypotheses of Theorem 1.1, let $`z_1,\mathrm{},z_k`$ be distinct points of $`U`$, and $`\gamma _1,\mathrm{},\gamma _k`$ disjoint paths running from the boundary of $`U`$ to $`z_1,\mathrm{},z_k`$ respectively. Let $`h(z_1)`$ denote the height of a point of $`P_ϵ`$ lying within $`O(ϵ)`$ of $`z_1`$. Then $$\underset{ϵ0}{lim}𝔼(h_0(z_1)\mathrm{}h_0(z_k))=\underset{\epsilon _1,\mathrm{},\epsilon _k\{\pm 1\}}{}\epsilon _1\mathrm{}\epsilon _k_{\gamma _1}\mathrm{}_{\gamma _k}\underset{i,j[1,k]}{det}\left(F_{\epsilon _i,\epsilon _j}(z_i,z_j)\right)dz_1^{(\epsilon _1)}\mathrm{}dz_k^{(\epsilon _k)},$$ (5) where $`dz_j^{(1)}=dz_j`$ and $`dz_j^{(1)}=d\overline{z_j}`$, and $$F_{\epsilon _i,\epsilon _j}(z_i,z_j)=\{\begin{array}{cc}0\hfill & \text{if }i=j\hfill \\ F_+(z_i,z_j)\hfill & \text{if }(\epsilon _i,\epsilon _j)=(1,1)\hfill \\ F_{}(z_i,z_j)\hfill & \text{if }(\epsilon _i,\epsilon _j)=(1,1)\hfill \\ \overline{F_{}(z_i,z_j)}\hfill & \text{if }(\epsilon _i,\epsilon _j)=(1,1)\hfill \\ \overline{F_+(z_i,z_j)}\hfill & \text{if }(\epsilon _i,\epsilon _j)=(1,1).\hfill \end{array}$$ ## 3 Proof of Theorem 1.1 When $`U`$ is the upper half plane with basepoint at $`\mathrm{}`$, the derivative of the analytic Neumann Green’s function is (see (3)) $$d\stackrel{~}{g}_N(z_1,z_2)=\frac{dz_1}{2\pi (z_1z_2)}+\frac{d\overline{z_1}}{2\pi (\overline{z_1}z_2)}.$$ Thus from (4) we have $`F_+(z_1,z_2)=\frac{2}{\pi (z_2z_1)}`$ and $`F_{}(z_1,z_2)=\frac{2}{\pi (z_2\overline{z_1})}`$. Let $`p,qU`$. From Proposition 2.2 we have $`lim_{ϵ0}𝔼(h_0(p)h_0(q))=`$ $$=_{\gamma _1,\gamma _2}\left\|\begin{array}{cc}0& F_+(z_1,z_2)\\ F_+(z_2,z_1)& 0\end{array}\right\|𝑑z_1𝑑z_2_{\gamma _1,\gamma _2}\left\|\begin{array}{cc}0& F_{}(z_1,z_2)\\ \overline{F_{}(z_2,z_1)}& 0\end{array}\right\|𝑑\overline{z_1}𝑑z_2$$ $$_{\gamma _1,\gamma _2}\left\|\begin{array}{cc}0& \overline{F_{}(z_1,z_2)}\\ F_{}(z_2,z_1)& 0\end{array}\right\|𝑑z_1𝑑\overline{z_2}+_{\gamma _1,\gamma _2}\left\|\begin{array}{cc}0& \overline{F_+(z_1,z_2)}\\ \overline{F_+(z_2,z_1)}& 0\end{array}\right\|𝑑\overline{z_1}𝑑\overline{z_2}.$$ Plugging in for $`F_\pm `$ gives $$\underset{ϵ0}{lim}𝔼(h_0(p)h_0(q))=\frac{4}{\pi ^2}_{\gamma _1}_{\gamma _2}\frac{1}{(z_2z_1)^2}𝑑z_1𝑑z_2+\frac{4}{\pi ^2}_{\gamma _1}_{\gamma _2}\frac{1}{(z_2\overline{z_1})^2}𝑑\overline{z_1}𝑑z_2+$$ $$+\frac{4}{\pi ^2}_{\gamma _1}_{\gamma _2}\frac{1}{(\overline{z_2}z_1)^2}𝑑z_1𝑑\overline{z_2}\frac{4}{\pi ^2}_{\gamma _1}_{\gamma _2}\frac{1}{(\overline{z_2}\overline{z_1})^2}𝑑\overline{z_1}𝑑\overline{z_2}$$ $$=\frac{8}{\pi ^2}\text{Re}\mathrm{log}\left(\frac{\overline{p}q}{pq}\right).$$ (6) Note that this is $`\frac{16}{\pi }g_D(p,q)`$ where $`g_D`$ is the Dirichlet Green’s function on $`U`$ (see (2)). Now let $`p_1,\mathrm{},p_k`$ be distinct points in the upper half plane $`U`$. By symmetry of $`h_0`$, if $`k`$ is odd the moment $`𝔼(h_0(p_1)\mathrm{}h_0(p_k))`$ is zero. We therefore assume $`k`$ is even. In (5), the matrix has $`ij`$ entry $$F_{\epsilon _i,\epsilon _j}(z_i,z_j)=\frac{2}{\pi (z_j^{(\epsilon _j)}z_i^{(\epsilon _i)})}.$$ Such a matrix has a simple determinant: ###### Lemma 3.1 For $`k`$ even let $`M`$ be the $`k\times k`$ matrix $`M=(m_{ij})`$ with $`m_{ii}=0`$ and $`m_{ij}=\frac{1}{x_jx_i}`$ when $`ij`$. Then $$det(M)=\frac{1}{(x_{\sigma (1)}x_{\sigma (2)})^2(x_{\sigma (3)}x_{\sigma (4)})^2\mathrm{}(x_{\sigma (k1)}x_{\sigma (k)})^2},$$ (7) where the sum is over all $`(k1)!!`$ possible pairings $`\{\{\sigma (1),\sigma (2)\},\mathrm{},\{\sigma (k1),\sigma (k)\}\}`$ of $`\{1,\mathrm{},k\}`$. This lemma also appears in . Proof. The proof is by induction on $`k`$. The formula clearly holds when $`k=2`$. For $`k>2`$, the determinant is a rational function of $`x_1`$ with a double pole at $`x_1=x_2`$; we can write $$det(M)=\frac{c_2}{(x_1x_2)^2}+\frac{c_1}{(x_1x_2)}+c_0+O(x_1x_2).$$ The coefficient $`c_1`$ is zero since the determinant is even under the exchange of $`x_1`$ and $`x_2`$ (exchange the first two rows and exchange the first two columns). The coefficient $`c_2`$ is the determinant of $`M_{12}`$, the matrix obtained from $`M`$ by deleting the first two rows and columns. Therefore the right and left-hand sides of (7) both represent rational functions (in each variable) with the same poles and residues; hence they differ by a constant. This constant is zero by homogeneity. $`\mathrm{}`$ Combining the lemma with Proposition 2.2 gives the following. ###### Proposition 3.2 Let $`U`$ be a Jordan domain with smooth boundary. Let $`p_1,\mathrm{},p_kU`$ (with $`k`$ even) be distinct points. We have $$\underset{ϵ0}{lim}𝔼(h_0(p_1)\mathrm{}h_0(p_k))=\left(\frac{16}{\pi }\right)^{k/2}\underset{\text{pairings }\sigma }{}g_D(p_{\sigma (1)},p_{\sigma (2)})\mathrm{}g_D(p_{\sigma (k1)},p_{\sigma (k)}).$$ Proof. When $`U`$ is the upper half plane this follows by combining Proposition 2.2 with Lemma 3.1 and the calculation (6). For arbitrary $`U`$, equation (6) shows that $`𝔼(h_0(p_1)h_0(p_2))=\frac{16}{\pi }g_D^U(p_1,p_2)`$ (where $`g_D^U`$ is the Dirichlet Green’s function on $`U`$) by conformal invariance of the height moments and of $`g_D`$. This completes the proof. $`\mathrm{}`$ The proof of Theorem 1.1 is completed as follows. Let $`f_{n_1},\mathrm{},f_{n_k}`$ be (not necessarily distinct) eigenvectors of $`\mathrm{\Delta }`$ with Dirichlet boundary conditions. Let $`C_{n_j}^{(ϵ)}`$ be the real-valued random variable $`C_{n_j}^{(ϵ)}=ϵ^2_{xV_ϵ}h_0(x)f_{n_j}(x)`$, where the sum is over the vertices $`V_ϵ`$ of $`P_ϵ`$, and $`f_{n_j}(x)`$ is $`f_{n_j}`$ evaluated at the vertex $`x`$. We have $`\underset{ϵ0}{lim}𝔼(C_{n_1}^{(ϵ)}\mathrm{}C_{n_k}^{(ϵ)})=`$ $`=`$ $`\underset{ϵ0}{lim}𝔼\left({\displaystyle \underset{x_1V_ϵ}{}}ϵ^2h_0(x_1)f_{n_1}(x_1)\mathrm{}{\displaystyle \underset{x_kV_ϵ}{}}ϵ^2h_0(x_k)f_{n_k}(x_k)\right)`$ $`=`$ $`\underset{ϵ0}{lim}{\displaystyle \underset{x_1V_ϵ}{}}\mathrm{}{\displaystyle \underset{x_kV_ϵ}{}}ϵ^2f_{n_1}(x_1)\mathrm{}ϵ^2f_{n_k}(x_k)𝔼\left(h_0(x_1)\mathrm{}h_0(x_k)\right)`$ $`=`$ $`\left({\displaystyle \frac{16}{\pi }}\right)^{k/2}{\displaystyle _U}\mathrm{}{\displaystyle _U}f_{n_1}(x_1)\mathrm{}f_{n_k}(x_k){\displaystyle \underset{\sigma }{}}g_D(x_{\sigma (1)},x_{\sigma (2)})\mathrm{}g_D(x_{\sigma (k1)},x_{\sigma (k)})`$ $`=`$ $`\left({\displaystyle \frac{16}{\pi }}\right)^{k/2}{\displaystyle \underset{\sigma }{}}{\displaystyle _U}\mathrm{}{\displaystyle _U}f_{n_1}(x_1)\mathrm{}f_{n_k}(x_k){\displaystyle \underset{m_1,\mathrm{},m_{k/2}}{}}{\displaystyle \frac{f_{m_1}(x_{\sigma (1)})f_{m_1}(x_{\sigma (2)})}{\lambda _{m_1}}}\mathrm{}{\displaystyle \frac{f_{m_{k/2}}(x_{\sigma (k1)})f_{m_{k/2}}(x_{\sigma (k)})}{\lambda _{m_{k/2}}}}`$ $`=`$ $`\left({\displaystyle \frac{16}{\pi }}\right)^{k/2}{\displaystyle \underset{\sigma }{}}{\displaystyle \frac{\delta _{\sigma (1),\sigma (2)}}{(\lambda _{\sigma (1)})}}\mathrm{}{\displaystyle \frac{\delta _{\sigma (k1),\sigma (k)}}{(\lambda _{\sigma (k1)})}}.`$ By Wick’s theorem , these are exactly the moments for a set of independent Gaussians of mean zero and variances $`\frac{16}{\pi \lambda _i}`$. Now to conclude we invoke the following standard probability lemma : ###### Lemma 3.3 () A sequence of (multidimensional) random variables whose moments converge to the moments of a Gaussian, converges itself to a Gaussian. This completes the proof. ## 4 Proof of Proposition 1.2 Since $`X`$ and $`Y`$ are Gaussians (each being the sum of Gaussians), and have mean $`0`$, it suffices to compute their variances. But $`𝔼(X^2)`$ $`=`$ $`{\displaystyle _U}{\displaystyle _U}\omega (z_1)\omega (z_2)𝔼(F(z_1)F(z_2))`$ $`=`$ $`{\displaystyle _U}{\displaystyle _U}\omega (z_1)\omega (z_2)g_D^U(z_1,z_2)`$ $`=`$ $`{\displaystyle _V}{\displaystyle _V}f^{}\omega (y_1)f^{}\omega (y_2)g_D^U(f(y_1),f(y_2))`$ $`=`$ $`{\displaystyle _V}{\displaystyle _V}f^{}\omega (y_1)f^{}\omega (y_2)g_D^V(y_1,y_2)`$ $`=`$ $`𝔼(Y^2)`$ where we used the conformal invariance of the Green’s function, $`g_D^U(f(y_1),f(y_2))=g_D^V(y_1,y_2)`$. This completes the proof.
warning/0002/astro-ph0002190.html	ar5iv	text	# Constraining the Lifetime of Quasars from their Spatial Clustering ## 1. Introduction A long outstanding problem in cosmology is the synchronized evolution of the quasar population over the redshift range $`0<z<\mathrm{\hspace{0.33em}5}`$. Observations in the optical (Pei 1995) and radio (Shaver et al. 1994) show a pronounced peak in the abundance of bright quasars at $`z2.5`$; recent X–ray observations (Miyaji et al. 2000) confirm the rapid rise from $`z=0`$ towards $`z2`$, but have not shown evidence for a decline at still higher redshifts. Individual quasars are widely understood to consist of supermassive black holes (BHs) powered by accretion (Lynden-Bell 1967; Rees 1984). A plausible timescale for quasar activity is then the Eddington time, $`4\times 10^7`$ ($`ϵ`$/0.1) yr, the e-folding time for the growth of a BH accreting mass at a rate $`\dot{M}`$, while shining at the Eddington luminosity with a radiative efficiency of $`L=L_{\mathrm{Edd}}=ϵ\dot{M}c^2`$. The lifetime $`t_Q`$ of the luminous phase of quasars can be estimated directly, by considering the space density of quasars and galaxies. At $`z2`$, the ratio $`n_Q/n_G3\times 10^3`$ implies the reassuringly close value of $`t_Qt_{\mathrm{Hub}}n_Q/n_G10^7`$ yr (Blandford 1999 and references therein). These lifetimes are significantly shorter than the Hubble time, suggesting that the quasar population evolves on cosmic time–scales by some mechanism other than local accretion physics near the BH. It is tempting to identify quasars with halos condensing in a cold dark matter (CDM) dominated universe, as the halo population naturally evolves on cosmic time–scales (Efstathiou & Rees 1988; Haiman & Loeb 1998; Kauffmann & Haehnelt 2000). Furthermore, quasars reside in a subset of all galaxies, while the redshift–evolution of the galaxy population as a whole (qualitatively similar to that of bright quasars) has been successfully described by associating galaxies with dark halos (e.g. Lacey & Cole 1993; Kauffmann & White 1993). A further link between galaxies and quasars comes from the recent detection, and measurements of the masses of massive BHs at the centers of nearby galaxies (Magorrian et al. 1998; van der Marel 1999). These arguments suggest that the evolution of the quasar population can indeed be described by “semi–analytic” models, associating quasars with dark matter halos. In this type of modeling, the quasar lifetime plays an important role. The quasar phase in a single halo could last longer ($`t_\mathrm{Q}10^8`$yr), with correspondingly small $`M_{\mathrm{bh}}/M_{\mathrm{halo}}`$ ratios, or last shorter ($`t_\mathrm{Q}10^6`$yr), with larger BH formation efficiencies (Haiman & Loeb 1998; Haehnelt et al. 1998). Note that although recent studies have established a correlation between the bulge mass $`M_{\mathrm{bulge}}`$ and BH mass $`M_{\mathrm{bh}}`$, this correlation leaves a considerable uncertainty in the relation between $`M_{\mathrm{bh}}`$ and the mass $`M_{\mathrm{halo}}`$ of its host halo. If the initial density fluctuations are Gaussian with a CDM power spectrum, the clustering of collapsed halos is a function of their mass – rarer, more massive halos cluster more strongly (Kaiser 1984; Mo & White 1996). Hence, measurements of quasar clustering are a potentially useful probe of both BH formation efficiencies and quasar lifetimes (La Franca et al. 1998, Haehnelt et al. 1998). In this paper, we assess the feasibility of breaking the above degeneracy, and inferring quasar lifetimes, from the statistics of clustering that will be available from the Sloan Digital Sky Survey (SDSS, Gunn & Weinberg 1995) and Anglo-Australian Telescope Two-Degree-Field (2dF, Boyle et al. 1999). Previous works (e.g. Stephens et al. 1997; Sabbey et al. 1999) have yielded estimates suggesting that quasars are clustered more strongly than galaxies. However, the current uncertainties are large, especially at higher redshifts, where clustering has been found to decrease (Iovino & Shaver 1988; Iovino et al. 1991), to stay constant (Andreani & Cristiani 1992; Croom & Shanks 1996), or to increase with redshift (La Franca et al. 1998). As a result, no strong constraints on the life–time can be obtained yet. The key advance of forthcoming surveys over previous efforts is two–fold. Because of their sheer size, i.e. the large number of quasars covering a large fraction of the sky, both shot–noise and sample variance can be beaten down, significantly reducing the statistical uncertainties. Furthermore, the large sample-size will eliminate the need to combine data from different surveys with different selection criteria, hence allowing cleaner interpretation. Recent measurements of the local massive black hole density have stimulated discussions of a radiative efficiency which is much lower than the usual $`0.1`$ (e.g. Haehnelt et al. 1999). A convincing constraint on the lifetime of quasars could be therefore highly interesting, as this might have implications for the local accretion physics near the BH. This paper is organized as follows. In § 2, we summarize our models for the quasar luminosity function, and in § 3 we compute the quasar correlation function in these models. In § 4, we compare the model predictions with presently available data, and in § 5 we assess the ability of future optical redshift surveys to discriminate between the various models. In § 6, we repeat our analysis in the soft X–ray band, and examine the contribution of quasars to the X-ray background, and its auto–correlation. In § 7, we address some of the caveats arising from our assumptions, and in § 8 we summarize our conclusions and the implications of this work. ## 2. Models for the Quasar Luminosity Function In this section, we briefly summarize our model for the luminosity function (LF) of quasars, based on associating quasar BHs with dark halos. Our treatment is similar to previous works (Haiman & Loeb 1998; Haehnelt et al. 1998), although differs in some of the details. A more extensive treatment is provided in the Appendix. The main assumption is that there is, on average, a direct monotonic relation between halo mass $`M_{\mathrm{halo}}`$ and average quasar luminosity $`L_{M,z}`$, which we parameterize using the simple power–law ansatz: $$\overline{L}_{M,z}=x_0(z)M_{\mathrm{halo}}\left(\frac{M_{\mathrm{halo}}}{M_0}\right)^{\alpha (z)}.$$ (1) Here $`x_0(z)`$ and $`\alpha (z)`$ are “free functions”, whose values are found by the requirement that the resulting luminosity function agrees with observations. As explained in the Appendix, our model has one free parameter, the lifetime $`t_Q`$, which uniquely determines $`x_0(z)`$ and $`\alpha (z)`$ in any given background cosmology. We assume the background cosmology to be either flat ($`\mathrm{\Lambda }`$CDM) with $`(\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_\mathrm{m},h,\sigma _{8\mathrm{h}^1},n)=(0.7,0.3,0.65,1.0,1.0)`$ or open (OCDM) with $`(\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_\mathrm{m},h,\sigma _{8\mathrm{h}^1},n)=(0,0.3,0.65,0.82,1.3)`$. In LCDM, we find $`(\mathrm{log}[x_0/\mathrm{L}_{}\mathrm{M}_{}^1],\alpha )(1,0.4)`$ and $`(0.2,0.1)`$ for lifetimes of $`t_Q=10^8`$ and $`t_Q=10^{6.5}`$yr, respectively. Similarly, in OCDM, we find $`(\mathrm{log}[x_0/\mathrm{L}_{}\mathrm{M}_{}^1],\alpha )(1,0.25)`$ and $`(0.2,0.25)`$ for these two lifetimes. In Figure 1, we demonstrate the agreement between the LF computed in our models with the observational data at two different redshifts, $`z=2`$ and $`z=3`$. For reference, the upper labels in this figure show the apparent magnitudes in the SDSS $`g^{}`$ band, assuming that the intrinsic quasar spectrum is the same as the mean spectrum in the Elvis et al. (1994) quasar sample. The photometric detection threshold<sup>1</sup><sup>1</sup>1See http://www.sdss.org/science/tech\_summary.html. of SDSS is $`g^{}22.6`$, corresponding to a BH mass of $`10^8\mathrm{M}_{}`$ at $`z=3`$ and a three times smaller mass at $`z=2`$. As the figure shows, the overall quality of the fits is excellent; for reference, the dashed lines show the ad–hoc empirical fitting formulae from Pei (1995). Similarly accurate match to the quasar LF is achieved at different redshifts, and in the models assuming an OCDM cosmology. Figure 1 shows, in particular, that the fits obtained from the power–law ansatz adopting either a short (solid lines) or a long (dotted lines) lifetime are nearly indistinguishable; hence modeling the LF by itself does not constrain the quasar lifetime within the limits $`10^{6.5}`$ yr $`<t_Q<\mathrm{\hspace{0.33em}10}^8`$ yr. Before considering constraints on the lifetime from clustering, it is useful to point out that estimates for both upper and lower limits on $`t_Q`$ follow from the observed luminosity function alone. Lower limit on $`t_Q`$. A halo of mass $`M_{\mathrm{halo}}`$ is unlikely to harbor a BH more massive than $`6\times 10^3(\mathrm{\Omega }_b/\mathrm{\Omega }_0)M_{\mathrm{halo}}=6\times 10^4M_{\mathrm{halo}}`$, where $`6\times 10^3`$ is the ratio $`M_{\mathrm{bh}}/M_{\mathrm{bulge}}`$ found in nearby galaxies (Magorrian et al. 1998) because $`M_{\mathrm{bulge}}`$ cannot be larger than $`(\mathrm{\Omega }_b/\mathrm{\Omega }_0)M_{\mathrm{halo}}`$. This maximal BH could at best emit $`10\%`$ of the Eddington luminosity in the B–band, implying $`L/M_{\mathrm{halo}}<\mathrm{\hspace{0.33em}3}\mathrm{L}_{}/\mathrm{M}_{}`$. We find (cf. Fig. 7) that our short lifetime model with $`t_Q=10^{6.5}`$ yr nearly reaches this limit; models with shorter lifetimes would require unrealistically large $`L/M_{\mathrm{halo}}`$ ratios. Upper limit on $`t_Q`$. Long lifetimes, on the other hand, require the ratio $`L/M_{\mathrm{halo}}`$ to be small; this can lead to unrealistically large halo masses. The brightest quasars detected at redshifts $`z23`$ have luminosities as large as $`L10^{14}\mathrm{L}_{}`$ (Pei 1995). We find (cf. Fig. 7) that in order to avoid the host halo masses of these bright quasars to exceed $`10^{15}\mathrm{M}_{}`$, the lifetime cannot be longer than $`10^8`$ yr. An alternative, standard argument goes as follows. The black hole mass grows during the quasar phase as $`e^{t_Q/t_E}`$ where $`t_E=4\times 10^7(ϵ/0.1)`$ yr is the Eddington time. Assuming a conservative initial black hole mass of $`1M_{}`$, a lifetime longer than $`10^9`$ yr would give final black hole masses that are unacceptably large. ## 3. The Clustering of Quasars As demonstrated in the previous section, equally good fits can be obtained to the luminosity function of quasars, assuming either a short or a long lifetime, and a power-law dependence of the mean quasar luminosity on the halo mass $`\overline{L}M^\alpha `$. In this section, we derive the clustering of quasars in our models, and demonstrate that they depend significantly on the assumed lifetime. The halos are a biased tracer of the underlying mass distribution, customarily expressed by $`P_{\mathrm{halo}}(k)=b^2P(k)`$ where $`P_{\mathrm{halo}}`$ and $`P`$ are the halo and mass power spectra as a function of wavenumber $`k`$. The bias parameter for halos of a given mass $`M`$ at a given redshift $`z`$ is given by (Mo & White 1996) $$b(M,z)=1+\frac{1}{\delta _c}\left[\left(\frac{\delta _c}{\sigma (M)D(z)}\right)^21\right],$$ (2) where $`D(z)`$ is the linear growth function, $`\sigma (M)`$ is the r.m.s. mass fluctuation on mass–scale $`M`$ (using the power spectrum of Eisenstein & Hu 1999), and $`\delta _c1.68`$ is the usual critical overdensity in the Press–Schechter formalism (see Jing 1999 and Sheth & Tormen 1999 for more accurate expressions for $`b(M,z)`$ for low $`M`$, which we find not to affect our results here). The bias associated with quasars with luminosity $`L`$ in our models is given by averaging over halos of different masses associated with this luminosity. Following equation 19, we obtain $`b(L,z)=`$ $`\left[{\displaystyle \frac{d\varphi }{dL}}(L,z)\right]^1\times {\displaystyle _0^{\mathrm{}}}𝑑M{\displaystyle \frac{dN}{dM}}(M,z)`$ $`b(M,z){\displaystyle \frac{dg}{dL}}(L,\overline{L}_{M,z})f_{\mathrm{on}}(M,z).`$ We show in Figure 2 the resulting bias parameter $`b(L,z)`$ in the models corresponding to Figure 1, with short and long lifetimes, and at redshifts $`z=2`$ and 3. As expected, quasars are more highly biased in the long lifetime model, by a ratio $`b`$(long)/$`b`$(short)$`>\mathrm{\hspace{0.33em}2}`$. In the $`\mathrm{\Lambda }`$CDM case, at the detection threshold of SDSS, we find $`b`$(long)$`3`$ at $`z=3`$ and $`b`$(long)$`2`$ at $`z=2`$. Bright quasars with $`g^{}=17`$ are predicted to have a bias at $`z=3`$ as large as $`b=10`$. For reference, we also show in this figure the bias parameters obtained in the OCDM cosmology, which are significantly lower than in the $`\mathrm{\Lambda }`$CDM case. The number of quasars observed at a fixed flux implies an intrinsically larger number of sources if OCDM is assumed, because the volume per unit redshift and solid angle in an open universe is smaller. This lowers the corresponding halo mass and therefore the bias. ## 4. Comparison with Available Data As emphasized in § 1, the presently available data leave considerable uncertainties in the clustering of quasars. Nevertheless, it is interesting to contrast the results of the previous section with preliminary results from the already relatively large, homogeneous sample of high–redshift quasars in the 2dF survey (Boyle et al. 2000). Our predictions are obtained by relating the apparent magnitude limit to a minimum absolute luminosity at a given redshift: $`\mathrm{log}[L_{\mathrm{min}}(z)/L_{B,}]=0.4[5.48B+5\mathrm{log}(d_\mathrm{L}(z)/\mathrm{pc})5]`$. This relation assumes no K–correction, justified by the nearly flat quasar spectra ($`\nu F_\nu =`$ const) at the relevant wavelengths (e.g. Elvis et al. 1994; Pei 1995). In our model, the correlation length $`r_0`$ is given implicitly by $$\xi _q(r_0)\overline{b}^2(z)D^2(z)\xi _m(r_0)=1,$$ (4) where $`\xi _m(r)`$ is the usual dark matter correlation function, and $`\overline{b}(z)`$ is the value of the bias parameter $`b(L,z)`$ as determined in the previous section, but now averaged over all quasars with magnitudes brighter than the detection limit, $$\overline{b}(z)=\left[_{L_{\mathrm{min}}(z)}^{\mathrm{}}𝑑L\frac{d\varphi }{dL}\right]^1_{L_{\mathrm{min}}(z)}^{\mathrm{}}𝑑L\frac{d\varphi }{dL}b(L,z).$$ (5) In Figure 3 we show the resulting correlation lengths in the long and short lifetime models. Also shown is a preliminary data–point with 1$`\sigma `$ error–bars from the 2dF survey, based on $``$3000 quasars with apparent magnitudes $`B<20.85`$ (Croom et al. 1999). The upper panel shows the results in our fiducial $`\mathrm{\Lambda }`$CDM model with predictions for this magnitude cut. The published results for $`r_0`$ are cosmology dependent, and we have simply converted them for our cosmological models by taking the corresponding average of the redshift-distance and angular-diameter-distance – this crude treatment is adequate given the large measurement errors. Our models generically predict a gradual increase of the correlation length with redshift (“positive evolution”). The clustering is dominated by the faintest quasars near the threshold luminosity; as a result, the fixed magnitude-cut of $`B=20.85`$ corresponds to more massive, and more highly clustered halos at higher redshifts. There are two additional effects that determine the redshift–evolution of clustering: (1) quasars of a fixed luminosity are more abundant towards high–$`z`$, requiring smaller halo masses to match their number density, and (2) halos with a fixed mass are more highly clustered towards high-$`z`$. We find, however, that these effects are less important than the increase in $`M_{\mathrm{halo}}`$ caused by fixing the apparent magnitude threshold, which gives rise to the overall positive evolution. The clustering in the long–lifetime model is stronger, and evolves more rapidly than in the short–lifetime case. As we can see, the present measurement error-bars are large – the whole range of life-times from $`10^{6.5}`$ to $`10^8`$ yr is broadly consistent with the data, to within $`2\sigma `$. For the $`\mathrm{\Lambda }`$CDM model, the 2dF data-point is consistent with a lifetime of around $`10^{7.7\pm 0.8}`$ yr; in the OCDM model the lifetime is somewhat higher, $`10^{8\pm 0.8}`$ yr. It is worth emphasizing here that the constraints on lifetime from clustering depends on the underlying cosmological parameters one assumes, which future large scale structure measurements (from e.g. the microwave background, galaxy surveys and the Lyman–$`\alpha `$ forest) will hopefully pin down to the accuracy required here. Lastly, we note that observational results on $`r_0`$ are commonly obtained by a fit to the two-point correlation of the form $`\xi (r)=(r/r_0)^\gamma `$. Since $`r_0`$ is the correlation length where the correlation is unity, we expect our formalism to begin to fail on such a scale, because neither linear fluctuation growth nor linear biasing holds. On the other hand, it is also unclear whether $`r_0`$ as presently measured from a 2-parameter fit to the still rather noisy observed two-point correlation necessarily corresponds to the true correlation length. While our crude comparison with existing data in Figure 3 suffices given the large measurement errors, superior data in the near future demand a more refined treatment, which is the subject of the next section. ## 5. Expectations from the SDSS and 2dF Although existing quasar clustering measurements still allow a wide range of quasar lifetimes, and do not provide tight constraints on our models, forthcoming large quasar samples from SDSS or the complete 2dF survey are ideally suited for this purpose. Here we estimate the statistical uncertainties on the derived lifetimes, using the three–dimensional quasar power spectrum $`P_Q(k)`$ derived from these surveys. The variance of the power spectrum is computed by $$\delta P_\mathrm{Q}^2(k)=n_k^1[\overline{b}^2P(k)+\overline{n}]^2,$$ (6) where the large scale fluctuations are assumed to be Gaussian, $`n_k`$ is the number of independent modes, $`\overline{n}`$ is the mean number density of observed quasars, $`P_Q(k)=\overline{b}^2P(k)`$ is the quasar power spectrum and $`P(k)`$ is the mass power spectrum (Feldman, Kaiser & Peacock 1994). For a survey of volume $`V`$, and a $`k`$-bin of size $`\mathrm{\Delta }k`$, we use $`n_k=k^2\mathrm{\Delta }kV/4\pi ^2`$. The fractional variance is therefore $`n_k^1\{1+1/[\overline{n}\overline{b}P(k)]\}^2`$. In terms of minimizing this error, increasing the luminosity cut of a survey has the advantage of raising the bias $`\overline{b}`$, but has the disadvantage of decreasing the abundance $`\overline{n}`$. In practice, $`\overline{b}`$ changes relatively slowly with mass (slower than $`\overline{b}M`$) whereas $`\overline{n}`$ varies with mass much more rapidly ($`\overline{n}1/M`$, or steeper if $`M>M_{}`$). As a result, we find that for our purpose of determining the clustering and the quasar lifetime, it is better to include more (fainter) quasars. The power spectrum $`P_Q(k)`$ of quasars in $`\mathrm{\Lambda }`$CDM is shown in Figure 4 at two different redshifts near the peak of the comoving quasar abundance, $`z=2`$ and $`z=3`$. We assume that redshift slices are taken centered at each redshift with a width of $`\mathrm{\Delta }z=0.5`$ (which enters into the volume $`V`$ above). Results are shown in the long and short lifetime models, together with the expected $`1\sigma `$ error bars from SDSS (crosses). Also shown in the lower panel are the expected error bars from 2dF (open squares), which are slightly larger because of the smaller volume (for SDSS, we assume an angular coverage of $`\pi `$ steradians, and for 2dF, an area of $`0.23`$ steradians). We do not show error bars for 2dF beyond $`k0.01\mathrm{Mpc}/\mathrm{h}`$ because larger scales would likely be affected by the survey window. We also only show scales where the mass power spectrum and biasing are believed to be linear. As these figures show, the long and short $`t_Q`$ models are easily distinguishable with the expected uncertainties both from the SDSS and the 2dF data, out to a scale of $`100`$Mpc. The measurement errors at different scales are independent (under the Gaussian assumption), and hence, when combined, give powerful constraints e.g. formally, even models with lifetimes differing by a few percent can be distinguished with high confidence using the SDSS. However, systematic errors due to the theoretical modeling are expected to be important at this level, which we will discuss in §7. In Figure 4, we have used the magnitude cuts for spectroscopy, i.e. $`B=20.85`$ for 2dF and $`B=20.4`$ ($`g^{}19`$) for SDSS. If photometric redshifts of quasars are sufficiently accurate, the magnitude cuts can be pushed fainter, further decreasing the error-bars – although it is likely that the photometric redshift errors will be large enough that one can only measure the clustering in projection, in some well-defined redshift-bin picked out using color information. Color selection of $`z>3`$ quasars has already proven to be highly effective (Fan et al. 1999); a high redshift sample can give valuable information on the evolution of the quasar clustering (see Fig. 3). ## 6. Quasar Clustering in X-ray Both the luminosity function (e.g. Miyaji et al. 2000), and the clustering of quasars (e.g. Carrera et al. 1998) has been studied in the X–ray band, analogously to the optical observations described above. At present, the accuracy of both quantities are inferior to that in the optical. Nevertheless, it is interesting to consider the clustering of quasars in the X–ray regime, because (1) applying the exercise outlined above to a different wavelength band provides a useful consistency check on our results, (2) observations with the Chandra X–ray Observatory (CXO) and XMM can potentially<sup>2</sup><sup>2</sup>2 see http://asc.harvard.edu and http://xmm.vilspa.esa.es, respectively probe quasars at redshifts higher than currently reached in the optical, and (3) X-ray observations are free from complications due to dust extinction, although other forms of obscuration are possible (see §7). In addition, we will consider here the auto–correlation of the soft X–ray background (XRB) as another potential probe of quasar clustering and lifetime. The formalism we presented in § 2 and § 3 is quite general, and we here apply it to the soft X–ray luminosity function (XRLF) from Miyaji et al. (2000). The details of the fitting procedure are given in the Appendix. In analogy with the optical case, we find that the clustering of X-ray selected quasars depends strongly on the lifetime. As an example, including all quasars whose observed flux is above $`3\times 10^{14}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, we find the correlation length at $`z=2`$ to be $`4h^1`$Mpc in the short lifetime, and $`11h^1`$Mpc in the long lifetime case ($`\mathrm{\Lambda }`$CDM). Current data probes the clustering of X-ray quasars only at low redshifts ($`z<\mathrm{\hspace{0.33em}1}`$, see Carrera et al. 1998), where our models suffer from significant uncertainties due to the subhalo problem discussed in §7. Constraints at $`z>\mathrm{\hspace{0.33em}2}`$ could be available in the future from CXO and XMM, provided that a large area of the sky is surveyed at the improved sensitivities of these instruments. We next focus on the quasar contribution to the soft X–ray background and its auto–correlation, which as we will see is dominated by quasar contributions at somewhat higher redshifts. The mean comoving emissivity at energy $`E`$ from all quasars at redshift $`z`$, typically in units of $`\mathrm{keV}\mathrm{cm}^3\mathrm{s}^1\mathrm{sr}^1`$, is given in our models by $$\overline{j}(E,z)=\frac{1}{4\pi }_0^{\mathrm{}}𝑑L\frac{d\varphi }{dL}L_X(E,L),$$ (7) where $`L_X(E,L)`$ is the luminosity (in $`\mathrm{keV}\mathrm{s}^1\mathrm{keV}^1`$) at the energy $`E`$ of a quasar whose luminosity at $`(1+z)`$ keV is $`L`$. We have used the mean spectrum of Elvis et al. (1994) to include a small K–correction when computing the background at observed energies $`E1`$keV. The mean background is the integral of the emissivity over redshift, $$\overline{I}(E)=_0^{\mathrm{}}\frac{d\chi }{(1+z)}\overline{j}(E_z,\chi ),$$ (8) where $`\overline{I}`$ is typically given in units of $`\mathrm{keV}\mathrm{cm}^2\mathrm{s}^1\mathrm{sr}^1\mathrm{keV}^1`$, $`E_z=E(1+z)`$, and $`\chi `$ is the comoving distance along the line of sight. If $`\delta (z)`$ is the mass fluctuation at some position at redshift $`z`$, then the fluctuation of the emissivity at the same position and redshift is given by $`b_X(z)\delta (z)\overline{j}(E_z,z)`$, where we have defined the X-ray emission–weighted bias $`\overline{b}_X(z)`$ as $$\overline{b}_X(z)=\frac{1}{\overline{j}(E_z,z)}\frac{1}{4\pi }_0^{\mathrm{}}𝑑L\frac{d\varphi }{dL}L_X(E_z,L)b_X(L,z).$$ (9) For simplicity, we compute the auto–correlation $`w_\theta `$ of the XRB using the Limber approximation, together with the small angle approximation, as: $`C_{\mathrm{}}(E)`$ $`=`$ $`\overline{I}(E)^2{\displaystyle \frac{d\chi }{r_\chi ^2}W^2(E_\chi ,\chi )P_0(\mathrm{}/r_\chi )}`$ (10) $`w_\theta (E)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}d\mathrm{}}{2\pi }C_{\mathrm{}}(E)J_0(\mathrm{}\theta )},`$ where $`C_{\mathrm{}}`$ is the angular power spectrum, $`J_0`$ is the zeroth order Bessel function, $`P_0(\mathrm{}/r_\chi )`$ is the linear mass power spectrum today, $`r_\chi `$ is the angular diameter distance ($`=\chi `$ for a flat universe), and $`W(E_\chi ,\chi )=\overline{j}(E_z,z)\overline{b}_X(z)D(z)/(1+z)`$. When the power spectrum is measured in practice, shot-noise has to be subtracted or should be included in the theoretical prediction, whereas the same is not necessary for the angular correlation except at zero-lag. Our models predicts the correct mean background spectrum $`\overline{I}(E)`$, computed from equation 8, at $`E=1`$ keV. We have included all quasars down to the observed 1 keV flux of $`2\times 10^{17}\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, i.e. we used our models to extrapolate the XRLF to two orders of magnitude fainter than the ROSAT detection threshold for discrete sources (Hasinger & Zamorani 1997), to make up the remaining $`50\%`$ of the XRB at 1keV. Our models predict a faint–end slope that is steeper than the Miyaji et al. (2000) fitting formulae, allowing faint quasars to contribute half of the background. The emissivities peak at $`z2`$, coinciding with the peak of the XRLF, implying that our model produces most of the XRB, as well as its auto–correlation signal at $`z2`$. Note that the known contribution from nearby galaxy clusters is $`10\%`$ (Gilli et al. 1999), which we ignore here. In Figure 5, we show our predictions for the two point angular correlation $`w_\theta `$ of the XRB from quasars at 1 keV. Most measurements at the soft X-ray bands have yielded only upper limits, which are consistent with our predictions (e.g. variance at $`<\mathrm{\hspace{0.33em}0.12}`$ at a scale of $`10`$ arcmin. and $`E=0.92`$keV, from Carrera et al. 1998; see also references therein). Soltan et al. (1999) obtained angular correlations significantly higher than previous results (dashed curve), which, taken at face value, would imply quasar lifetimes $`t_Q10^8`$ yr. However, the results of Soltan et al. (1999) could be partially explained by galactic contamination (Barcons et al. 2000). We therefore view this measurement as an upper limit, which is consistent with models using both lifetimes we considered. Figure 5 shows that $`w_\theta `$ predicted in the long and short lifetime models differ by a factor of $`2`$ on angular scales of 0.1-1 degrees, offering another potential probe of quasar lifetimes, provided that $`w_\theta `$ can be measured more accurately in the future, and that the contribution to the clustering signal from nearby non-quasar sources (e.g. clusters) is small or can be subtracted out. Finally, we note that there have been detections of of clustering on several-degree-scales at the hard X-ray bands ($`210`$ keV) from the HEAO satellite (Treyer et al. 1998) – while a prediction for such energies would be interesting (see also Lahav et al. 1997), it would require an extrapolation of the X-ray spectrum, since we normalize by fitting to the soft-Xray luminosity function. ## 7. Further Considerations We have shown above that the quasar lifetime could be measured to high precision, using the soon available large samples of quasars at $`z<\mathrm{\hspace{0.33em}3}`$; either from the 2dF or the SDSS survey. This precision, however, reflects only the statistical errors in the simple model we have adopted for relating quasars to dark matter halos. The main hindrance in determining the quasar lifetime will likely be systematic errors; here we discuss how several potential complications could affect the derived lifetime. Obscured sources. Considerations of the hard X–ray background have led several authors to suggest the presence of a large population of “absorbed” quasars, necessary to fit the hard slope and overall amplitude of the background. Although not a unique explanation for the XRB, this would imply that the true number of quasars near the faint end of both the optical and soft X–ray LF is $`10`$ times larger than what is observed; 90% being undetected due to large absorbing columns of dust in the optical, and neutral hydrogen in the soft X–rays (see, e.g. Gilli et al. 1999). Unless the optically bright and dust–obscured phases occur within the same object (Fabian & Iwasawa 1999), this increase would have a direct effect on our results, since we would then need to adjust our fitting parameters to match a $`10`$ times higher quasar abundance. We find that this is easily achieved by leaving $`x_0`$ and $`\alpha `$ unchanged, and instead raising the lifetime from $`10^{6.5}`$ to $`10^{7.5}`$ yr in the short lifetime model, and from $`10^8`$ to $`10^{8.6}`$ yr in the long lifetime model. In the latter case, a 10-fold increase in the duty cycle requires only an increase in $`t_Q`$ by a factor of $`4`$, owing to the shape of the age distribution $`dp_a/dt`$ (Lacey & Cole 1993). Our results would then hold as before, but they would describe the two cases of $`t_Q=10^{7.5}`$ and $`10^{8.6}`$ yr. Interestingly, this scenario would imply that the quasar lifetime could not be shorter than $`t_Q10^{7.5}`$ yr, simply based on the abundance of quasars (cf. § 2). Future infrared and hard X-ray observations should help constrain the abundance of obscured sources, and reduce this systematic uncertainty. Multiple BH’s in a single halo. A possibility that could modify the simple picture adopted above is that a single halo might host several quasar black holes. A massive (e.g. $`10^{14}\mathrm{M}_{}`$) halo corresponds to a cluster of galaxies; while the Press–Schechter formalism counts this halo as a single object. If quasar activity is triggered by galaxy–galaxy mergers, a massive Press-Schechter halo, known to contain several galaxies, could equally well host several quasars (e.g. Cavaliere & Vittorini 1998). There is some observational evidence of perhaps merger driven double quasar activity (Owen et al. 1985; Comins & Owen 1991). One could therefore envision that quasars reside in the sub–halos of massive “parent” halos – a scenario that would modify the predicted clustering. To address these issues in detail, one needs to know the mass–function of sub–halos within a given parent halo, as well as the rate at which they merge and turn on. In principle, this information can be extracted from Monte Carlo realizations of the formation history of halos in the extended Press–Schechter formalism (i.e. the so called merger tree) together with some estimate of the time-scale for mergers of sub–halos based on, for instance, dynamical friction (e.g. Kauffmann & Haehnelt 2000). Here, we consider two toy models that we hope can bracket the plausible range of clustering predictions. To simplify matters, we ignore the scatter in $`L`$$`M`$ in the following discussion. Suppose one is interested in quasars of a luminosity $`L`$ at redshift $`z_0`$, which correspond to Press-Schechter halos of mass $`M_0`$ in our formalism as laid out in §2. This choice of $`M_0`$ matches the abundance of quasars, expressed approximately as $`L(d\mathrm{\Phi }/dL)M_0[dN(M_0,z_0)/dM_0](t_Q/t_{\mathrm{Hub}})`$ (the merger, or activation rate of halos is approximated as $`t_{\mathrm{Hub}}^1`$, where $`t_{\mathrm{Hub}}`$ is the Hubble time), and implies the bias $`b_0(L)b(M_0,z_0)`$. In model A, we suppose the Press-Schechter halos are identified at some earlier redshift $`z_1`$ – these would be sub–halos of those Press-Schechter halos identified at $`z_0`$. Quasars of luminosity $`L`$ now correspond to sub–halos of mass $`M_1`$. The abundance of these sub-halos is given by the Press-Schechter mass function $`dN(M_1,z_1)/dM_1`$, which is related to the mass function at $`z_0`$ by $`dN(M_1,z_1)/dM_1=_{M_1}^{\mathrm{}}𝑑M[dN(M,z_0)/dM][d𝒩(M_1,z_1\|M,z_0)/dM_1]`$ where $`dM_1\times d𝒩(M_1,z_1\|M,z_0)/dM_1`$ is the average number of $`M_1\pm dM_1/2`$ sub–halos within parent halos of mass $`M`$, given by (e.g. Sheth & Lemson 1999) $`d𝒩(M_1,z_1\|M,z_0)/dM_1=`$ (11) $`{\displaystyle \frac{M}{M_1}}{\displaystyle \frac{1}{\sqrt{2\pi }}}{\displaystyle \frac{\delta _1\delta _0}{[\sigma ^2(M_1)\sigma ^2(M)]^{3/2}}}`$ $`\mathrm{exp}\left\{{\displaystyle \frac{(\delta _1\delta _0)^2}{2[\sigma ^2(M_1)\sigma ^2(M)]}}\right\}{\displaystyle \frac{d\sigma ^2(M_1)}{dM_1}},`$ where $`\delta _1=\delta _c/D(z_1)`$ and $`\delta _0=\delta _c/D(z_0)`$ (see eq. 2). To match the abundance of quasars at luminosity $`L`$, we impose the condition that $`M_1[dN(M_1,z_1)/dM_1]`$ $`M_0[dN(M_0,z_0)/dM_0]`$ , which determines $`M_1`$ given $`z_1`$, $`M_0`$ and $`z_0`$. The bias of the quasars is no longer $`b_0(L)b(M_0,z_0)`$, but is instead given by $`b_{\mathrm{eff}}^A(L,z_1)=\left[{\displaystyle \frac{dN(M_1,z_1)}{dM_1}}\right]^1`$ (12) $`{\displaystyle _{M_1}^{\mathrm{}}}𝑑M{\displaystyle \frac{dN(M,z_0)}{dM}}{\displaystyle \frac{d𝒩(M_1,z_1\|M,z_0)}{dM_1}}b(M,z_0).`$ We show in Fig. 6 the ratio of $`b_{\mathrm{eff}}^A(L,z_1)/b_0(L)`$ as a function of $`z_1`$, for $`z_0=3`$ and $`z_0=2`$ respectively, and for a range of masses $`M_0`$ which are representative of the halos that dominate our clustering signal in previous discussions. It is interesting how the bias $`b_{\mathrm{eff}}^A(L,z_1)`$ is not necessarily larger than our original bias $`b_0(L)`$, despite the fact that the bias of sub–halos should be boosted by their taking residence in bigger halos. This is because the relevant masses here (e.g. $`M_0`$) are generally large, and we find that the number of halos of mass $`M_0`$ at $`z_1`$ is always smaller than the number of halos of the same mass at $`z_0<z_1`$. Hence, to match the observed abundance of quasars at the same $`L`$, $`M_1`$ must be chosen to be smaller than $`M_0`$. As Fig. 6 shows, this could, in some cases, more than compensate the increase in clustering due to massive parent halos. Because of these two opposing effects, the bias does not change by more than about $`50\%`$ even if one considers $`z_1`$ as high as $`10`$. This translates into a factor of $`2`$ uncertainty in our predictions for the quasar power spectrum. Our clustering predictions for the short and long lifetime models differ by about a factor of $`5`$, implying that the lifetime can still be usefully constrained at $`z>\mathrm{\hspace{0.33em}2}`$. As we can see from Fig. 6, at lower redshifts, or equivalently lower $`M_0/M_{}`$, our predictions for the quasar power spectrum should be more uncertain. One might imagine modifying the above model by allowing mergers to take place preferentially in massive parents, and therefore boosting the predicted bias. In model B, we adopt a more general procedure of matching the observed quasar abundance by $`L(d\mathrm{\Phi }/dL)=M_1_{M_1}^{\mathrm{}}𝑑M[dN(M,z_0)/dM][d𝒩(M_1,z_1\|M,z_0)/dM_1]`$ $`(t_Q/t_{\mathrm{Hub}})f(M_1,M)`$ where $`(t_Q/t_{\mathrm{Hub}})f(M_1,M)`$ is the probability that an $`M_1`$ sub–halo residing within a parent halo of $`M_0`$ harbors an active quasar of luminosity $`L`$. It is conceivable that $`f`$ increases with the parent mass $`M_0`$ – a more massive parent might encourage more quasar activity by having a higher fraction of mergers or collisions. The following heuristic argument shows that one might expect $`[d𝒩(M_1,z_1\|M,z_0)/dM_1]f(M_1,M)`$ to scale approximately as $`M^{4/3}`$. Let $`N_h`$ be the number of sub–halos inside a parent halo of mass $`M`$. The rate of collisions is given by $`N_h^2v_h\sigma _h/R^3`$ where $`v_h`$ is the velocity of the sub–halos, $`\sigma _h`$ is their cross-section and $`R^3`$ is the volume of the parent halo. Using $`N_hM`$ (which can be obtained from eq. 11 in the large $`M`$ limit), $`v_h\sqrt{M/R}`$ (virial theorem) and $`R^3M`$ (fixed overdensity of $`200`$ at the redshift of formation), the rate of collisions scales with the parent mass as $`M^{4/3}`$, if one ignores the possibility that $`\sigma _h`$ might depend on the parent mass as well. A similar scaling of $`M^{1.3}`$ has been observed in simulations of the star-burst model for Lyman-break objects (Kolatt et al. 1999, Weschler et al. 1999). To model the enhanced rate of collisions inside massive parent halos, we can simply modify model A by using $`f(M_1,M)=(M/M_1)^{1/3}`$. The effective bias is given by $`b_{\mathrm{eff}}^B(L,z_1)=`$ (13) $`\left[{\displaystyle _{M_1}^{\mathrm{}}}𝑑M{\displaystyle \frac{dN(M,z_0)}{dM}}{\displaystyle \frac{d𝒩(M_1,z_1\|M,z_0)}{dM_1}}\left({\displaystyle \frac{M}{M_1}}\right)^{\frac{1}{3}}\right]^1`$ $`\times {\displaystyle _{M_1}^{\mathrm{}}}dM{\displaystyle \frac{dN(M,z_0)}{dM}}{\displaystyle \frac{d𝒩(M_1,z_1\|M,z_0)}{dM_1}}\left({\displaystyle \frac{M}{M_1}}\right)^{\frac{1}{3}}b(M,z_0).`$ We find that the above prescription does not significantly alter our conclusions following from model A: the $`M^{1/3}`$ enhancement of the activation rate inside massive parent halos turns out to be relatively shallow, and translate to a small effect in the bias. Finally, we note that $`z_1`$ above could in principle depend on $`M_0`$ and $`M_1`$, a possibility that would require further modeling and is not pursued here. Galaxies without BH’s. Another possibility that could modify our picture is that only a fraction $`f<1`$ of the halos harbor BH’s; the duty cycle could then reflect this fraction, rather than the lifetime of quasars. Although there is evidence (e.g. Magorrian et al. 1998) that most nearby galaxies harbor a central BH, this is not necessarily the case at redshifts $`z=23`$: the fraction $`f`$ of galaxies hosting BH’s at $`z=23`$ could, in principle have merged with the fraction $`1f`$ of galaxies without BH’s, satisfying the local constraint. Using the extended Press–Schechter formalism (Lacey & Cole 1993), one can compute the rate of mergers between halos of various masses. On galaxy–mass scales, the merger rates at $`z=23`$ are comparable to the reciprocal of the Hubble time (cf. Fig. 5 in Haiman & Menou 2000), implying that a typical galaxy did not go through numerous major mergers between $`z=23`$ and $`z=0`$, i.e. that the fraction $`f`$ cannot be significantly less than unity at $`z=23`$. A more detailed consideration of this issue is beyond the scope of this paper; we simply note that the lifetimes derived here scale approximately as $`1/f`$, where $`f`$ is likely of order unity. Larger scatter in $`L/M`$. The scatter $`\sigma `$ we assumed around the mean relation between quasar luminosity and halo mass is motivated by the scatter found empirically for the $`M_{\mathrm{bh}}M_{\mathrm{bulge}}`$ relation (Magorrian et al. 1998). It is interesting to consider the sensitivity of our conclusions to an increased $`\sigma `$. In general, scatter raises the number of quasars predicted by our models, by an amount that depends on the slope of the underlying mass function $`dN/dM`$. As a result, increasing $`\sigma `$ raises, and flattens the predicted LF. We find that an increase of $`\sigma `$ from 0.5 to 1 (an additional order of magnitude of scatter) can be compensated by a steeper $`\overline{L}_M`$ relation, typically replacing $`\alpha `$ with $`\alpha 0.5`$. As a result of the increase in $`\sigma `$, quasars with a fixed $`L`$ are, on average, associated with larger, and more highly biased halos. Nevertheless, we find that the mean bias $`\overline{b}`$ of all sources above a fixed flux (cf. eq. 5), and therefore the correlation length $`r_0`$, is unchanged by the increased scatter (at the level of $`3\%`$). The reason for the insensitivity of $`r_0`$ to the amplitude of the scatter can be understood as follows. The mean bias $`\overline{b}`$ of all sources with $`L>L_{\mathrm{min}}`$ is dominated by the bias $`b(L)`$ of sources near the threshold $`L_{\mathrm{min}}`$. The latter is obtained by averaging $`b(M)`$ over halos of different masses (cf. eq. 3), and it is dominated by the bias of the smallest halos within the width of the scatter, i.e. of halos with mass $`M_{\mathrm{min}}\overline{M}/10^\sigma `$, where $`\overline{M}`$ defines the mean relation between $`L_{\mathrm{min}}`$ and halo mass i.e. $`L_{\mathrm{min}}`$ = $`\overline{L}(\overline{M})`$, and $`\sigma `$ quantifies the scatter (see eq. 16). As mentioned before, increasing the scatter makes the luminosity function flatter, which means to match the observed abundance of halos at a fixed luminosity $`L_{\mathrm{min}}`$, one has to choose a higher $`\overline{M}`$. In other words, $`\overline{M}`$ scales up with the scatter, and it turns out to scale up approximately as $`10^\sigma `$, making $`M_{\mathrm{min}}`$ and hence the effective bias roughly independent of scatter. We note that the relation between quasar luminosity and halo mass can, in principle be derived from observations, by measuring $`M_{\mathrm{halo}}`$ for the hosts of quasars (e.g. by weak lensing, or by finding test particles around quasars, such as nearby satellite galaxies). Mass and redshift dependent lifetime. In all of the above, we have assumed that the quasar lifetime is a single parameter, independent of the halo mass. This is not unreasonable if the Eddington time, the timescale for the growth of black hole mass, is indeed the relevant time–scale, $`4\times 10^7`$ ($`ϵ`$/0.1) yr. Implicit in such reasoning is that the active phase of the quasar is coincident with the phase where the black hole gains most of its mass. This is not the only possibility; see Haehnelt et al. (1999) for more discussions. One can attempt to explore how $`t_Q`$ depends on halo mass by applying our method to quasars grouped into different absolute luminosity ranges, but the intrinsic scatter in the mass-luminosity relation should be kept in mind. We emphasize, however, since we fit the luminosity function and clustering data at the same redshift, there is no need within our formalism to assume a redshift independent lifetime. In fact, performing our exercise as a function of redshift could give interesting constraints on how $`t_Q`$ evolves with redshift. ## 8. Conclusions In this paper, we have modeled the quasar luminosity function in detail in the optical and X–ray bands using the Press–Schechter formalism. The lifetime of quasars $`t_Q`$ enters into our analysis through the duty–cycle of quasars, and we find that matching the observed quasar LF to dark matter halos yields the constraint $`10^6<t_Q<\mathrm{\hspace{0.33em}10}^8`$ yr: smaller lifetimes would imply overly massive BH’s, while longer lifetimes would necessitate overly massive halos. This range reassuringly brackets the Eddington timescale of $`4\times 10^7`$ ($`ϵ`$/0.1) yr. The main conclusion of this paper is that the lifetime (and hence $`ϵ`$, if the Eddington time is the relevant timescale for quasar activity) can be further constrained within this range using the clustering of quasars: for quasars with a fixed luminosity, longer $`t_Q`$ implies larger host halo masses, and higher bias. We find that as a result, the correlation length $`r_0`$ varies strongly with the assumed lifetime. Preliminary data from the 2dF survey already sets mild constraints on the lifetime. Depending on the assumed cosmology, we find $`t_Q=10^{7.7\pm 0.8}`$ yr ($`\mathrm{\Lambda }`$CDM) or $`t_Q=10^{8\pm 0.8}`$ yr (OCDM) to within $`1\sigma `$ statistical uncertainty. These values are also found to satisfy upper limits on the auto–correlation function of the soft X–ray background. Forthcoming large quasar samples from SDSS or the complete 2dF survey are ideally suited for a study of quasar clustering, and they can, in principle constrain the quasar lifetime to high accuracy, with small statistical errors. We expect the modeling of the quasar–halo relation, as well as the possible presence of obscured quasars, to be the dominant sources of systematic uncertainty. Not discussed in depth in this paper is the possibility of using higher moments (such as skewness), which our models also make definite predictions for, and will be considered in a future publication. Remarkably, our best determination of the lifetime of quasars might come from the statistics of high–redshift quasars, rather than the study of individual objects. Near the completion of this work, we became aware of a similar, independent study by P. Martini & D. Weinberg. We thank M. Haehnelt, K. Menou, E. Quataert, U-L. Pen., U. Seljak, R. Sheth and D. Spergel for useful discussions, T. Miyaji for his advice on the X–ray luminosity function, and T. Shanks and S. Croom for discussions on the 2dF survey. Support for this work was provided by the DOE and the NASA grant NAG 5-7092 at Fermilab, by the NSF grant PHY-9513835, by the Taplin Fellowship to LH and by NASA through the Hubble Fellowship grant HF-01119.01-99A to ZH, awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5-26555. ## Appendix In this Appendix, we describe our models for the luminosity function (LF) of quasars, based on associating quasar BHs with dark halos. Our main assumption is that there is, on average, a direct monotonic relation between halo mass and quasar light. Our treatment is similar to previous works (Haiman & Loeb 1998; Haehnelt et al. 1998), but differs in some of the details. We adopt the parameterization of the observational LF in the optical B band, given in the redshift range $`0<z<\mathrm{\hspace{0.33em}4}`$ by Pei (1995). We assume the background cosmology to be either flat ($`\mathrm{\Lambda }`$CDM) with $`(\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_\mathrm{m},h,\sigma _{8\mathrm{h}^1},n)=(0.7,0.3,0.65,1.0,1.0)`$ or open (OCDM) with $`(\mathrm{\Omega }_\mathrm{\Lambda },\mathrm{\Omega }_\mathrm{m},h,\sigma _{8\mathrm{h}^1},n)=(0,0.3,0.65,0.82,1.3)`$. The LF quoted by Pei (1995) is scaled appropriately with cosmology by keeping $`(d\varphi /dL)dVd_\mathrm{L}^2`$=const, where $`dV`$ is the volume element, and $`d_\mathrm{L}`$ is the luminosity distance), so that $`(d\varphi /dL)dL`$ is the comoving abundance in $`\mathrm{Mpc}^3`$ of quasars with B–band luminosity $`L`$ (in solar units $`L_{B,}`$). The comoving abundance $`dN/dM(M,z)`$ of dark halos is assumed to follow the Press–Schechter (1974) formalism. We assume that each halo harbors a single quasar that turns on when the halo forms, i.e. triggered by merger (e.g. Percival & Miller 1999), and shines for a fixed lifetime $`t_Q`$ (relaxing these assumptions is discussed below in section 7). The duty–cycle $`f_{\mathrm{on}}`$ of halos with mass $`M`$ at redshift $`z`$, is then given by the fraction of these halos younger than $`t_Q`$. The distribution of ages $`dp_a/dt(M,z,t)`$ for halos of mass $`M`$ existing at redshift $`z`$ is obtained using the extended Press-Schechter formalism, which assumes that the halo formed at the epoch when it acquired half of its present mass (Lacey & Cole 1993). The duty–cycle, which is the probability that a dark matter halo of a given mass harbors an active quasar, is simply $$f_{\mathrm{on}}(M,z)=_0^{t_Q}𝑑t\frac{dp_a}{dt}(M,z,t).$$ (14) A model in which the quasar turns on/off more gradually (as expected if the mass of the BH grows significantly during the luminous quasar phase) is equivalent to one having additional scatter in the ratio $`L/M`$, which is discussed in § 7. We next relate the quasar luminosity to the mass of its host halo. We define $`dp/dL(M,L,z)`$ to be the probability that a halo of mass $`M`$ at redshift $`z`$ hosts a quasar with luminosity $`L`$, and express this quantity as $$\frac{dp}{dL}(L,M,z)=\frac{dg}{dL}(L,\overline{L}_{M,z})f_{\mathrm{on}}(M,z).$$ (15) Here $`dg/dL(L,\overline{L}_{M,z})`$ is the probability distribution of luminosities associated with the subset of halos of mass $`M`$ harboring a live quasar (normalized to $`_0^{\mathrm{}}𝑑L𝑑g/𝑑L=1`$), and $`\overline{L}_{M,z}`$ is the mean quasar luminosity for these halos. In the limit of a perfect intrinsic correlation, we would have $`dg/dL(L,\overline{L}_{M,z})=\delta (L\overline{L}_{M,z})`$; more realistically, this correlation will have non–negligible scatter. Lacking an a–priori theory for this scatter, we here simply assume that it follows the same functional form as the scatter found empirically for the $`M_{\mathrm{bh}}M_{\mathrm{bulge}}`$ relation (Magorrian et al. 1998), and we set $$\frac{dg}{dL}(L,\overline{L}_M)\mathrm{exp}((\mathrm{log}L/\overline{L}_{M,z})^2/2\sigma ^2).$$ (16) For reference, we note that the empirical scatter between $`M_{\mathrm{bh}}`$ and $`M_{\mathrm{bulge}}`$ gives $`\sigma 0.5`$, it is not yet clear, however, what fraction of this scatter is intrinsic vs. instrumental (van der Marel 1999). One might expect the scatter in the $`L`$$`M_{\mathrm{halo}}`$ relation not to be significantly larger, since (i) for a sufficiently high fueling rate, the luminosity $`L`$ corresponding to $`M_{\mathrm{bh}}`$ is likely to always be near the Eddington limit, and (ii) at least for disk galaxies, the bulge luminosity correlates well with the total luminosity $`L_{\mathrm{tot}}`$ ($`\sigma 0.5`$, see e.g. Andredakis & Sanders 1994); $`L_{\mathrm{tot}}`$ is tightly correlated with the velocity dispersion $`\sigma _v`$ through the Tully-Fisher relation (e.g. Raychaudhury et al. 1997) as is the total halo mass to $`\sigma _v`$ (Eisenstein & Loeb 1996). Nevertheless, in § 7 below we will investigate the consequences of an increased scatter. We note that an extension of the models presented here, by following the merger histories of halos and their BH’s can, in principle, be used to estimate the scatter in $`L/M_{\mathrm{halo}}`$. Cattaneao, Haehnelt & Rees (1998) have used this approach to fit the observed relation $`M_{\mathrm{bh}}/M_{\mathrm{bulge}}`$, including its scatter. Under the above assumptions, the cumulative probability that a halo of mass $`M`$ hosts a quasar with luminosity equal to or greater than $`L`$ is given by $$p(L,M,z)=f_{\mathrm{on}}(M,z)_L^{\mathrm{}}𝑑L\frac{dg}{dL}(L,\overline{L}_{M,z}),$$ (17) and matching the observed cumulative quasar LF requires $$_L^{\mathrm{}}𝑑L\frac{d\varphi }{dL}(L,z)=_0^{\mathrm{}}𝑑M\frac{dN}{dM}(M,z)p(L,M,z),$$ (18) or alternatively, matching the differential LF gives $$\frac{d\varphi }{dL}(L,z)=_0^{\mathrm{}}𝑑M\frac{dN}{dM}(M,z)\frac{dg}{dL}(L,\overline{L}_{M,z})f_{\mathrm{on}}(M,z).$$ (19) Equation 18 or 19, together with equations 14, 16, and 17 implicitly determines the function $`\overline{L}_{M,z}`$, once the quasar lifetime $`t_Q`$ and magnitude of scatter $`\sigma `$ are specified. In general, these equations would need to be solved iteratively. In practice, we have found that sufficiently accurate solutions (given the error bars on the observational LF in Pei 1995; cf. Fig. 1) can be found by using the simple power–law ansatz: $$\overline{L}_{M,z}=x_0(z)M_{\mathrm{halo}}\left(\frac{M_{\mathrm{halo}}}{M_0}\right)^{\alpha (z)},$$ (20) where the coefficients $`x_0(z)`$ and $`\alpha (z)`$ depend on $`t_Q`$, $`\sigma `$, and the underlying cosmology (Haiman & Loeb 1998; Haehnelt et al. 1998). In summary, assuming a fixed scatter $`\sigma `$, our model has only one free parameter, the quasar lifetime $`t_Q`$. We emphasize that our parameterization in equation 20 is purely phenomenological – it gives us a convenient way to relate the quasar luminosity to the host halo mass ($`\overline{L}_{M,z}`$). In reality, the quasar luminosity likely depends on the details of its immediate physical environment (e.g. gas supply, magnetic fields, angular momentum distribution, etc.), in addition to the halo mass. Our description includes these possibilities only in allowing a non–negligible scatter around the mean relation $`\overline{L}_{M,z}`$. The rationale behind this choice is that the average properties of the physical environment should ultimately be governed by the halo mass (or circular velocity), as expected within the picture of structure formation via hierarchical clustering. A useful check on the physical implications of equation 20 is obtained by assuming that the luminosity $`\overline{L}_{M,z}`$ is produced by a BH of mass $`M_{\mathrm{bh}}`$, shining at the Eddington limit $`L_{\mathrm{Edd}}=(4\pi G\mu m_pc/\sigma _T)M_{\mathrm{bh}}`$. In the mean spectrum of a sample of quasars with detections from radio to X–ray bands (Elvis et al. 1994), $`7\%`$ of the bolometric luminosity is emitted in the rest–frame $`B`$ band, resulting in $`L=0.07L_{\mathrm{Edd}}=5\times 10^3\mathrm{L}_{\mathrm{B},}(\mathrm{M}_{\mathrm{bh}}/\mathrm{M}_{})`$. Equation 20 then translates into a relation between the mass of a BH and its host halo, $$\overline{M}_{\mathrm{bh}}=10^{3.7}x_0(z)M_{\mathrm{halo}}\left(\frac{M_{\mathrm{halo}}}{M_0}\right)^{\alpha (z)}.$$ (21) As an example, Haehnelt et al. (1999) argue that the central BH mass is determined by a radiative feedback from the central BH that would unbind the disk in a dynamical time. Their derived scaling corresponds to $`\alpha =2/3`$ and $`x_0(1+z)^{5/2}`$, not far from what we find for the long–lifetime case (cf. Figure 7 and discussion below). In Figure 7, we show the values of the parameters $`x_0(z)`$ and $`\alpha (z)`$ obtained in our models when two different quasar lifetimes are assumed, $`t_Q=10^{6.5}`$ (solid curves) and $`t_Q=10^8`$ yr (dotted curves). The filled dots show the parameters in $`\mathrm{\Lambda }`$CDM, and the empty dots in the OCDM cosmology. We have set the arbitrary constant $`M_0=10^{12}\mathrm{M}_{}`$ in both cases. Note that $`t_Q`$ determines both $`\alpha `$ and $`x_0`$, and therefore the values of $`\alpha `$ and $`x_0`$ are correlated. In general, the fitting parameters show little evolution in the range $`2<z<4`$, around the peak of the quasar LF. According to equation 21, the corresponding BH masses in, e.g. a $`10^{12}\mathrm{M}_{}`$ halo at $`2<z<4`$ are $`M_{\mathrm{bh}}4\times 10^4M_{\mathrm{halo}}=4\times 10^8\mathrm{M}_{}`$ and $`M_{\mathrm{bh}}2\times 10^5M_{\mathrm{halo}}=2\times 10^7\mathrm{M}_{}`$ in the short and long lifetime models, respectively. The fitting procedure described above can be repeated in the X–ray bands. We therefore fit the XRLF using equation 20 analogously to the optical case, except $`\overline{L}_{M,z}`$ now denotes the X–ray luminosity at 1 keV, quoted in units of $`\mathrm{erg}\mathrm{s}^1`$. Note that the XRLF in Miyaji et al. (2000) is quoted a function of luminosity at observed 1 keV, i.e. no K–correction is applied (alternatively, the XRLF can be interpreted as the rest–frame luminosity function of sources with an average intrinsic photon index of 2). Figure 8 show the resulting fitting parameters $`x_0`$ and $`\alpha `$ in the $`\mathrm{\Lambda }`$CDM cosmology, analogous to those shown in Figure 8 for the optical case. It is apparent that both parameters have a somewhat behavior different from that in the optical. This reflects the fact that the mean quasar spectrum must evolve with redshift, or at least is black-hole/halo mass dependent: if every quasar had the same spectrum, or at least a similar X–ray/optical flux ratio, the fitting parameters derived from the optical and X–ray LF would differ only by a constant in $`x_0`$. For our purpose of deriving clustering, it is sufficient to treat $`x_0`$ and $`\alpha `$ as phenomenological fitting parameters, and we do not address the physical reason for the apparent spectral evolution (see Haiman & Menou for a brief discussion). It is important to note that the simple power-law ansatz in equation 20 with the parameters shown in Figure 8 adequately fits only the faint end of the XRLF. In the optical case, the entire range of observed luminosities is well matched by our models (cf. Fig 1). In comparison, the well–fitted range in X–rays typically extends from the detection threshold to up to 2-3 orders of magnitude in luminosity (i.e. typically upto $`3\times 10^{45}\mathrm{erg}\mathrm{s}^1`$), depending on redshift, and our models underestimate the abundance of still brighter quasars. One might then consider searching for a different ansatz to replace equation 20 that fits the entire range of the observed LF. However, we have verified that the rare quasars with these high luminosities would contribute negligibly both to the clustering signals, or the XRB investigated here. Therefore, we did not consider further improvements over equation 20, since this would not change our results.
warning/0002/cond-mat0002005.html	ar5iv	text	# Violation of ensemble equivalence in the antiferromagnetic mean-field XY model ## 1 Introduction The relation between microscopic dynamics and macroscopic thermodynamic behaviour can nowadays be explored in full detail in computer simulations Chaos . This allows one to test important hypotheses of statistical mechanics. Formost among these is equivalence of ensembles. It is widely accepted that the constant energy microcanonical ensemble gives the same results for average values as the constant temperature canonical ensemble. This is true under certain conditions Rue , among them the most important is that interactions must be short-ranged. If interactions are long-range and attractive, as for gravitating systems, all of thermodynamics breaks down due to the non extensivity of thermodynamic potentials Pad90 . For systems of Coulomb charges of opposite signs, screening effects may help in the construction of thermodynamics, but the problem is not trivially solved, even for classical systems plasmas , and each case must be examined separately. In this context, mean-field models occupy a special status. Here, thermodynamic potentials can be made extensive if the thermodynamic limit is performed by rescaling the coupling with system size and letting the range of the interaction go to infinity Hem . Mean-field models are quite well studied in the canonical ensemble and serve as a zero-order characterization of phase transitions. Although their solution is often trivial for systems without disorder, they may hide important subtleties for disordered and frustrated systems Par . Far fewer studies of mean-field models exist in the microcanonical ensemble. This may be the reason why, until now, ensemble equivalence has not been questioned for these models. In fact, one might have already expected some surprises on the basis of the exact solution of a model by Hertel and Thirring Thi70 , where in the mean-field limit in the presence of extensive thermodynamic potentials ensemble inequivalence was explicitly shown. In this paper we present a detailed study of the low temperature/energy phase of a model of classical rotators, whose potential energy is that of the mean-field antiferromagnetic XY model. The ground state of this model is highly degenerate, and while in the canonical ensemble equilibrium states are disordered at all temperatures (i.e. rotators do not display any directional organization), microcanonical ensemble simulations show the presence of a “bimodal” state where rotators are mainly grouped in two “clusters”, pointing in directions separated by an angle $`\pi `$. This dynamical effect also has thermodynamical consequences: the order parameter measuring the degree of clustering is non-zero in the limit of zero energy and the energy temperature relation is not that predicted by the canonical ensemble. We find that, although the energy temperature relation $`T=\alpha U`$ is linear, as in the canonical case, the coefficient is $`\alpha 1.3`$ and not the canonical value $`\alpha =2`$. From the point of view of dynamics, the rotators can be separated into two groups: a slow group which oscillates around the bi-cluster and a group of almost freely rotating rotators, which we call a “gas”, following an analogy with particle motion. In Section 2 we introduce the model. In Section 3 we discuss the unusual properties of its ground state and in the following Section 4 we present the main controversial points related to ensemble inequivalence. In Section 5 we present an analytical expression for the probability density function (PDF) for the orientation of a rotator, while the system is in the bi-cluster state, and we explain the consequences for the statistical properties in the system in the microcanonical ensemble. In Section 6 we discuss the main features of the dynamics of the moments of the PDF. The paper ends with some conclusions and perspectives. ## 2 The model We consider classical rotators denoted by the angle $`\theta _i`$, $`i=1,\mathrm{},N`$, which all interact with each other, with an antiferromagnetic coupling $`1/N`$ and an external field $`h`$ $$V=\frac{1}{2N}\underset{i,j}{}\mathrm{cos}(\theta _i\theta _j)h\underset{i}{}\mathrm{sin}\theta _i.$$ (1) We will mostly restrict ourselves to $`h=0`$. After defining the complex order parameters (with $`i`$ the imaginary unit) $$M_k=\frac{1}{N}\underset{n}{}\mathrm{exp}(ik\theta _n)=\|M_k\|\mathrm{exp}(i\psi _k),$$ (2) the potential can be rewritten as $$V=N\left(\frac{\|M_1\|^2}{2}h\|M_1\|\mathrm{sin}\psi _1\right)$$ (3) The microcanonical ensemble is obtained by adding a kinetic energy term to the above potential Kogut ; Antoni $$H=K+V,$$ (4) with $$K=\underset{n=1}{\overset{N}{}}\frac{p_n^2}{2}.$$ (5) In this formulation the model can also be thought of as describing a system of particles with unitary mass, interacting through the mean-field coupling $`V`$ (the names rotator and particle will be equivalently used throughout the paper). This model has been first introduced in Ref. Kogut for 2D nearest-neighbour couplings and then studied in the antiferromagnetic mean-field context in Ref. Antoni . The total energy $`E=UN`$ ($`U`$ being the energy density) is fixed by the initial conditions and is conserved in time, while temperature is defined through the time averaged kinetic energy (see Ref. temp for alternative definitions of temperature), $`T=2<K>/N`$, where $`<>=\underset{t\mathrm{}}{lim}{\displaystyle \frac{1}{t}}{\displaystyle _0^t}`$. If $`h=0`$, as is usually the case, then the total momentum $`P={\displaystyle \underset{n}{}}p_n`$ is also conserved. We set $`P=0`$ in order to avoid ballistic center of mass motion when $`h=0`$. The equations of motion $$\ddot{\theta }_n=\|M_1\|\left[\mathrm{sin}(\theta _n\psi _1)+h\mathrm{cos}\theta _n\right]$$ (6) have been integrated using an improved fourth-order symplectic scheme Mac . The algorithm is $`𝒪(N)`$, provided one first computes $`M_1`$ in the central loop. During the time evolution, we sample the instantaneous values of the order parameters (2) up to $`k=20`$ and we compute their running time averages. System size was varied from $`N=100`$ to $`N=10^4`$. We have also performed canonical Monte-Carlo simulations, using the Metropolis algorithm, for comparison. The initial conditions were of two classes: i) homogeneous state $`\theta _n=(2\pi n)/N`$, which can be shown to be marginally stable (see below), to which we add either a small spatially random perturbation $`\theta _n\theta _n+r_n`$ and/or a small momentum $`p_n=r_n`$ with zero average; ii) homogeneous state $`\theta _n=(2\pi n)/N`$ with $`p_n=A\mathrm{sin}\theta _n`$, which leads to a faster induction of the bimodal state we want to study. These initial conditions lead to the same bimodal state, discussed below, at low energy. The canonical solution of this model for $`h=0`$ is sketched in Ref. Antoni . It uses the Hubbard-Stratonovich trick to decouple the rotators and the thermodynamic limit is performed by a saddle-point technique. The result is that all moments (2) vanish in the $`N\mathrm{}`$ limit, including the magnetization $`M_1`$, which implies, on the one hand, that $$T=2U,$$ (7) on the other, that since in this limit $$M_k=_0^{2\pi }𝒫(\theta )\mathrm{exp}(ik\theta )𝑑\theta ,$$ (8) the PDF $`𝒫(\theta )`$ is flat. Hence, the bimodal state is absent in the canonical ensemble, which we have confirmed by Monte-Carlo simulations. Most of the discussions below will concentrate on the reasons for the different findings in the microcanonical and canonical statistical ensembles. ## 3 Ground State and Statistical Properties in the Canonical Ensemble The long ranged interactions mean that it is impossible to satisfy all the antiferromagnetic bonds at once and the model is highly frustrated. In zero field, the frustration is minimized for configurations with $`M_1=0`$, giving a ground state energy of $`U=E/N=0`$. The ground state is infinitely degenerate, as there is an infinity of ways of minimizing the frustration. For example, grouping the rotators into pairs with angles $`\theta _i`$ and $`\theta _i+\pi `$ ensures that $`M_1=0`$ for all configurations of the pairs. The pairs do not have to be arranged in an ordered way however, as any disordered arrangement will equally give $`M_1=0`$. Neither is the ground state manifold restricted to pairs: one can equally construct ground states from groups of three, four, five … rotators separated by angles of $`2\pi /3,2\pi /4,2\pi /5\mathrm{}`$. By moving rotators in clusters whose total angle is zero the system can evolve from one ground state to another remaining on the constant energy hypersurface. The high dimensional ground state manifold follows from the fact that the ground state condition requires the two constraints $`M_{1x}=M_{1y}=0`$ only. One can therefore expect a ground state to possess $`N2`$ unconstrained degrees of freedom, which is easily confirmed by calculating the Hessian $`J_{i,j}=^2V/\theta _i\theta _j`$. For example, for the perfectly homogeneous ground state at $`h=0`$, with $`\theta _i=(2\pi i)/N`$ $$J_{i,j}=\frac{1}{2N}\mathrm{cos}\left(\frac{2\pi (ij)}{N}\right).$$ (9) The matrix indeed has two non-zero equal eigenvalues, $`1/4`$ and $`N2`$ zero eigenvalues. Collective organization of the particles into reduced symmetry states can reduce the number of constraints even further. For example, in our system, if the rotational symmetry is broken and the spins lie along a single axis the ground state condition is reduced to a single constraint $`M_x=0`$, leading, at the harmonic level to $`N1`$ free degrees of freedom Moessner . This result is confirmed by calculating the Hessian (9), on a perfect bi-cluster groundstate with $`N/2`$ rotators at $`\theta =0`$ and $`N/2`$ at $`\theta =\pi `$. One finds only one non-zero eigenvalue, $`1/2`$, corresponding to the counter vibrating motion on the circle of the two groups of particles; all the other eigenvalues are zero. The model we study is an extreme case of a collective paramagnet Villain , or classical spin liquid Moessner ; CHS ; a system that remains disordered with no evidence of spin freezing down to the limit of zero temperature. The special points on the ground state manifold with a reduced number of constraints can dominate the partition function, in the canonical ensemble, leading to an “Order by Disorder” transition Villain2 to a reduced symmetry state. However, the mode counting arguments of Moessner and Chalker Moessner predict this to happen only if the number of liberated modes at the special points exceeds the number of zero-modes in a state with full symmetry. This is certainly not the case here and as we have confirmed by Monte Carlo simulation, Order by Disorder does not occur. We have simulated between $`N=10`$ and $`N=10^4`$ rotators down to temperatures $`T<10^4`$. At all temperatures, the system remains perfectly disordered, with no evidence of bi-cluster formation. The number of zero modes can be directly verified by measuring the specific heat at constant field, $`C_h`$, at low temperature, as each quadratic mode makes a contribution $`1/2`$ (in units of $`k_B`$), while each zero mode makes a contribution zero. We find $`C_h=1.0`$, for the $`N`$ rotator system, as expected for two regular modes. There is therefore no evidence of the system preferring states with a single quadratic mode. ## 4 Statistical Properties in the Microcanonical Ensemble The system reserves a surprise, when studied in the microcanonical ensemble, as we do not observe ensemble equivalence. For the classes of initial conditions cited above, the low energy state of the system is not one with a homogeneous distribution of angles, rather we observe the formation of a bimodal structure, with enhanced probability for two angles separated by a distance $`\pi `$ Antoni . As the energy goes to zero, the asymptotically reached state has broken symmetry, but is not perfectly bi-modal, as we discuss below. Nevertheless, when the Hessian for such a state is calculated numerically, we find that it also has only one non-zero eigenvalue. It seems therefore that the formation of the bicluster in the microcanonical ensemble, realises the condition of minimizing the number of non-zero modes. After a transient time $`\tau `$, the formation of a stable bi-cluster is revealed by a non-zero value of the second moment $`\|M_2\|`$ of $`𝒫(\theta )`$. The early time evolution of $`\|M_2\|`$ shows the typical behavior observed for the growth of the mean-field in self-consistent dynamical models Tenny with an initial exponential growth, followed by a saturation reached after a damped oscillatory motion has died away, see Figs. 1(a) and 1(b). In Figs. 1(c) and 1(d) we show the space-time plot of $`𝒫(\theta ,t)`$, which reveals the origin of the oscillations in the density waves originated periodically from the bi-cluster, preceding stabilization. The darkest colours correspond to the highest densities. A very sharply peaked but unstable bi-cluster forms over a short time period. The structure disperses with well defined instability edges that propagate from the cluster center and re-appears in a quasi-periodic manner with slowly lengthening period. The dispersing cluster appears to interact with propagating fronts from previous incarnations. The result is that the dispersion is successively slower and the fronts less well defined for following quasi-periods until finally the bi-cluster stabilizes. In addition Fig. 1(d) shows that once this coherent structures has emerged and is stabilized, it propagates around the circle with apparently ballistic dynamics. In Fig. 2(a), we show $`\|M_2\|`$, averaged over times later than $`\tau `$, as a function of temperature $`T`$, for various system sizes. For low temperatures ($`T<10^4`$), $`\|M_2\|`$ approaches the value $`\|M_2\|0.5`$ and never decays, even in extremely long simulations (times as long as $`10^8`$ in proper units<sup>1</sup><sup>1</sup>1The appropriate linear time scale of the system is $`2\pi =𝒪(1)`$). On the contrary, in the high temperature regime ($`T>10^2`$) the bi-cluster never forms, $`\|M_2\|0`$ for all times and the PDF remains flat, as in Monte-Carlo simulations. In the intermediate regime, $`10^4<T<10^2`$, the bi-cluster forms but is less well-defined, corresponding to progressively smaller values of $`\|M_2\|`$. As the temperature is increased there is a smooth transition to the homogeneous state observed in the canonical ensemble. In the next Section we will further discuss the internal structure of the bi-cluster. The choice of initial conditions discussed above is dictated by the need to progressively increase the energy starting from a ground state. One may wonder what happens for more general initial conditions. We therefore looked at statistics of the final configuration, starting from many different realizations of a uniformly random distribution of angles and zero momenta. For $`N=100`$, at temperatures around $`T10^510^4`$, the bimodal state was always reached over $`500`$ initial states. For $`N=1000`$ and $`T10^6`$, the totality of 250 random initial states went to the bimodal state. It means that the bi-cluster is fully attractive for this class of initial states, which is quite generic. As the formation of the bi-cluster is a unique property of the deterministic microcanonical system, it is an example of violation of the equivalence hypothesis between microcanonical and canonical behaviour. One might therefore expect this inequivalence to show up in other measurable quantities. This is indeed the case, for example we find that the standard relationship between energy and temperature is modified. In the canonical ensemble, as there are only two harmonic modes, equipartition of energy at low temperature leads to the energy temperature relation $`E=NT/2+T`$ ($`N`$ quadratic modes for the kinetic energy and 2 for the potential energy), which implies that $`T=2U+𝒪(N^1)`$, see (7). The eventual Order by Disorder averaging over parts of the phase space with only one quadratic modes (or any finite number of them) produces only $`1/N`$ corrections to this relation. In the microcanonical simulation, in the presence of the bi-cluster we find the anomalous relation $`T1.3U`$ as shown in Fig. 2(b). Once the temperature exceeds $`T>10^2`$, the bi-cluster does not form and the energy-temperature relation crosses over to the canonical expression. In the presence of the bi-cluster, the kinetic energy is therefore much smaller than one should expect and the mean potential energy far in excess of that predicted by equipartition for the single non-zero mode associated with the bicluster. In fact, from this result, we find that the configurational contribution to the heat capacity $`C_h<V>/T(0.35N)/1.3`$ is an extensive quantity, in complete contradiction to the predictions of mode counting in the canonical ensemble. The system manages to put a macroscopic amount of energy in the single non-zero mode corresponding to the vibrating motion of the two groups of rotators and gives the system a low-dimensional aspect with the bi-cluster taking on many characteristics of a two-particle system. Fluctuations in the magnetization $`\|M_1\|`$, whose mean value is identically zero in the ground state illustrate this point. In zero field the potential energy is $`V=N\|M_1\|^2/2`$, hence the extensive nature of the configurational heat capacity implies that mean value of the magnetization $`<\|M_1\|>`$ should be of order unity and not of order $`1/\sqrt{N}`$ as one might expect in an uncorrelated paramagnetic system. We have indeed observed ordered features in the time dependence of $`M_1`$ which are consistent with the non-fluctuational value of this quantity. ## 5 Structural and Dynamical Details of the Bi-Cluster The bi-cluster never becomes perfectly formed, with $`\|M_2\|=1`$, even in the limit $`U0`$. Rather, the angular distribution of rotators retains a width, with a certain population of rotators homogeneously distributed. In Fig. 3 we show the angular PDF $`𝒫(\theta )`$ for a system of $`N=10^4`$ particles and three different temperatures: one very low (Fig. 3(a)) and two in the transition region (Fig. 3b) and (c)). The distribution is peaked for two angles, separated by $`\pi `$, but there remains non-zero probability of finding angles between the two peaks. The degradation of the bicluster is progressive as the temperature is increased. We have made the experimental observation that the moments of the PDF assume the following values as the energy is decreased $`<\|M_k\|>`$ $``$ $`0\text{for odd k}`$ (10) $`<\|M_k\|>`$ $`=`$ $`1/\|k\|\text{for even k with}<\|M_0\|>=1,`$ (11) as shown in the insets of Fig. 3 for the first $`k=20`$ modes. This observation offers no contradiction to our finding that $`<\|M_1\|>`$ is independent of system size. As the bi-cluster is stable at low temperatures only, the numerical value of $`<\|M_1\|>`$ remains small, even though it is an intensive quantity. By summing the Fourier series for $`𝒫(\theta )`$, $$𝒫(\theta )=\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}<\|M_k\|>\mathrm{exp}(ik\theta ),$$ (12) one gets Grad $$𝒫(\theta )=\frac{1}{2\pi }\left(1\mathrm{log}(2\|\mathrm{sin}\theta \|)\right).$$ (13) This analytical formula is superimposed on to the numerical data in Fig. 3(a) with no free parameter apart from a shift of $`\theta `$, due to the motion of the bi-cluster (see below). The agreement is impressive, and although we have no theory for this result, we may well say we have a solution. We expect the analytical formula (13) to become exact in the $`U0`$ limit. The internal structure of the bi-cluster is further shown in Fig. 4(a), where the state at time $`t=2.510^5`$, with $`N=10^4`$ and $`T=1.4310^5`$, is displayed in the $`(\theta ,p)`$ single-particle phase-space (so-called Boltzmann $`\mu `$-space). A large fraction of particles get stuck in the bi-cluster and have small values of $`\|\theta \|`$, while others develop much larger $`\theta `$ values. This illustrates why we speak of a “gas-cluster” coexistence: particles appear to belong to two distinct groups; those in the bi-cluster, which perform oscillations around the two centers separated by angle $`\pi `$, and those in the gas, which have ballistic dynamics and travel over large distances. Once projected onto the $`p`$-axis the distribution is symmetric (giving, for instance, zero average momentum), but as we see it in the $`(\theta ,p)`$ space it is skewed, with, on average, bigger momenta for those particles which have traveled the furthest. This picture is further clarified if we fold the $`\theta `$-axis onto $`[0,2\pi ]`$ (Fig. 4b). The oscillating motion of the particles of the cluster is now revealed by the spiral arms visible in the bi-cluster centers (see the inset). The particles in the cluster appear now to be distributed in a sinusoidal band of lesser density. Indeed, if we look inside the region around $`\theta =0`$ (Fig. 5a), we see that the particles are quite well distributed along the phase line $$p=A\mathrm{sin}(\theta +\varphi ),$$ (14) where $`\varphi `$ is a sliding phase related to the bi-cluster motion and $`A`$ an energy related amplitude. This gives us a hint as to how to discover the dynamical mechanism at work in the process of particle evaporation from the bi-cluster. In fact, differentiating (14) with respect to time and using $`\dot{\theta }=p`$, one gets for $`u=2\theta `$ the equation of motion $$\ddot{u}=A^2\mathrm{sin}u,$$ (15) which describes a pendulum with gravity pointing upwards, $`u=0`$ being a saddle-point. This system is integrable and cannot produce the phase-distribution in Fig. 5(a). We must, therefore introduce some sort of perturbation capable of producing the sinusoidal layer observed in Fig. 4(b). A well known mechanism for producing a stochastic layer in pendulum motion is the introduction of a finite time step, as done in the Chirikov standard map. We have therefore decided to iterate the map $`u^{}`$ $`=`$ $`u+\mathrm{\Delta }tp_u`$ (16) $`p_u^{}`$ $`=`$ $`p_u+A^2\mathrm{\Delta }tu^{},`$ (17) with $`\mathrm{\Delta }t=1`$, taking an ensemble of initial points homogeneously distributed in a small square around the unstable point $`u=p_u=0`$. A snapshot of these points at iteration $`3000`$ is shown in Fig. 5(b); the particles have been stretched along the separatrix layer and distributed along it in a similar way to that shown in Fig. 5(a). Two features are responsible for the stretching: the unstable manifold distributes the particles away from the saddles, while the stable manifold attracts them to the saddles. The inhomogeneous distribution of the points is a result of these two mechanisms. One feature of the evaporation in the full model, which is not at all reproduced by the map (17), is the fact that time evolution seems to select only one of the two lobes of the separatrix; in other words, the symmetry $`pp`$ seems to be broken in the full model, while stretching in both the lobes is present in the map (17). This map is of course unable to explain why the bi-cluster forms and why it is stable. That is, if the evaporation mechanism were the only process present then the bi-cluster would be progressively depleted. This is in contrast with numerical simulations for low temperature values, which show the temporal stability of this collective state. Hence, a condensation mechanism on the bi-cluster should also exist, in order to establish a steady-state. Indeed, Fig. 4(b) shows that in the single-particle phase space two elliptic points are present at the centers of the bi-cluster, therefore in the bordering chaotic region of these two points a “trapping” effect could be present. If formula (14) is approximately verified, then the knowledge of the PDF at low temperatures (13) would allow us to compute the relationship between the temperature and the constant $`A`$. Performing the integral, it turns out that $$T=<p^2>=\frac{A^2}{4}.$$ (18) This relation is very well verified numerically for low temperature values. Moreover, particle momenta appear to follow a distribution which, although non-Gaussian, has a variance given by (18). The distribution shows a sharp peak in the center, due to the core of the bi-cluster, which is far from Gaussian. However, the wings of the distribution, which are due to the “gas” do follow a Gaussian distribution. One might propose that total momentum conservation plays a role in the observed phenomenology. An easy way to test this idea is by adding a small external magnetic field $`h`$, which removes the constraint. We observe that, for small $`h`$, the bi-cluster still forms and is stable at low temperatures. The conservation law is not therefore relevant for cluster formation. On increasing $`h`$, the distribution $`𝒫(\theta )`$ is modified, as shown in Fig. 6. The bi-cluster lies along the field axis and becomes asymmetric, with the number of rotators lying parallel and antiparallel to the field direction becoming unequal. The antiparallel cluster is continuously depleted until a single cluster eventually forms. This route towards a single cluster is somewhat counter intuitive if one thinks of a Néel ordered, unfrustrated antiferromagnet. Such a system would minimize its energy by aligning the bi-cluster perpendicular to the field and relaxing the two halves continuously out of the plane in a symmetric way. The bi-cluster would be destroyed by the two clusters, of equal size, rotating continuously onto the field direction. ## 6 Moment dynamics As we have already remarked, there is no implicit contradiction between the proposed formula for the moments of the PDF (11) and the presence of a modified energy temperature relation. The latter means that $`<\|M_1^2\|>`$ is an intensive quantity rather than being $`𝒪(1/N)`$ as in the in the canonical ensemble. Hence, although all the other odd modes of the PDF vanish in the $`N\mathrm{}`$ limit, this is not true for $`<\|M_1\|>`$. Rather, it remains finite in the low temperature regime, increasing linearly $`U`$. Solution (11) is therefore exact, only when $`U0`$. However, because of the smallness of $`U`$ in this regime, $`<\|M_1\|>`$ is a small quantity with respect to $`<\|M_2\|>`$ and equation (11) is perfectly valid. We expect that $`M_1`$ will display interesting dynamical behaviour, given its rather unexpected extensive nature, when the bi-cluster is formed. The $`(`$Re$`(M_1)`$,Im$`(M_1))`$ phase-plane is shown in Fig. 7(a) for $`N=10^4`$ and $`U=1.1310^5(T=1.4310^5)`$. Successive points are joined by lines to show the relevant properties of the motion. The phase-point has a fast oscillatory motion through zero and a much slower rotatory motion centered on zero. The fast motion is due to the vibrations of the bi-cluster around the equilibrium positions (the two components of $`M_1`$ cross zero in phase). The slow rotational motion is due to the rigid rotation of the bi-cluster. This latter motion is further revealed by the dynamics of $`M_2`$ in Fig. 7(b). Again we join successive phase-points with a line, showing that the phase point in the $`M_2`$ plane is rotating around the center, maintaining a fixed radius $`\|M_2\|0.5`$. The time dependence of the phase $`\psi _2`$ of $`M_2`$ is such that, over the long time span, the average is zero, but preliminary measurements of the variance $`\sigma ^2`$ show that the motion is ballistic, rather than diffusive, $`\sigma ^2t^2`$. This picture is not only qualitative, but also quantitative. Indeed, if we average the fast oscillatory motion of $`\|M_1\|`$ in Fig. 7(a), we get an estimate of $`<\|M_1\|^2>`$. Taking $`<\|M_1\|^2>/U0.7`$, we get the correction to the energy temperature relation $`T(20.7)U=1.3U`$ as shown in Fig. 2(b). At the transition temperature where $`\|M_2\|`$ is decreasing, both the motion of $`M_1`$ and that of $`M_2`$ become more erratic, revealing the beginning of the region where the bi-cluster is progressively depleting in time. The other odd moments of the PDF always show an erratic motion around zero, with the variance of the cloud of points decreasing as $`N`$ is increased; i.e. higher odd modes are not intensive. The higher even moments of the PDF show a pattern similar to that of Fig. 7(b), with a progressively reduced radius; the motion of the higher phases has also ballistic features. ## 7 Conclusions and perspectives The antiferromagnetic mean-field classical rotator system is an ideal laboratory to study the relation between microcanonical and canonical ensembles. Although it has a trivial canonical solution, the randomly uniform state at all temperatures, the high degeneracy of its ground state induces nontrivial dynamical effects, which are revealed in the microcanonical ensemble. Instead of maintaining a random distribution of the rotators, the Hamiltonian dynamics selects a bimodal state, where the rotators are oriented along angles at distance $`\pi `$ with some spread, a bi-cluster. We have introduced an order parameter which reveals the formation of this state in the low energy phase and we have empirically obtained an analytical formula for the probability distribution function in angle, which perfectly fits the numerical data. In addition to this first remarkable difference between the two ensembles, we have also shown that the energy temperature relation is modified and, although still linear, it has a different slope in the two ensembles. The origin of this behavior lies in the extensive amount of energy which Hamiltonian dynamics puts into the oscillatory vibrating motion of the bi-cluster. Mean-field models are characterized by self-consistency. Indeed, a feature of our model is that rotators generate themselves the mean-field in which they move. Therefore, in the large $`N`$ limit, the dynamics of our model should depend only on the interaction of the single rotator with the mean-field. As already claimed for other models (e.g. beam-plasma instability and vorticity defect model Tenny ) self-consistency effectively reduces the number of degrees of freedom. This is why many properties of the dynamics of such a complex $`N`$-body system as ours resemble those of a “simple” perturbed pendulum. This is also the origin of the ordering of the rotators in a bi-cluster. There should be entropic reasons why the rotators prefer this state rather than choosing the disordered state, which is instead selected in the canonical ensemble. Many questions remain to be explored, but the most compelling one concerns the careful description and explanation of the dynamics. The perturbed pendulum analogy must be further investigated and a thorough study of the time evolution of the moments of the probability distribution function in angle should allow a better understanding of the low-dimensional properties of the dynamics. ###### Acknowledgements. We thank M. Droz, J. Farago, M-C Firpo, M. Paliy and Z. Racz for useful discussions. S.R. thanks ENS-Lyon, INFN and INFM for financial support. P.H. thanks INFN for financial support. This work was performed using the computing resources of the Pôle Scientifique de Modélisation Numérique (PSMN) of ENS Lyon and of the DOCS-INFM group in Florence.
warning/0002/cond-mat0002421.html	ar5iv	text	# Spin stripe dimerization in the 𝑡-𝐽 model. ## I Introduction It is widely believed that the 2D $`tJ`$ model is relevant to the low energy physics of high-temperature superconductors. This is why investigation of this model is of great interest both for theory and experiment. In spite of great efforts during more than a decade there is no full understanding of the phase diagram of the $`tJ`$ model, however some facts are well established. At zero doping the model is equivalent to the Heisenberg model on a square lattice which has long range Neel order . Doping by holes destroys the order. A simplified picture of noninteracting holes leads to the Neel state instability with respect to spirals at arbitrary small but finite doping . However more sophisticated numerical calculations which take into account renormalization of the hole Green’s function under the doping indicate that the Neel order is stable below some critical hole concentration $`x_{c1}`$ . In the Neel phase ($`x<x_{c1}`$), in all waves except s-wave, there is magnon mediated superconducting pairing between holes. . It is also clear that at very small hopping there is phase separation in the model because separation leads to reduction of the number of destroyed antiferromagnetic links The purpose of the present work is to elucidate spin structure of the ground state at $`x>x_{c1}`$. The most important hint comes from experiment: indications of stripes in the high-$`T_c`$ materials . Another important hint is a remarkable stability of the spin dimerized phase in the frustrated $`J_1J_2`$ model. The idea of such state for this model was first formulated by Read and Sachdev , and was then confirmed by further work . The stability of such a configuration implies that the lattice symmetry is spontaneously broken and the ground state is four-fold degenerate. Such a route towards quantum disorder is known rigorously to take place in one dimension, where the Lieb-Schultz-Mattis (LSM) theorem guarantees that a gapped phase always breaks the translational symmetry . Some time ago Affleck suggested that the LSM theorem can be extended to higher dimensions, and the gapped states of quantum systems necessarily break the discrete symmetries of the lattice . The example of the $`J_1J_2`$ model provides further support for this idea. There have been several attempts to consider the spin-dimerized phase in a doped Heisenberg antiferromagnet. For this purpose Affleck and Marston analyzed Hubbard-Heisenberg model in the weak-coupling regime, Grilli, Castellani and G. Kotliar considered $`SU(N)`$, $`N\mathrm{}`$, $`tJ`$ model, and very recently Vojta and Sachdev considered $`Sp(2N)`$, $`N\mathrm{}`$, $`tJ`$ model. These works indicated a stability of the spin-dimerized phase in some region of parameters, providing a very important guiding line. However relevance of these results to ”physical regime” of the $`tJ`$ model remained unclear. Stability of the spin-dimer order for the $`tJ`$ model has been demonstrated in the paper . The only small parameter used in the analysis was hole concentration with respect to the half filling. In the present work we continue studies in the same direction applying various techniques. To be confident in the results we prove stability of the dimer phase by two independent methods: 1) Calculation of the magnon Green’s function in the dimerized phase (Green Function Method), 2) Comparison of the ground state energies of the doped Neel state and the dimerized state (Direct Energy Method). The second approach is very simple physically and technically. The first approach is more technically involved, but it allows us to calculate the critical concentration more precisely. To incorporate some experimental data we consider $`tt^{}t^{\prime \prime }J`$ model defined by the following Hamiltonian $$H=t\underset{ij\sigma }{}c_{i\sigma }^{}c_{j\sigma }t^{}\underset{ij_1\sigma }{}c_{i\sigma }^{}c_{j_1\sigma }t^{\prime \prime }\underset{ij_2\sigma }{}c_{i\sigma }^{}c_{j_2\sigma }+\underset{ij\sigma }{}J_{ij}\left(𝐒_i𝐒_j\frac{1}{4}n_in_j\right).$$ (1) $`c_{i\sigma }^{}`$ is the creation operator of an electron with spin $`\sigma `$ $`(\sigma =,)`$ at site $`i`$ of the two-dimensional square lattice. The $`ij`$ represents nearest neighbor sites, $`ij_1`$ \- next nearest neighbor (diagonal), and $`ij_2`$ represents next next nearest sites. The spin operator is $`𝐒_i=\frac{1}{2}c_{i\alpha }^{}\sigma _{\alpha \beta }c_{i\beta }`$ and the number density operator is $`n_i=_\sigma c_{i\sigma }^{}c_{i\sigma }`$. The $`c_{i\sigma }^{}`$ operators act in the Hilbert space with no double electron occupancy. Antiferromagnetic interactions $`J_{ij}>0`$ are arranged in a stripe pattern shown in Fig. 1: solid links correspond to $`J_{ij}=J_{}=J(1+\delta )`$, and dashed links correspond to $`J_{ij}=j=J(1\delta )`$. For real cuprates the antiferromagnetic interaction is isotropic, $`\delta =0`$. However for theoretical analysis of the magnon Green’s function it is convenient to consider nonzero $`\delta `$ and later set $`\delta 0`$. The antiferromagnetic exchange measured in two magnon Raman scattering is $`J=125meV`$. Calculation of the hopping matrix elements has been done by Andersen et al . They consider a two-plane situation, and the effective matrix elements are slightly different for symmetric and antisymmetric combinations of orbitals between planes. After averaging over these combinations we get: $`t=386meV`$, $`t^{}=105meV`$, $`t^{\prime \prime }=86meV`$. Below we set $`J=1`$, in these units $$t=3.1,t^{}=0.8,t^{\prime \prime }=0.7$$ (2) These values are confirmed by the analysis of photoemission (PES) data for charge transfer insulator Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>. In further numerical estimates we will use the values (2), but we will also consider “pure” $`tJ`$ model which corresponds to $`t^{}=t^{\prime \prime }=0`$. The rest of the paper is organized as follows. In Sec. II we describe magnon Green’s function at zero doping, and remind ideas of the Brueckner technique used in the work. In Sec III the single-hole dispersion and wave function are considered. This section is important for both Green’s Function Method and for Direct Energy Method. Sections IV and V contain the main results of the paper: In Sec. IV we demonstrate stabilization of the dimerization by doping using Green’s Function Method, and in Sec. V we come to the same conclusion using Direct Energy Method. Sec. VI addresses the quantum phase transition from the dimerized liquid to the Normal Fermi liquid at high doping. In Sec VII we discuss shape of the Fermi surface and distribution of the photoemission intensity. Sec. VIII summarizes the work. Some technical details concerning so called “triple” diagrams are discussed in Appendix. ## II Zero Doping At half filling ($`n_i=1`$) the Hamiltonian (1) is equivalent to a Heisenberg model which has already been studied: for $`\delta >\delta _c0.303`$ the ground state is a quantum state with gapped spectrum, and for $`\delta <\delta _c`$ there is spontaneous Neel ordering with gapless spin waves. In order to study the stability of the dimer phase we first derive an effective Hamiltonian in terms of bosonic operators creating spin-wave triplets (magnons) $`t_{i\alpha }^{}`$, $`\alpha =x,y,z`$ and fermionic operators creating holes $`b_\sigma ^{}`$, $`a_\sigma ^{}`$, $`\sigma =,`$ from the spin singlets shown in Fig. 1. This Hamiltonian consists of four parts: the spin-wave part $`H_t`$, the hole part $`H_h`$, the spin-wave-hole interaction $`H_{th}`$, and the hole-hole interaction $`H_{hh}`$. Let us start from $`H_t`$. Similar effective theories have been derived in Refs. and we only present the result: $`H_t`$ $`=`$ $`H_2+H_3+H_4+H_U,`$ (3) $`H_2`$ $`=`$ $`{\displaystyle \underset{𝐤,\alpha }{}}\left\{A_𝐤t_{𝐤\alpha }^{}t_{𝐤\alpha }+{\displaystyle \frac{B_𝐤}{2}}\left(t_{𝐤\alpha }^{}t_{𝐤\alpha }^{}+\text{h.c.}\right)\right\}`$ (4) $`,H_3`$ $`=`$ $`{\displaystyle \underset{1+2=3}{}}\text{R}(𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐})ϵ_{\alpha \beta \gamma }t_{𝐤_\mathrm{𝟏}\alpha }^{}t_{𝐤_\mathrm{𝟐}\beta }^{}t_{𝐤_\mathrm{𝟑}\gamma }+\text{h.c.}`$ (5) $`H_4`$ $`=`$ $`{\displaystyle \underset{1+2=3+4}{}}\text{T}(𝐤_\mathrm{𝟏}𝐤_\mathrm{𝟑})(\delta _{\alpha \delta }\delta _{\beta \gamma }\delta _{\alpha \beta }\delta _{\gamma \delta })t_{𝐤_\mathrm{𝟏}\alpha }^{}t_{𝐤_\mathrm{𝟐}\beta }^{}t_{𝐤_\mathrm{𝟑}\gamma }t_{𝐤_\mathrm{𝟒}\delta }.`$ (6) We also introduce an infinite repulsion on each site, in order to enforce the kinematic constraint $`t_{i\alpha }^{}t_{i\beta }^{}=0`$. $$H_U=U\underset{i,\alpha \beta }{}t_{i\alpha }^{}t_{i\beta }^{}t_{i\beta }t_{i\alpha },U\mathrm{}$$ (7) The following definitions are used in (3): $`A_𝐤=J_{}+B_𝐤,`$ (8) $`B_𝐤=j(\mathrm{cos}k_y0.5\mathrm{cos}k_x),`$ (9) $`\text{T}(𝐤)=j(0.25\mathrm{cos}k_x+0.5\mathrm{cos}k_y),`$ (10) $`\text{R}(𝐩,𝐪)=0.25j(\mathrm{sin}q_x\mathrm{sin}p_x).`$ (11) Throughout the paper we work in the Brillouin zone of the dimerized lattice. At zero doping ($`n_i=1`$) $`H_t`$ is an exact mapping of the original Hamiltonian (1). To analyze this case it is enough to apply the Brueckner technique . The result for the normal spin-wave Green’s function reads: $$G_N(𝐤,\omega )=\frac{\omega +\stackrel{~}{A}_𝐤(\omega )}{\{\omega +\stackrel{~}{A}_𝐤(\omega )\}\{\omega \stackrel{~}{A}_𝐤(\omega )\}+\stackrel{~}{B}_𝐤^2}$$ (12) where $`\stackrel{~}{A}_𝐤(\omega )=A_𝐤+\mathrm{\Sigma }_4^N(𝐤)+\mathrm{\Sigma }_{Br}^{(1)}(𝐤,\omega ),`$ (13) $`\stackrel{~}{B}_𝐤(\omega )=B_𝐤+\mathrm{\Sigma }_4^A(𝐤).`$ (14) Normal $`\mathrm{\Sigma }_4^N`$ and anomalous $`\mathrm{\Sigma }_4^A`$ self-energies are caused by the quartic interaction $`H_4`$ and the most important contribution $`\mathrm{\Sigma }_{Br}^{(1)}`$ comes from the Brueckner diagrams as described in . Strictly speaking there is also some contribution to the self-energy caused by the ”triple” interaction $`H_3`$. However this contribution is very small (see, e.g. Ref.) and therefore we neglect it. Expansion of the self-energy in powers of $`\omega `$ near $`\omega =0`$ gives quasiparticle residue and spin-wave spectrum $`Z_𝐤=\left(1{\displaystyle \frac{\mathrm{\Sigma }_{Br}^{(1)}}{\omega }}\right)^1,`$ (15) $`\omega _𝐤=Z_𝐤\sqrt{[\stackrel{~}{A_𝐤}(0)]^2\stackrel{~}{B_𝐤}^2}.`$ (16) Expressions for effective Bogoliubov parameters $`u_𝐤`$ and $`v_𝐤`$ are given in . The spin-wave gap $`\mathrm{\Delta }=\omega _{𝐤_\mathrm{𝟎}}`$, $`𝐤_\mathrm{𝟎}=(0,\pi )`$, obtained as a result of a selfconsistent solution of Dyson’s equations is plotted in Fig.2 (line at $`x=0`$). The critical value of the explicit dimerization (point where the gap vanishes) $`\delta _c=0.298`$ is in agreement with results of series expansions and quantum Monte Carlo simulations . The validity of the Brueckner approximation is justified by the smallness of the gas parameter $`n_t=_\alpha t_{i\alpha }^{}t_{i\alpha }`$. At the critical point $`n_t=0.13`$. ## III Single Hole Dispersion Consider now doping by holes. On the single dimer $`\|s`$ the hole can exist in symmetric (bonding) and antisymmetric (antibonding) states. Corresponding fermionic operators $`b_{i\sigma }^{}`$, and $`a_{i\sigma }^{}`$ $`\sigma =,`$ creating hole from the singlet $`\|s_i`$ are defined as: $`b_\sigma ^{}\|s`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(c_{2,\sigma }^{}+c_{1,\sigma }^{})\|0,`$ (17) $`a_\sigma ^{}\|s`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(c_{2,\sigma }^{}c_{1,\sigma }^{})\|0,`$ (18) where $`1`$ and $`2`$ numerate the dimer sites. We will see later that doping suppresses spin quantum fluctuations of the dimerized state so that $`n_t=_\alpha t_{i\alpha }^{}t_{i\alpha }`$ at $`x>x_c`$ does not exceed 0.07. This is why we can neglect these fluctuations and consider pure dimerized spin liquid. For comparison we can say that even for spin ladder influence of the quantum fluctuations on the hole dispersion is not strong in spite of the fact that in this case $`n_t0.3`$. In leading approximation the wave function of a hole with given quasimomentum is of the form $$\|1=\frac{1}{\sqrt{N_2}}\underset{n}{}e^{i\mathrm{𝐤𝐫}_𝐧}b_n^{}\|S,$$ (19) where $`\|S`$ is the spin dimerized state, index $`n`$ numerates the dimers and $`N_2=N/2`$ is number of sites in the dimerized lattice. We remind that throughout the paper we work in the Brillouin zone of the dimerized lattice. Sometimes the results are transfered to the usual lattice, but then it is specially pointed. The dispersion corresponding to the state (19) can be easily calculated considering all possible hoppings between the dimers. The result is $$ϵ_1(𝐤)=1\|H\|1=t+(t+t^{})\mathrm{cos}k_y+(\frac{1}{2}t+t^{\prime \prime })\mathrm{cos}k_x+t^{}\mathrm{cos}k_x\mathrm{cos}k_y+t^{\prime \prime }\mathrm{cos}2k_y.$$ (20) There is also an additional t-independent constant in the dispersion $$e_0=1.75J.$$ (21) The constant arises because a hole destroys one spin dimer (this gives $`0.75J`$) and four links $`1/4n_in_j`$ (see Hamiltonian (1)). In the present section we ignore $`e_0`$ because it gives just a shift of the entire dispersion. However in Section V, calculating total energy of the system we will restore $`e_0`$. The wave function (19) as well as the dispersion (20) are renormalized due to virtual admixture of antibonding states and magnon excitations. Now we calculate this admixture to show that it is small. In this calculation we follow the approach developed earlier for the doped spin ladder . The Hamiltonian (1) admixes following states to the wave function (19) $`\|2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}a_n^{}\|S,`$ (22) $`\|3`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{n+y}^{})\|S,`$ (23) $`\|4`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{ny}^{})\|S,`$ (24) $`\|5`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{n+x}^{})\|S,`$ (25) $`\|6`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{nx}^{})\|S,`$ (26) $`\|7`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{n+x}^{})\|S,`$ (27) $`\|8`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{nx}^{})\|S,`$ (28) $`\|9`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{n+x+y}^{})\|S,`$ (29) $`\|10`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{n+xy}^{})\|S,`$ (30) $`\|11`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{nx+y}^{})\|S,`$ (31) $`\|12`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{nxy}^{})\|S,`$ (32) $`\|13`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{n+x+y}^{})\|S,`$ (33) $`\|14`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{n+xy}^{})\|S,`$ (34) $`\|15`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{nx+y}^{})\|S,`$ (35) $`\|16`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}a_{nxy}^{})\|S,`$ (36) $`\|17`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{n+y}^{})\|S,`$ (37) $`\|18`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{N_2}}}{\displaystyle \underset{n}{}}e^{i\mathrm{𝐤𝐫}_𝐧}(t_n^{}b_{ny}^{})\|S.`$ (38) The state $`\|2`$ is similar to $`\|1`$ with replacement of bonding orbital to the antibonding one. The states $`\|3\|18`$ describe excited triplet (magnon) on the dimer $`n`$ and a hole on the one of the closest dimers, for example $`n+x`$ denotes the dimer on the right, $`n+x+y`$ denotes the up-right dimer, etc. The brackets $`(t^{}a^{})`$ denote that spins of the magnon and the hole are combined to the total spin $`1/2`$ and z-projection $`1/2`$: $`\|1/2,1/2`$. Calculation of matrix elements of the Hamiltonian is straightforward. The diagonal matrix elements are $`2\|H\|2`$ $`=`$ $`t+(tt^{})\mathrm{cos}k_y({\displaystyle \frac{1}{2}}tt^{\prime \prime })\mathrm{cos}k_xt^{}\mathrm{cos}k_x\mathrm{cos}k_y+t^{\prime \prime }\mathrm{cos}2k_y,`$ (39) $`3\|H\|3`$ $`=`$ $`4\|H\|4=t+J_{}j/2,`$ (40) $`5\|H\|5`$ $`=`$ $`6\|H\|6=t+J_{}j/4,`$ (41) $`7\|H\|7`$ $`=`$ $`8\|H\|8=t+J_{}j/4,`$ (42) $`9\|H\|9`$ $`=`$ $`10\|H\|10=11\|H\|11=12\|H\|12=t+J_{},`$ (43) $`13\|H\|13`$ $`=`$ $`14\|H\|14=15\|H\|15=16\|H\|16=t+J_{},`$ (44) $`17\|H\|17`$ $`=`$ $`18\|H\|18=t+J_{}j/2.`$ (45) The nonzero nondiagonal matrix elements are $`2\|H\|1`$ $`=`$ $`i{\displaystyle \frac{t}{2}}\mathrm{sin}k_x,`$ (46) $`3\|H\|1`$ $`=`$ $`4\|H\|1^{}={\displaystyle \frac{\sqrt{3}}{2}}(tt^{})+{\displaystyle \frac{\sqrt{3}}{4}}je^{ik_y},`$ (47) $`5\|H\|1`$ $`=`$ $`6\|H\|1^{}={\displaystyle \frac{\sqrt{3}}{4}}(t2t^{\prime \prime }){\displaystyle \frac{\sqrt{3}}{8}}je^{ik_x},`$ (48) $`7\|H\|1`$ $`=`$ $`8\|H\|1^{}={\displaystyle \frac{\sqrt{3}}{4}}t+{\displaystyle \frac{\sqrt{3}}{8}}je^{ik_x},`$ (49) $`9\|H\|1`$ $`=`$ $`10\|H\|1=11\|H\|1=12\|H\|1={\displaystyle \frac{\sqrt{3}}{4}}t^{},`$ (50) $`13\|H\|1`$ $`=`$ $`14\|H\|1=15\|H\|1=16\|H\|1={\displaystyle \frac{\sqrt{3}}{4}}t^{},`$ (51) $`5\|H\|2`$ $`=`$ $`6\|H\|2^{}={\displaystyle \frac{\sqrt{3}}{4}}t+{\displaystyle \frac{\sqrt{3}}{8}}je^{ik_x},`$ (52) $`7\|H\|2`$ $`=`$ $`8\|H\|2^{}={\displaystyle \frac{\sqrt{3}}{4}}(t+t^{\prime \prime }){\displaystyle \frac{\sqrt{3}}{8}}je^{ik_x},`$ (53) $`9\|H\|2`$ $`=`$ $`10\|H\|2=11\|H\|2=12\|H\|2={\displaystyle \frac{\sqrt{3}}{4}}t^{},`$ (54) $`13\|H\|2`$ $`=`$ $`14\|H\|2=15\|H\|2=16\|H\|2={\displaystyle \frac{\sqrt{3}}{4}}t^{},`$ (55) $`17\|H\|2`$ $`=`$ $`18\|H\|2^{}={\displaystyle \frac{\sqrt{3}}{2}}(t+t^{})+{\displaystyle \frac{\sqrt{3}}{4}}je^{ik_y},`$ (56) $`4\|H\|3`$ $`=`$ $`{\displaystyle \frac{(tt^{})}{2}}e^{ik_y},`$ (57) $`18\|H\|17`$ $`=`$ $`{\displaystyle \frac{(t+t^{})}{2}}e^{ik_y}.`$ (58) Diagonalization of the Hamiltonian matrix can be performed numerically. This gives the quasiparticle dispersion $`ϵ_𝐤`$ and the quasiparticle residue $`Z_𝐤^{(h)}`$ which by definition is equal to the weight of the state $`\|1`$ (see eq. (19) in the exact wave function . The minimum energy for different sets of parameters is shown in the fifth column of the Table I. We remind that we use units $`J=1`$. In the last column we show position of the minimum. For comparison we also display the minimum value of $`ϵ_1`$ (see eq. (20)), which always is at $`𝐩_0=(\pi ,\pi )`$. Table I. $`t`$ $`t^{}`$ $`t^{\prime \prime }`$ $`ϵ_{1,min}`$ $`ϵ_{min}`$ $`𝐩_0`$ 3.1 -0.8 0.7 -7.75 -9.70 $`(\pi ,0.82\pi )`$ 3. 0. 0. -7.50 -9.03 $`(\pi ,\pi )`$ 2. 0. 0. -5.0 -5.95 $`(\pi ,\pi )`$ 1. 0. 0. -2.5 -2.91 $`(\pi ,\pi )`$ Let us denote the hole concentration by $`x=n/N`$, where $`n`$ is number of holes, and $`N`$ is number of sites. Hence on-site electron occupation number is $`n_i=1x`$. Concentration of holes in terms of the dimerized lattice is two times larger $`n/(0.5N)=2x`$. At this stage we neglect interaction of holes between themselves, hence we consider them as an ideal Fermi gas, and the Fermi surface can be easily found from the condition $$2\frac{dk_xdk_y}{(2\pi )^2}=2x,$$ (59) where integration is performed inside the Brillouin zone of the dimerized lattice. The Fermi surface at $`x=0.1`$ and hopping parameters given in (2) is shown in Fig. 3 by solid line. We stress that in this figure we put $`k_x/2\pi `$ along the horizontal axis and $`k_y/\pi `$ along the vertical axis. It means that the picture corresponds to a quadrant of the Brillouin zone of the original lattice. The Fermi surface corresponding to the bare dispersion (20) is shown by the dashed line. The two curves are very close and this proves that the dispersion renormalization is small. The wave function renormalization is also not large and the quasiparticle residue is close to unity: it is $`Z^{(h)}=0.83`$ at the bottom of the band, and at the Fermi surface $`Z^{(h)}=0.80`$. So admixture of the states (22) to the bare state $`\|1`$ is relatively small. In this admixture the states $`\|3`$, $`\|4`$, $`\|7`$, and $`\|8`$ clearly dominate. In a reasonable approximation the wave function can be written as $$\psi 0.9\|10.22(\|3+\|4\|7+\|8).$$ (60) Admixture of other components is even smaller. For comparison in Fig. 3 we show also by dot-dashed line the Fermi surface for “pure” $`tJ`$ model ($`t/J=3`$, $`t^{}=t^{\prime \prime }=0`$, doping is the same, $`x=0.1`$). We see that the additional hoppings influence substantially shape of the Fermi surface. ## IV The Spin-Wave-Hole Interaction. Stabilization of the Dimer Order The magnon-hole interaction $`H_{th}`$ can be easily calculated in the way similar to that for doped spin-ladder . This interaction consists of two parts. The first one is interaction of a hole and a magnon positioned at different dimers. This is a relatively weak interaction which can be neglected . The second part, which gives the main effect, comes from the constraint that a hole and a magnon can not coexist at the same dimer: $`t_{i\alpha }^{}b_{i\sigma }^{}=t_{i\alpha }^{}a_{i\sigma }^{}=0`$. To deal with this constraint we introduce, similarly to (7), an infinite repulsion $$H_{U1}=U\underset{i,\alpha \sigma }{}t_{i\alpha }^{}t_{i\alpha }(b_{i\sigma }^{}b_{i\sigma }+a_{i\sigma }^{}a_{i\sigma }),U\mathrm{}.$$ (61) The exact hole-magnon scattering amplitude caused by this interaction can be found via Bethe-Salpeter equation shown in Fig.4a. This scattering amplitude is similar to that for magnon-magnon scattering . Solution of the Bethe-Salpeter equation gives $$\mathrm{\Gamma }(E,𝐤)=\left(\underset{𝐪}{}\frac{Z_𝐪u_𝐪^2}{E\omega _𝐪ϵ_1(𝐤𝐪)}\right)^1,$$ (62) where $`E`$ and $`𝐤`$ is total energy and total momentum of the incoming particles. It has been demonstrated in the previous section that the renormalized hole dispersion $`ϵ_𝐤`$ is close to the bare one $`ϵ_1(𝐤)`$. This is why we use in eq. (62) the bare dispersion (20). We remind that concentration of holes in terms of the dimerized lattice is $`2x1`$, and this is the gas parameter of the magnon-hole Brueckner approximation. In the previous section we have shown that the holes are concentrated in the pocket in the vicinity of $`𝐩_\mathrm{𝟎}=(\pi ,\pi )`$. Therefore the magnon normal self-energy described by the diagram Fig. 4b is $$\mathrm{\Sigma }_{Br}^{(2)}(𝐤,\omega )=2x\mathrm{\Gamma }[\omega +ϵ_1(𝐩_\mathrm{𝟎}),𝐤+𝐩_\mathrm{𝟎}]$$ (63) It is instructive to consider first the limit which allows an analytical solution: $`jJ_{}`$, $`\sqrt{2}\pi x1`$. Bare magnon dispersion in this case is $`\omega _𝐤J_{}+j(\mathrm{cos}k_y0.5\mathrm{cos}k_x)`$ and hence the integrals in (62,63) can be calculated analytically with logarithmic accuracy. This gives $$\mathrm{\Sigma }_{Br}^{(2)}(𝐤,\omega )\frac{2\sqrt{2}\pi x(t+j)}{\mathrm{ln}(12.5/\mu )+i\pi \theta (\delta \omega )},$$ (64) where $`\delta \omega ={\displaystyle \frac{1}{t+j}}\left[\omega \omega _𝐤+{\displaystyle \frac{j}{t+j}}(\omega _𝐤\omega _{𝐤_\mathrm{𝟎}})\right],`$ (65) $`\mu =\mathrm{max}(\|\delta \omega \|,\sqrt{2}\pi x),`$ (66) and $`\theta (\delta \omega )`$ is a step function. The magnon Green’s function is $$G(𝐤,\omega )=\frac{1}{\omega \omega _𝐤\mathrm{\Sigma }_{Br}^{(2)}(𝐤,\omega )}.$$ (67) For illustration the spectral function $`ImG(\omega )`$ at $`𝐤=𝐤_\mathrm{𝟎}=(0,\pi )`$, $`t/j=3`$ and different $`x`$ is plotted in Fig.5. There are several conclusions from formula (64) and Fig. 5: 1) doping pushes the spin-wave spectrum up, 2) the effect is increasing with hopping $`t`$, 3) finite width appears, 4) there is only a logarithmic dependence on the infrared cutoff. Let us stress the importance of the point (4). It means that the effect is practically independent of the long-range dynamics. Moreover, near the critical point ($`\mathrm{\Delta }=0`$) the situation is even better: the spin-wave spectrum is linear and even the logarithmic divergence disappears. Thus in the 2D case there is separation of scales which justifies Brueckner approximation. If we tried to apply the described approach to the 1D case (say a doped spin ladder) we would get into trouble: power infrared divergence appears in Brueckner diagram and hence there is no justification for gas approximation. Let us also comment on the point (3) (width). There is also a ”triple” contribution to the magnon self-energy shown in Fig. 6a,b. This is a long-range contribution, and it can be shown that it does not influence position of the critical point ($`\mathrm{\Delta }=0`$) in linear in x approximation. This fact will be proved in Appendix and this is why here we neglect the “triple” diagrams. However note that these diagrams influence the width of the magnon spectral function. In the general case there are two contributions to the Brueckner self-energy: $`\mathrm{\Sigma }_{Br}^{(1)}`$, which is due to the magnon-magnon constraint, and $`\mathrm{\Sigma }_{Br}^{(2)}`$ which is due to the magnon-hole constraint. To find the spin wave spectrum one has to solve selfconsistently Dyson’s equation for Green’s function (12), as it is described in Ref. . Results for the spin-wave ”gap” $`\mathrm{\Delta }`$ as a function of explicit dimerization $`\delta `$ for different hole concentrations $`x`$ and $`t/J=3`$ are plotted in Fig. 2. These curves are practically independent of the longer range hoppings $`t^{}`$ and $`t^{\prime \prime }`$. Strictly speaking at $`x0`$ the $`\mathrm{\Delta }`$ is not a gap because of the large decay width. What we plot is the position of the center of gravity of the magnon spectral function. However at $`\mathrm{\Delta }0`$ the width vanishes, and therefore the critical regime is uniquely defined. It is clear from Fig.2 that at $`t/J=3`$ and $`x>x_{c1}0.090`$ the ”gap” remains finite even at $`\delta =0`$. This is regime of spontaneous dimerization. Critical concentrations for other values of $`t/J`$ are presented in Table II. Table II. $`t/J`$ 1 2 3 $`x_{c1}`$ 0.132 0.106 0.090 Thus the doping stabilizes the dimerized phase. The larger the hopping $`t`$, the stronger the effect of stabilization. (The same follows from eq. (64).) This statement is true only if $`t/J10`$. At $`t/J10`$ there is a crossover to quasiparticles with higher spin (hole-magnon bound states) which indicates transition to the Nagaoka regime. The small parameter of the Brueckner approximation is concentration of holes in the dimerized lattice: $`2x`$. Therefore at $`t/J=3`$ one should expect $`20\%`$ accuracy in calculation of $`x_{c1}`$. Note that the value of $`x_{c1}`$ is close to that found in from the Neel state. Another important gas parameter is density of spin fluctuations $`n_t`$. It also proved to be small: At the critical point, $`\delta =0`$, $`x=x_c=0.09`$, the density is $`n_t0.07`$. ## V Direct Comparison of Ground State Energies of the Neel State and the Spin Dimerized State In the previous section we have demonstrated quantum phase transition from the Neel state to the spin dimerized state at some critical hole concentration $`x_{c1}`$. We calculated the magnon Green’s function in the dimerized phase, and the transition point was identified as the point where the magnon gap vanishes. This is a rigorous approach which allows to determine $`x_{c1}`$ with relatively high precision ($`20\%`$), however it is rather involved technically. An alternative method is direct comparison of the ground state energies of these two states. This is a very simple method which does not require introduction of the explicit dimerization. So throughout this section $`\delta =0`$. The energy per site for the undoped Neel state is $`1.17=0.670.5`$, where the first term is Heisenberg energy (see e. g. Ref. ), and the second contribution comes from the $`1/4n_in_j`$ term in the Hamiltonian (1). We remind that in our units J=1. In this section we consider only “pure” $`tJ`$ model, i. e. $`t^{}=t^{\prime \prime }=0`$. Energy of a single hole injected into the Neel background is $`3.17t+2.83t^{0.27}+1`$, see e. g. Ref., where the last term comes from the $`1/4n_in_j`$ term in the Hamiltonian (1). Therefore in linear in x approximation energy of the doped Neel state is $$E_{Neel}/N=1.17+(3.17t+2.83t^{0.27}+1)x.$$ (68) Energy of the undoped columnar dimerized state without account of quantum fluctuations is $`(0.3750.5)N`$, where $`0.375`$ is Heisenberg energy and $`0.5`$ comes from the $`1/4n_in_j`$ term in the Hamiltonian (1). Quantum fluctuations push this energy down. Using perturbation theory one can find that in linear in $`n_t`$ (triplet density) approximation this shift is $`0.5n_tN`$. According to previous section at the critical point $`n_t0.07`$ and therefore the shift is tiny. Energy of a single hole injected into the spin dimerized background has been found in section III, $`ϵ=e_0+ϵ_{min}`$, where $`e_0`$ is given by eq. (21) and the values of $`ϵ_{min}`$ are presented in Table I. For $`1t4`$ one can fit $`ϵ_{min}`$ as $`ϵ_{min}3.0t`$. Altogether this gives following energy of the doped dimerized state (linear in x approximation) $$E_{Dimer}/N=0.91+(3t+1.75)x.$$ (69) At $`x=0`$ the Neel state energy (68) is lower than that of the dimerized state (69). The critical concentration (the transition point to the dimerized state) is defined by the condition $`E_{Neel}=E_{Dimer}`$. This gives following values of the critical concentration: Table III. $`t/J`$ 1 2 3 $`x_{c1}`$ 0.14 0.11 0.10 The values of $`x_{c1}`$ in Table III are somewhat overestimated. The matter is that the single hole energy for the Neel state is known with high precision, at the same time similar energy for the dimerized state has been found in Section III with trial wave function which contains only 18 components. The true energy is lower than the variational one. To estimate this uncertainty we refer to the doped spin ladder. There is a variational calculation for this system which is similar to to the calculation in Sec. III, and there are also exact numerical simulations . Comparison shows that the variational method underestimate $`ϵ_{min}`$ (see Table I) by 10-15%. Taking this as an estimate we should replace the term $`3t`$ in eq. (69) by $`3.3t`$. Then we come to following values of the critical concentration Table IV. $`t/J`$ 1 2 3 $`x_{c1}`$ 0.12 0.09 0.08 Altogether results for $`x_{c1}`$ presented in Tables II, and III, IV and derived by absolutely different methods are in remarkable agreement. This gives a very strong confirmation of the phase transition to the columnar spin dimerized state. Physical reason for stability of the spin dimerized state is especially evident after the energy considerations: this is the gain in the hole kinetic energy, it is “easier” to propagate in the dimerized background. ## VI Spin-Dimer Order Parameter and Transition to the Normal Fermi Liquid at High Doping. We have discussed the transition to the Neel state at hole concentration $`x<x_{c1}`$. It is clear that at large $`x`$ ($`x>x_{c2}`$) there is a 2nd order phase transition to the normal Fermi liquid. Let us define the spin-dimer order parameter as $$\rho =𝐒_2𝐒_3𝐒_1𝐒_2,$$ (70) where sites 1, 2, and 3 are shown in Fig. 1. For perfect dimerized state $`\rho =3/4`$. There are two mechanisms for reduction of the order parameter. The first one is due to spin quantum fluctuations which approximately give $`\rho 3/4n_t`$. The second mechanism is direct effect of doping. Naive estimate is $`\rho 3/4(12x)`$, however the hole wave function (60) is slightly different from the bare one and because of this the coefficient in the naive formula is slightly renormalized. All together this gives $$\rho =\frac{3}{4}[12x(1+6\alpha ^2)]n_t,$$ (71) where $`\alpha =0.22`$ is the admixture coefficient in eq. (60). This formula is derived in dilute gas approximation, i.e at $`2x,n_t1`$, however for an estimate we can extend it to large $`x`$. Setting $`n_t0.05`$ we find from (71) that $`\rho `$ vanishes at $`x_{c2}0.36`$. We repeat that this is only an estimate because the approach assumes that $`2x1`$. The phase diagram of the $`tJ\delta `$ model at zero temperature is presented in Fig. 7 Because of the mobile holes the dimerized spin liquid at $`x_{c1}<x<x_{c2}`$ is a conducting state. Stability of this state is a very robust effect because it is due to the high energy correlations (typical energy scale $`2t`$). There are also low energy effects with typical energy scale $`2tx`$ which can lead to hole-hole pairing, charge stripes, etc. We do not consider these effects in the present work because they are secondary with respect to the main one: spin dimerization. However we would like to note that there is a simple mechanism for charge stripes induced by the spin dimers: Because of the anisotropic dispersion, see Fig. 3, the charge response is enhanced at some momentum $`𝐩=(p_x,0)`$. The effect is very sensitive to additional hopping parameters $`t^{}`$ and $`t^{\prime \prime }`$, they can further enhance or suppress the response. ## VII Shape of the Fermi Surface and PES Intensity Shape of the Fermi surface for the dimerized state at the doping $`x=0.1`$ and hopping matrix elements from (2) is shown in Fig.3 by the solid line. In this figure we put $`k_x/2\pi `$ along the horizontal axis and $`k_y/\pi `$ along the vertical axis, where $`𝐤`$ is defined on the dimerized lattice. In terms of the Brillouin zone of the original lattice this is usual quadrant: $`0P_x\pi `$, $`0P_y\pi `$. To distinguish between the dimerized lattice and the original one we denote by capital letters momenta corresponding to the original lattice: $`P_x=k_x/2`$, $`P_y=k_y`$. In a real sample there are domains with one stripe dimerization and there are domains with stripes rotated by $`90^o`$. Therefore in an experiment one should see two superimposed Fermi surfaces. Corresponding picture for $`x=0.15`$ is shown in Fig. 8. The first impression is that it is very much different from what is usually observed in angular resolved photoemission (PES) experiments, see e. g. Ref. However let us calculate intensity of the photoemission. The photoeffect operator is of the form $$\widehat{A}=\underset{n}{}c_ie^{i\mathrm{𝐏𝐫}_𝐧},$$ (72) where summation is performed over sites of the square lattice. Amplitude of the hole creation from the dimerized background (oriented as it is shown in Fig. 1) is equal to $$A=s\|b_{}\widehat{A}\|s=\mathrm{cos}(P_x/2).$$ (73) Here we have taken into account that according to considerations in Section III wave function of the hole in bonding state is practically unrenormalized. We stress once more that P is a quasimomentum corresponding to the original lattice. According to (73) intensity of PES spectra ($`IA^2`$) is dropping quickly as $`P_x`$ is increasing. If one assumes that the $`tJ`$ model originates from the single band Hubbard model, then the corrections of the order of $`t/U`$ to the PES intensity can be calculated in a way suggested in Ref. . This gives $$I_𝐏\left(\mathrm{cos}\frac{P_x}{2}+\frac{J}{8t}\mathrm{cos}\frac{3P_x}{2}+\frac{J}{4t}\mathrm{cos}\frac{P_x}{2}\mathrm{cos}P_y\right)^2.$$ (74) The Hubbard repulsion $`U`$ is excluded from this formula using relation $`J=4t^2/U`$. According to (74) the PES intensity is strongly asymmetric at the Fermi surface. For example at $`x=0.15`$ the intensity at the right top corner of the Fermi surface, $`P_x=0.66\pi `$, $`P_y=\pi `$ (see Fig. 8) is 3.5 times smaller than that at the left top corner of the Fermi surface, $`P_x=0.34\pi `$, $`P_y=\pi `$. In real cuprate the asymmetry must be even stronger. The reason for further enhancement of the asymmetry is structure factor of the Zhang-Rice singlet. This is quite similar to the well understood situation in the charge transfer insulator Sr<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub>, see Ref. . Thus the angle resolved photoemission measurements are sensitive only to the “inner” parts (the parts closest to the $`\mathrm{\Gamma }`$ point: $`𝐏=(0,0)`$) of the Fermi surfaces shown in Fig. 8. Shape of this effective “Fermi surface” is very close to what is observed in numerous PES experiments. Another feature which agrees with experiment is width of the “quasiparticle” peak along (1,1) direction: it is always rather big because the peak arises as a superposition of two different peaks corresponding to two different Fermi surfaces. ## VIII Conclusions In conclusion, using the dilute gas approximation we have analyzed the phase diagram of the $`tJ`$ model and the stability of spin-dimerized phase. The main result of the work is phase diagram shown in Fig. 7. Without any explicit dimerization ($`\delta =0`$) the spin dimerized phase is stable at hole concentration $`x_{c1}<x<x_{c2}`$. At $`t/J=3`$ the critical concentrations are $`x_{c1}0.09`$, $`x_{c2}0.36`$. At $`x<x_{c1}`$ the system undergoes transition to the Neel state, and at $`x>x_{c2}`$ to the Normal Fermi liquid. To prove stability of the spin dimerized phase and to calculate critical concentrations we have used two independent approaches. The first one is based on the calculation of the magnon Green’s function. The second approach consists in direct comparison of ground state energies of the Neel state and the dimerized state. Both approaches demonstrate stability of the spin dimerized phase and give very close values of the critical concentrations. ###### Acknowledgements. I thank V. N. Kotov and M. Yu. Kuchiev for stimulating discussions. I am especially grateful to T. M. Rice who attracted my attention to the direct comparison of the ground state energies. ## IX Appendix: “Triple” Diagrams Purpose of the present section is to demonstrate that the “triple” diagrams shown in Fig. 6a,b do not influence position of the critical point $`x_{c1}`$ found in Section IV. The “triple” vertex is shown in Fig. 6c. In the present section we completely neglect small renormalization of the hole wave function considered in Section III. Therefore, the initial state in Fig. 6c is given by (19) and the final state is $$\|f=b_{p\sigma }^{}t_{q\alpha }^{}\|S=\frac{1}{N_2}\underset{n}{}e^{i\mathrm{𝐩𝐫}_𝐧}b_{n\sigma }^{}\underset{m}{}e^{i\mathrm{𝐪𝐫}_𝐦}t_{m\alpha }^{}\|S.$$ (75) Kinematic structure of the vertex is obvious $$\mathrm{\Gamma }_{𝐩,𝐪}=iAt_{𝐪\alpha }^{}[b_𝐩^{}\sigma _\alpha b_{𝐩+𝐪}],$$ (76) where $`\sigma _\alpha `$, $`\alpha =1,2,3`$ is vector of Pauli matrices, and $`b_𝐩^{}`$ is the hole wave function in spinor representation. Direct calculation of the matrix element $`f\|H\|1`$ and comparison with (76) gives following value of the coupling constant $$A=\frac{1}{2}\left[(t+2t^{}\mathrm{cos}p_y)\mathrm{sin}p_x+\frac{j}{2}\mathrm{sin}q_x\right].$$ (77) The normal “triple” magnon self-energy $`\mathrm{\Sigma }_{3n}(𝐪,\omega )`$ is shown in Fig. 6a. Since we are interested in the critical point $`x_{c1}`$ it is enough to calculate the self energy at zero frequency and at the momentum where the spin-wave gap vanishes: $`𝐪=𝐪_0=(0,\pi )`$. Straightforward calculation of the loop at $`t^{}=t^{\prime \prime }=0`$ gives $$\mathrm{\Sigma }_{3n}(𝐪_0,0)=tx.$$ (78) Unfortunately analytical calculation at the values of $`t^{}`$ and $`t^{\prime \prime }`$ from (2) is impossible because at these values the curvature of the hole dispersion along y-direction vanishes. However numerical calculation shows that at $`x=0.10.15`$ and $`t^{}`$, $`t^{\prime \prime }`$ from (2) the ”triple” self-energy is by a factor 1.5 smaller than one given by (78). Comparing the “triple” self-energy with the Brueckner one (63) we find that the ”triple” self-energy is by a factor 7 smaller. This is already enough to neglect the ”triple” contribution. However we would like to demonstrate that suppression of the ”triple” contribution is even stronger: its influence on the point of phase transition is exact(!) zero in linear in x approximation. This interesting fact is related to anomalous ”triple” self energy shown in Fig. 6b. Simple consideration based on the structure of the vertex (76) shows that following exact relation takes place $$\mathrm{\Sigma }_{3a}(𝐪_0,0)=\mathrm{\Sigma }_{3n}(𝐪_0,0).$$ (79) According to (15) the magnon spectrum found without ”triple” diagram is of the form $$\omega _𝐪\sqrt{\stackrel{~}{A}^2\stackrel{~}{B}^2},$$ (80) where $`\stackrel{~}{A}`$ arises from the normal terms in the effective Hamiltonian and $`\stackrel{~}{B}`$ arises from the anomalous terms. At the critical point $`\stackrel{~}{A}=\stackrel{~}{B}`$ and hence the excitation energy vanishes. With account of ”triple” diagrams the relation (80) should be rewritten as $$\omega _𝐪\sqrt{\left(\stackrel{~}{A}+\mathrm{\Sigma }_{3n}\right)^2\left(\stackrel{~}{B}+\mathrm{\Sigma }_{3a}\right)^2}.$$ (81) It is clear that because of (79) the dispersion (81) vanishes exactly at the same point as (80). This proves our statement that the ”triple” self energy does not influence position of the transition point. However away from the transition point when the spin-wave gap is increasing the ”triple” self energy is getting more important. Note that similar situation takes place with ”triple” diagrams considered in Ref. for $`J_1J_2`$ model.
warning/0002/cond-mat0002280.html	ar5iv	text	# Remanent magnetization of high-temperature Josephson junction arrays ## Abstract In this work we study the remanent magnetization exhibited by tridimensional disordered high-T<sub>c</sub> Josephson junction arrays excited by an AC magnetic field. The effect, as predicted by numerical simulations and previously verified for a low-T<sub>c</sub> array of Nb, occurs in a limited range of temperatures. We also show that the magnetized state can be excited and detected by two alternative experimental routines. In a recent work we have demonstrated that Josephson junction arrays (JJAs) fabricated from granular Nb may exhibit a magnetic remanence, $`M_r`$, upon excitation by a magnetic field. As predicted, the magnetized state occurs in a window of temperatures, whose extent depends on the critical current, $`J_c`$, of the junctions. Also, there is a threshold value for the magnetic field in order to drive the JJA to the state where flux is retained after suppression of the field <sup>,</sup>. In Ref. 1 we have also shown that the profile of $`M_r`$ is sensitive to the critical current dispersion, what stresses the prospective use if this effect as a suitable tool for determining the critical current distribution of the array, $`N(J_c)`$. This contribution presents selected parts of a systematic study of the remanent magnetization displayed by our newly produced high-temperature tridimensional disordered arrays (3D-DJJAs), fabricated from granular YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. The experimental results confirm the predictions, revealing that the remanence develops in a limited range of temperature. Granular YBCO material used to fabricate the arrays was prepared employing a modified method of polymeric precursors. This route consists of mixing oxides and carbonates in stoichiometric amounts dissolved in HNO<sub>3</sub>, and then to an aqueous citric acid solution. A metallic citrate solution is then formed, to which ethylene glycol is added, resulting in a blue solution which was neutralized to pH~7 with ethylenediamine. This solution was turned into a gel and subsequently decomposed to a solid by heating at 400 <sup>o</sup>C. The sample was heat-treated at 850 <sup>o</sup>C for 12 h in air with several intermediary grindings, in order to prevent undesirable phase formations. Then, it was pressed into a pellet using controlled uniaxial (5,000 kgf/cm<sup>2</sup>) pressure and sintered at 950 <sup>o</sup>C for 6 h in O<sub>2</sub>. This pellet is a 3D-DJJA, in which the junctions are weakly coupled grains, i.e., weak-links (WLs) formed by a sandwich of YBCO grains and intergrain material. As a consequence of the uniaxial pressure, samples produced in this way are anisotropic, a feature that can be either enhanced, by using higher pressures, or reduced, by applying isostatic pressures. Also, thermal treatment plays a fundamental role on creation and control of WLs and anisotropy, as will be thoroughly discussed elsewhere. The sample studied here exhibits all characteristic features of a genuine 3D-JJA, the most significant of which are shown in Fig.1: the main picture is a low-field measurement of a positive magnetization (Wohlleben Effect), for $`H=0.02Oe`$. The inset displays a Fraunhofer pattern for the real part of the magnetic susceptibility $`\chi _{AC}`$. As demonstrated in Ref. 6, this is an indirect determination of $`J_c`$. To study the remanent state of the arrays, we employ two routines especially developed for detection and study of granular JJAs: the Temperature Scan Routine (TS) and the Field Scan Routine (FS). The core of both experimental procedures consists of two steps: i. the sample is submitted to an AC field ($`h`$) consisting of a train of sinusoidal pulses, after what h is kept null; ii. with $`h=0`$, the magnetic moment of the sample is measured. In the FS routine we measure $`\chi _{AC}(h)`$ performing steps (i) and (ii) as h is varied at a fixed temperature. On the other hand, the TS routine is employed to measure, at a fixed value of $`h`$, $`\chi _{AC}(T)`$ through steps (i) and (ii). All measurements were performed using a Quantum Design MPMS-5T SQUID magnetometer. Both routines were extensively explored, furnishing valuable results for the purposes of this work. In this short paper, however, we emphasize the similarities among results obtained employing the two alternative routines. Remaining parts of this study, including many other aspects of the problem, will be published elsewhere. As expected, the high-T<sub>c</sub> disordered array studied here, a heat-treated isotropic 3D-DJJA of YBCO, exhibits the predicted magnetic behavior. The magnetized state at zero field can be easily recognized from measurements using either one of the above mentioned routines. Remanence versus temperature curves normalized to peak values, are shown in Fig.2. The main graph represents a direct measurement of $`M_r(T)`$ on warming, employing the TS routine. Intragranular contributions were not subtracted from this curve, but are totally irrelevant, as we have certified by measuring unlinked grained material, for which no magnetic remanence was detected. As in the case of the Nb array reported previously, the remanence is intense for a limited temperature interval. In the present case, the temperature window for which the array is magnetically active has a quite long low-T tail, revealing that $`N(J_c)`$ is broad. It should be noticed that the magnetic response ceases at $`T^{}=83K`$, a temperature significantly smaller than $`T_c=90K`$. The inset shows a few representative points of an experiment performed using the FS routine for $`h_o=10mOe`$. The AC field was cycled up to $`3.8Oe`$ and then down to zero, after which values of $`\mathrm{\Delta }M_r`$ were calculated by subtraction between $`M_r(h_o)`$ increasing and decreasing $`h`$. The curve obtained is a replica of that measured directly using the TS routine. In conclusion, we have measured the predicted magnetic remanence of JJAs, using a 3D-DJJA fabricated from granular YBCO. The remanence is intense within a limited interval of temperatures. The $`M_r(T)`$ profile, which is sensitive to the critical current dispersion, reveals a fairly broad $`N(J_c)`$ for the array.
warning/0002/math0002169.html	ar5iv	text	# Minimal resolution of general stable vector bundles on ℙ² ## 1. Introduction In this paper we investigate stable vector bundles on the complex projective plane $`^2`$ by means of their minimal free resolution. The fundamental background for this study is the work of Bohnhorst and Spindler \[BS92\] and their idea to use admissible pairs to characterize the stability of rank-$`n`$ vector bundles on $`^n`$ and to give a stratification of the relative moduli space in constructible subsets. Two main difficulties arise: (i) to state a weak version of the Bohnhorst-Spindler theorems for rank $`r2`$ vector bundles on $`^2`$, (ii) to estimate the codimension of the constructible subsets of the moduli space. In section $`2`$ we make some general remarks on rank $`r`$ vector bundles on $`^n`$ with $`rn`$ to address (i), while (ii) is the object of the last section. We would like to thank G. Ottaviani for his invaluable guidance and V. Ancona for many useful discussions. ## 2. Admissible pairs and resolutions Let (1) $$0\underset{i=1}{\overset{k}{}}𝒪_^n(a_i)\stackrel{\Phi }{}\underset{j=1}{\overset{r+k}{}}𝒪_^n(b_j)0$$ be a free resolution of length $`1`$ of a rank $`r`$ vector bundle $``$ on $`^n`$. We assume that the two sequences $`a_i`$ and $`b_i`$ are indexed in nondecreasing order (2) $`a_1`$ $`a_2\mathrm{}a_k,`$ $`b_1`$ $`b_2\mathrm{}b_k\mathrm{}b_{r+k}.`$ We call $`(a,b)=((a_1,\mathrm{},a_k),(b_1,\mathrm{},b_{r+k}))`$ the pair associated to the resolution (1). If the resolution (1) is minimal, then we call $`(a,b)`$ the pair associate to the bundle $``$. Notice that the associated pair and the Betti numbers of a resolution encode exactly the same information; in particular $`\mathrm{max}(a_k1,b_{r+k})`$ is the regularity. ###### 2.1 Definition. The pair $`(a,b)`$ is said to be weakly admissible if (3) $$a_i>b_{n+i}\text{for all }i=1,\mathrm{},k$$ and admissible (or strongly admissible) if (4) $$a_i>b_{r+i}\text{for all }i=1,\mathrm{},k$$ For brevity we say that the resolution (1) is weakly or strongly admissible if the associated pair $`(a,b)`$ is. ###### 2.2 Example. The map $`\mathrm{\Phi }`$ can be expressed by a $`(r+k)\times k`$ matrix of forms $`(\varphi _{i,j})`$ of degree $`\mathrm{deg}(\varphi _{i,j})=(b_ia_j)`$. If $`(a,b)`$ is a strongly admissible pair and $`\omega _0\mathrm{}\omega _r`$ are linear forms in general position on $`^n`$, then the $`(r+k)\times k`$ matrix (5) $$(\varphi _{ij}):=\left[\begin{array}{ccc}\omega _0^{a_1b_1}& \mathrm{}& 0\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \mathrm{}& & \omega _0^{a_kb_k}\\ \omega _r^{a_1b_{r+1}}& & \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \omega _r^{a_kb_{r+k}}\end{array}\right]$$ defines a minimal free resolution and $`:=\mathrm{coker}\mathrm{\Phi }`$ is a vector bundle with associated pair $`(a,b)`$. If the pair $`(a,b)`$ is only weakly admissible, the same reasoning works for the $`\omega `$’s defined by $$(\omega _0,\mathrm{},\omega _n,\omega _{n+1},\mathrm{},\omega _r)=(x_0,\mathrm{},x_n,\mathrm{\hspace{0.17em}0},\mathrm{},\mathrm{\hspace{0.17em}0})$$ Admissibility was originally introduced by Bohnhorst and Spindler in \[BS92\] to characterize stability for vector bundles on $`^n`$ of homological dimension $`1`$ with rank $`r`$ equal to $`n`$. Note that for $`r=n`$ weakly and strongly admissibility coincide. In this section we are going to restate their results with some generalizations to the case of rank $`rn`$. ###### 2.3 Proposition. If $`rn`$, the following two condition on the resolution (1) are equivalent: 1. is minimal; 2. is weakly admissible and every constant entry of the matrix $`(\varphi _{i,j})`$ is zero. ###### Proof. Obviously, (2.) follows from (1.). For $`r=n`$ the statement was proved by Bohnhorst and Spindler (\[BS92\] proposition 2.3). Now suppose that $`r>n`$ and (1) is minimal. Since $`(b_{r+k})`$ is globally generated, Bertini’s theorem ensures that a generic map $`f:𝒪(b_{r+k})`$ is injective. Then, in the following commutative diagram columns and rows are exact and $`^{\prime \prime }`$ is locally free: (6) $$\begin{array}{ccccccccc}& & 0& & 0& & 0& & \\ & & & & & & & & & & \\ 0& & 0& & 𝒪\left(b_{r+k}\right)& \stackrel{Id}{}& 𝒪\left(b_{r+k}\right)& & 0\\ & & & & & & & & & & \\ 0& & _{i=1}^k𝒪\left(a_i\right)& & _{i=1}^{r+k}𝒪\left(b_i\right)& & & & 0\\ & & Id& & & & & & & & \\ 0& & _{i=1}^k𝒪\left(a_i\right)& & _{i=1}^{r+k1}𝒪\left(b_i\right)& & ^{\prime \prime }& & 0\\ & & & & & & & & & & \\ & & 0& & 0& & 0& & \end{array}$$ The minimality of the middle row yields the minimality of the last row. Using induction on $`r`$, we may assume that the last row is weakly admissible. Then, the middle row is also weakly admissible. ∎ As a consequence, to every vector bundle it corresponds a weakly admissible pair. Vice versa, for every weakly admissible pair, the examples 2.2 provide a way to construct a vector bundle with such associated pair. ###### 2.4 Theorem (Bohnhorst-Spindler \[BS92\]). Suppose $`r=n`$ and that the resolution (1) is admissible. Let $`c_1=a_ib_j`$ be the first Chern class and $`\mu =c_1/n`$ the slope of $``$. Then $``$ is stable (respectively semistable) if and only if (7) $$b_1>\mu \text{(resp. }b_1\mu \text{)}.$$ For the case of higher rank, we have no chances to extend the above arithmetical characterization since stable and unstable vector bundles may have the same associated pair. However, one implication still hold: ###### 2.5 Theorem. If the resolution (1) is weakly admissible (in particular if it is minimal) and $``$ is a stable (resp. semistable) vector bundle, then the associated pair $`(a,b)`$ is strongly admissible and $`b_1>\mu `$ (resp. $`b_1\mu `$). ###### Proof. We first prove that if $``$ is semistable, then $`b_{r+k}<a_k`$. In fact, if $`b_{r+k}a_k`$ then $``$ split as $`=𝒪(b_{k+r})^{\prime \prime }`$ and by weakly admissibility (8) $$\begin{array}{cc}\hfill \underset{i=1}{\overset{k}{}}a_i\underset{j=1}{\overset{r+k1}{}}b_j& =\underset{i=1}{\overset{n}{}}b_i+\underset{i=1}{\overset{k}{}}(a_ib_{n+i})\underset{i=n+k+1}{\overset{r+k1}{}}b_i\hfill \\ & nb_{r+k}+k(rn1)b_{r+k}\hfill \\ & >(1r)b_{r+k}\hfill \end{array}$$ then we have $`\mu (^{\prime \prime })>b_{r+k}=\mu (𝒪(b_{r+k}))`$ which contradicts the semistability of $``$. Now suppose that $``$ is semistable and $`a_sb_{r+s}`$ for some $`s`$ with $`1s<k`$ and let $`(s_01)`$ be the largest of such $`s`$. Since $$a_1\mathrm{}a_{s_0}b_{r+s_0}\mathrm{}b_{r+k},$$ the minor $`\mathrm{\Phi }^{\prime \prime }`$ of $`\mathrm{\Phi }`$, obtained by cutting off the last $`(ks_0)`$ rows and $`(ks_0)`$ columns, remains of maximal rank so $`^{\prime \prime }:=\mathrm{coker}\mathrm{\Phi }^{\prime \prime }`$ is a vector bundle. A surjective morphism $`^{\prime \prime }0`$ is defined by the diagram (9) $$\begin{array}{ccccccccc}0& & _{i=1}^k𝒪\left(a_i\right)& \stackrel{\mathrm{\Phi }}{}& _{j=1}^{r+k}𝒪\left(b_j\right)& & & & 0\\ & & & & & & & & & & \\ 0& & _{i=1}^{s_0}𝒪\left(a_i\right)& \stackrel{\mathrm{\Phi }^{\prime \prime }}{}& _{j=1}^{r+s_0}𝒪\left(b_j\right)& & ^{\prime \prime }& & 0\end{array}$$ where the first two vertical map are the natural projections. Observe that $`^{\prime \prime }`$ have the same rank $`r`$ as $``$ and (10) $$\begin{array}{cc}\hfill \mu (^{\prime \prime })& =\frac{1}{r}\left(\underset{i=1}{\overset{s_0}{}}a_i\underset{j=1}{\overset{r+s_0}{}}b_j\right)=\mu ()\frac{1}{r}\underset{i=s_0+1}{\overset{k}{}}(a_ib_{r+i})\hfill \\ & <\mu ().\hfill \end{array}$$ Then, $`^{\prime \prime }`$ must be semistable otherwise any torsionless quotient sheaf destabilizing it would also destabilize $``$. By induction on $`k`$, we may assume that the second row of (9) is strongly admissible. In particular, we have $`a_{s_0}>b_{r+s_0}`$ that gives a contradiction. Finally, if $``$ is stable (resp. semistable), then (11) $$\mathrm{H}^0((m))=0m\mu ()\text{(resp. }m<\mu ()\text{)}$$ but, from the exact sequence (12) $$0\underset{i=1}{\overset{k}{}}𝒪(a_i+b_1)\underset{j=1}{\overset{r+k}{}}𝒪(b_j+b_1)(b_1)0,$$ we have $`\mathrm{H}^0((b_1))0`$ then $`b_1>\mu ()`$ (resp. $`b_1\mu ()`$). ∎ ## 3. On vector bundles on $`^2`$ By Horrocks theorem, every rank $`r`$ vector bundle $``$ on $`^2`$ has homological dimension at most $`1`$, that is, if $``$ does not split in the direct sum of line bundles, then it is presented by a minimal free resolution of the form (13) $$0\underset{i=1}{\overset{k}{}}𝒪_^2(a_i)\stackrel{\Phi }{}\underset{j=1}{\overset{r+k}{}}𝒪_^2(b_j)0.$$ The Chern classes $`c_1`$, $`c_2`$ of $``$ are determined by $`a_i`$ and $`b_j`$ with the formulas (14) $`c_1`$ $`={\displaystyle \underset{i=1}{\overset{k}{}}}a_i{\displaystyle \underset{i=1}{\overset{k+r}{}}}b_i,`$ $`2c_2c_1^2`$ $`={\displaystyle \underset{i=1}{\overset{k}{}}}a_i^2{\displaystyle \underset{i=1}{\overset{k+r}{}}}b_i^2.`$ We denote by $``$ the set of all (strongly) admissible pairs $`(a,b)`$ associated to rank $`r`$-vector bundles on $`^2`$ with Chern classes $`c_1`$, $`c_2`$ satisfying the condition $`b_1>\mu =(a_ib_j)/r`$. Theorem 2.5 shows that the set $``$ contains the set of all possible associated pair to a stable vector bundle in $`𝔐_^2(r,c_1,c_2)`$ and coincides exactly with it for $`r=n`$. Then (15) $$𝔐_^2(r,c_1,c_2)=\underset{(a,b)}{}𝔐(a,b)$$ where $`𝔐(a,b)`$ will be the subset (possibly empty) of $`𝔐_^2(r,c_1,c_2)`$ of vector bundles with associated pair $`(a,b)`$. The following result was stated and proved by Bohnhorst and Spindler \[BS92\] for rank-$`n`$ vector bundles on $`^n`$ with homological dimension $`1`$, but their proof works on $`^2`$ for vector bundles of any rank without modifications. ###### 3.1 Theorem. For all $`(a,b)`$, the closed set $`\overline{𝔐(a,b)}`$ is an irreducible algebraic subset of $`𝔐_^2(r,c_1,c_2)`$ of dimension: (16) $$\begin{array}{cc}\hfill dim\overline{𝔐(a,b)}& =dim\mathrm{Hom}(F_1,F_0)+dim\mathrm{Hom}(F_0,F_1)\hfill \\ & dim\mathrm{End}(F_1)dim\mathrm{End}(F_0)+1\mathrm{\#}\{(i,j):a_i=b_j\},\hfill \end{array}$$ where $`F_0=_{j=1}^{k+r}𝒪(b_j),`$ $`F_1=_{i=1}^k𝒪(a_i).`$ ## 4. Natural pairs and general vector bundles We say that $`(a,b)=((a_1,\mathrm{},a_k),(b_1,\mathrm{},b_{r+k}))`$ is a natural pair if it is admissible and (17) $$b_{r+k}<a_1,a_kb_1+2.$$ Through this section, we are going to show that resolutions of general vector bundles have natural pairs: ###### 4.1 Theorem. One has $`\mathrm{codim}\overline{𝔐(a,b)}=0`$ if and only if $`(a,b)`$ is a natural pair. As a remarkable consequence we will derive a quite simple proof of the irreducibility of the moduli spaces of stable vector bundles on $`^2`$ (other proofs with different techniques can be found in \[Bar77a\], \[Ell83\]), \[HL93\], \[LeP79\], \[Mar78\]) and we will compute the regularity and the cohomology of their general elements. We recall that, since $`dim\mathrm{Ext}^2(,)=0`$ for any stable vector bundle $``$ on $`^2`$, the relative moduli $`𝔐_^2(r,c_1,c_2)`$ space is smooth of dimension (18) $$dim\mathrm{Ext}^1(,)=2rc_2(r1)c_1^2r^2+1.$$ Let us consider the function $`A(t):=h^2(𝒪(t))`$ and the finite differences of first and second order $`(\mathrm{\Delta }_uA)(t):=A(t+u)A(t)`$ and $`(\mathrm{\Delta }_v\mathrm{\Delta }_uA)(t):=(\mathrm{\Delta }_uA)(t+v)(\mathrm{\Delta }_uA)(t).`$ ###### 4.2 Lemma. Let $`(a,b)`$ be the admissible pair associated to a stable vector bundle $``$ on $`^2`$. Then (19) $$\begin{array}{cc}\hfill \mathrm{codim}\overline{𝔐(a,b)}& =\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \\ & +\underset{i,j=1}{\overset{r}{}}(\mathrm{\Delta }_{b_{i+r}a_i}\mathrm{\Delta }_{b_{j+r}a_j}A)(a_ib_{j+r})\hfill \end{array}$$ ###### Proof. Let (20) $$0F_1F_00$$ be the minimal resolution of $``$ where (21) $$F_0=\underset{j=1}{\overset{k+r}{}}𝒪(b_j),F_1=\underset{i=1}{\overset{k}{}}𝒪(a_i).$$ The stability of $``$ ensures the vanishing $`dim(\mathrm{Ext}^2(,))=h^2(^{})=0`$ so that $`h^2(F_{0}^{}{}_{}{}^{})=h^2(F_{1}^{}{}_{}{}^{})`$. Then, from (20) we easily find the following data: (22) $$\begin{array}{cc}\hfill h^0(F_{0}^{}{}_{}{}^{})& =h^0(F_{0}^{}{}_{}{}^{}F_0)h^0(F_{0}^{}{}_{}{}^{}F_1),\hfill \\ \hfill h^0(F_{1}^{}{}_{}{}^{})& =h^0(F_{1}^{}{}_{}{}^{}F_0)h^0(F_{1}^{}{}_{}{}^{}F_1),\hfill \\ \hfill dim(\mathrm{Ext}^1(,))& =h^1(^{})=\hfill \\ & =h^1(F_{0}^{}{}_{}{}^{})h^1(F_{1}^{}{}_{}{}^{})+\hfill \\ & +h^0(F_{1}^{}{}_{}{}^{})h^0(F_{0}^{}{}_{}{}^{})+1\hfill \end{array}$$ and from (16) we have (23) $$\begin{array}{cc}\hfill \mathrm{codim}\overline{𝔐(a,b)}& =dim(\mathrm{Ext}^1(,))dim\overline{𝔐(a,b)}=\hfill \\ & =h^1(F_{0}^{}{}_{}{}^{})h^1(F_{1}^{}{}_{}{}^{})+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \end{array}$$ Now, by splitting $`F_0`$ as $`𝒪(b_1)\mathrm{}𝒪(b_{r1})\stackrel{~}{F_0}`$ with $`\stackrel{~}{F_0}:=_{i=r+1}^{k+r}𝒪(b_i)`$, the above formula becomes (24) $$\begin{array}{cc}\hfill \mathrm{codim}\overline{𝔐(a,b)}& =\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \\ & +h^1(\stackrel{~}{F_0}^{})h^1(F_{1}^{}{}_{}{}^{})\hfill \\ & =\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \\ & +h^2(\stackrel{~}{F_0}^{}F_1)h^2(\stackrel{~}{F_0}^{}F_0)\hfill \\ & h^2(F_{1}^{}{}_{}{}^{}F_1)+h^2(F_{1}^{}{}_{}{}^{}F_0).\hfill \end{array}$$ Since $`h^2(\stackrel{~}{F_0}^{}F_0)=h^2(\stackrel{~}{F_0}^{}\stackrel{~}{F_0})`$ and $`h^2(F_{1}^{}{}_{}{}^{}F_0)=h^2(F_{1}^{}{}_{}{}^{}\stackrel{~}{F_0})`$ so (25) $$\begin{array}{cc}\hfill \mathrm{codim}\overline{𝔐(a,b)}& =\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \\ & +h^2(\stackrel{~}{F_0}^{}F_1)h^2(\stackrel{~}{F_0}^{}\stackrel{~}{F_0})\hfill \\ & h^2(F_{1}^{}{}_{}{}^{}F_1)+h^2(F_{1}^{}{}_{}{}^{}\stackrel{~}{F_0}).\hfill \end{array}$$ Finally, distributing the direct sums appearing in the definition of $`\stackrel{~}{F}_0`$ and $`F_1`$ the equation (25) becomes (26) $$\begin{array}{cc}\hfill \mathrm{codim}\overline{𝔐(a,b)}& =\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \\ & +\underset{i,j=1}{\overset{r}{}}[h^2(𝒪(b_{i+r}a_j))h^2(𝒪(b_{i+r}b_{j+r}))\hfill \\ & h^2(𝒪(a_ia_j))+h^2(𝒪(a_ib_{j+r}))]\hfill \\ & =\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}\hfill \\ & +\underset{i,j=1}{\overset{r}{}}(\mathrm{\Delta }_{b_{i+r}a_i}\mathrm{\Delta }_{b_{j+r}a_j}A)(a_ib_{j+r})\hfill \end{array}$$ We observe that natural pairs are parametrized by three integers $`s`$, $`k`$, $`\alpha `$ such that (27) $$k1\text{and}k+1\alpha k+r$$ as follows: the pair $`(a,b)_{s,k,\alpha }`$ corresponding to the triple $`(s,k,\alpha )`$ is the pair associated to a resolution of the form (28) $$0𝒪(s1)^k𝒪(s)^\alpha 𝒪(s+1)^{r+k\alpha }0$$ if $`\alpha 0`$, or of the form (29) $$0𝒪(s1)^{k+\alpha }𝒪(s)^\alpha 𝒪(s+1)^{r+k}0$$ if $`\alpha <0`$. We exclude the case $`\alpha =k`$ so that $`s`$ is the regularity of the pair, i.e. $`s=\mathrm{max}(a_k1,b_{r+k})`$. ###### Proof of theorem 4.1. It can be verified by direct computation from theorem 3.1 that, if $``$ has natural pair, then the codimension of $`\overline{𝔐(a,b)}`$ is zero. Conversely, let $`u`$, $`v`$ be two non-negative integers. Since all finite difference $`(\mathrm{\Delta }_uA)(t):=A(t+u)A(t)`$ are non decreasing functions of $`t`$, then (30) $$(\mathrm{\Delta }_v\mathrm{\Delta }_uA)(t)0$$ and by the previous lemma (31) $$\mathrm{codim}\overline{𝔐(a,b)}\underset{i=1}{\overset{r}{}}h^1((b_i))+\mathrm{\#}\{(i,j):a_i=b_j\}.$$ If $`\mathrm{codim}\overline{𝔐(a,b)}=0`$, we have $`a_kb_1+2`$ and $`\mathrm{\#}\{(i,j):a_i=b_j\}=0`$, since $`h^1((b_1))=0`$ implies $`h^2(F_1(b_1))=0`$. This forces $`(a,b)`$ to be a natural pair. ∎ ###### 4.3 Proposition. Let $`𝔐_^2(r,c_1,c_2)`$ be nonempty and (32) $$s:=\mathrm{max}\{\rho :r\rho ^2+2c_1\rho r\rho 2c_2c_1^2+c_11\},$$ or, equivalently, (32bis) $$s:=\mathrm{min}\{\rho :r\rho ^2+2c_1\rho +r\rho 2c_2c_1^2c_1\}.$$ If $`\alpha `$ and $`k`$ are defined by (33) $`\alpha `$ $`:=2c_2c_1^2+rrs^22c_1s,`$ $`k`$ $`:=(rs+c_1r+\|\alpha \|)/2,`$ then $`(a,b)_{s,k,\alpha }`$ is the only natural pair of $`𝔐_^2(r,c_1,c_2)`$. ###### Proof. This is a verification; we outline the main steps of the computation. In the first place, one must ensure that the natural pair $`(a,b)_{s,k,\alpha }`$ is actually associated to vector bundles in $`𝔐_^2(r,c_1,c_2)`$. This amount to show that from (14) the pair $`(a,b)_{s,k,\alpha }`$ has the appropriate Chern classes and that conditions (27) hold. From theorem 4.1, any pair $`(a,b)`$ such that $`dim\overline{𝔐(a,b)}=0`$ is a natural pair of the form $`(a,b)_{s,k,\alpha }`$. From resolutions (28) and (29) we find that $`\alpha `$, $`k`$ must satisfy (33). Then, it remains to verify that $`s`$ is uniquely determined from $`r`$, $`c_1`$, $`c_2`$ and satisfy (32). By substitution, the inequalities $`k<\alpha k+r`$ turn into (34) $$rs^2+2c_1sc_1rs+12c_2c_1^2rs^2+2c_1s+c_1+rs.$$ Since the intervals $`[rs^2+2c_1sc_1rs+1,rs^2+2c_1s+c_1+rs]`$ are disjoint for $`s`$ varying in $``$, then equations (32) and (32bis) give the only suitable value for $`s`$. ∎ ###### 4.4 Theorem. The moduli spaces of stable rank $`r`$ vector bundles on $`^2`$ are irreducible. ###### Proof. Moduli space of stable rank $`r`$ vector bundles on $`^2`$ are smooth. By the previous proposition they can have only one connected component. ∎ ###### 4.5 Corollary. The general element of $`𝔐_^2(r,c_1,c_2)`$ has natural cohomology. The above corollary justify the terminology “natural pair”. A different proof for it can be found in \[HL93\]. Now, we are going to give some inequalities on the regularity and the cohomology of stable vector bundles using proposition 4.3. In particular, for rank $`2`$ vector bundles, the next two corollaries give respectively a refined version of corollary 5.4 \[Bru80\] and proposition 7.1 \[Har78\]. ###### 4.6 Corollary. A general vector bundle $``$ in $`𝔐_^2(r,c_1,c_2)`$ has regularity $`\mathrm{reg}()=s`$, where $`s`$ is given by (32). ###### 4.7 Corollary. Let $`[]`$ be a vector bundle in $`𝔐=𝔐_^2(r,c_1,c_2)`$ and $`s`$ defined by (32). Then $`\mathrm{H}^0((t))0`$ if $`t`$ $`s`$ $`\text{when}rs^2+2c_1s+rs=2c_2c_1^2c_1,`$ $`t`$ $`s1`$ otherwise. The above inequality is sharp, in the sense that, if $``$ is general, it gives a necessary and sufficient condition. ###### Proof. Let $`((a_1,\mathrm{},a_k),(b_1,\mathrm{},b_{k+r}))`$ be the admissible pair associated to a vector bundle $``$ in $`𝔐`$. Then, one has $`\mathrm{H}^0((t))0`$ if and only if $`tb_10`$. By semicontinuity of cohomology groups and theorem 4.4, it is enough to restrict ourselves to the case where $``$ is general. So, by (28) and (29) one has $`\mathrm{H}^0((t))0`$ if and only if $`t`$ $`s`$ $`\text{if}\alpha =k+r`$ $`t`$ $`s1`$ otherwise and the condition $`\alpha =k+r`$ is equivalent to $`rs^2+2c_1s+rs=2c_2c_1^2c_1`$ by (33). ∎
warning/0002/astro-ph0002085.html	ar5iv	text	# 1 INTRODUCTION ## 1 INTRODUCTION The observational evidence for the existence of supermassive black holes ($`10^6`$$`10^9`$ times the mass of the sun, M) in the centers of active galaxies has been accumulating at an ever accelerating pace for the last few decades (e.g., Rees 1998; Blandford & Gehrels 1999). Seyfert (1943) first drew attention to a group of galaxies with unusual excitation conditions in their nuclei indicative of energetic activity. Among the twelve galaxies in his list was NGC 4258, which is the subject of much of this paper. Such galaxies, now known as galaxies with active galactic nuclei (AGN), have grown in membership and importance. Ironically, NGC 4258 no longer belongs to the class of Seyfert galaxies by modern classification standards (Heckman 1980), but it is still considered to have a mildly active galactic nucleus. Meanwhile, the study of AGN has become a major field in modern astrophysics. In the 1960s, galaxies with AGN were discovered with intense radio emission arising from jets of relativistic particles often extending far beyond the optical boundaries of the host galaxy. The central engine, the source of energy that powers such jets and other phenomena in the centers of galaxies, has long been ascribed to black holes (e.g., Salpeter 1964; Blandford & Rees 1992). There are two sources of energy for these phenomena: the gravitational energy from material falling onto the black hole and the spin energy of the black hole itself (Blandford & Znajek 1977). The direct evidence for black holes in AGN has come principally from observations of the motions of gas and stars in the extended environments of black holes. In the optical and infrared domains, the evidence for black holes from stellar measurements comes from an analysis of the velocity dispersion of stars as a function of distance from the dynamical centers of galaxies. In the case of our own Galactic center, the proper motions (angular velocities in the plane of the sky) of individual stars can be measured. These data show that there is a mass of about $`2.6\times 10^6`$ M within a volume of radius 0.01 pc (Genzel et al. 1997; Ghez et al. 1998). In addition, measurements by the Hubble Space Telescope of the velocity field of hydrogen gas in active galaxies indicate the presence of massive centrally condensed objects. Reviews of these data have been written by Faber (1999), Ho (1999), Kormendy and Richstone (1995), and others. In the X-ray portion of the spectrum, there is compelling evidence for black holes in AGN from the detection of the highly broadened iron K$`\alpha `$ line at 6.4 keV. The line is broadened by the gravitational redshift of gas as close as 3 Schwarzschild radii from the black hole. An example of an iron line profile in the galaxy MCG -6-30-15 is shown in Figure 1 (Tanaka et al. 1995). The linear extent of the emission region cannot be determined directly by the X-ray telescope, so it is not possible to estimate directly the mass of the putative black hole. Detailed analysis of the line profile suggests that the black hole is spinning (e.g., Bromley, Miller, & Pariev 1998). In the radio regime, a new line of inquiry has given unexpectedly clear and compelling evidence for black holes: the discovery of water masers orbiting highly massive and compact central objects. With the aid of very long baseline interferometry (VLBI), which provides angular resolution as fine as 200 microarcseconds ($`\mu `$as) at a wavelength of 1.3 cm and spectral resolution of 0.1 km s<sup>-1</sup> or less, the structure of accreting material around these central objects can be studied in detail. This paper describes the observations and the significance of these measurements of water masers in AGN. We begin with a brief description of cosmic masers and the interferometric techniques used to observe them. ## 2 COSMIC MASERS Intense maser action in cosmic molecular clouds was discovered in 1965 (Weaver et al. 1965) from observations of OH, and later from H<sub>2</sub>O, SiO, and CH<sub>3</sub>OH. In the case of water vapor, the commonly observed masers emit in the 6<sub>16</sub>–5<sub>23</sub> transition at 22235 MHz (1.35 cm wavelength). Most masers have been found to be associated with one of two types of objects: newly formed stars or evolved stars (e.g., Elitzur 1992; Reid & Moran 1988). Although very distinct, they share the characteristic of having envelopes of outflowing gas and dust (silicate material). The pump source in all cases is thought to come in the form of either shock waves or infrared radiation. More recently, masers have been found in the spiral arms of nearby galaxies and in AGN. Cosmic masers are similar to their laboratory counterparts on earth in that their intense radiation is produced by population inversion. However, cosmic masers are one-pass amplifiers and have little temporal or spatial coherence. The intensity of cosmic masers varies, often erratically, on timescales of hours to years. The underlying electric field is a Gaussian random process (Moran 1981). In an unsaturated maser (in astronomical terminology), the pumping is sufficiently strong that the microwave intensity does not affect the level populations, and the intensity increases exponentially through the masing medium. The input signal can be either a background source or the maser’s own spontaneous emission. In a saturated maser, one pump photon is needed for each microwave photon, and in a one-dimensional maser medium where beaming can be neglected, the intensity increases linearly with distance. Maser emission is expected to be beamed. Most masers are thought to be saturated, and this condition requires the least pump power. However, this assertion is difficult to prove observationally. Consider a simple geometry for a maser, a filamentary tube, shown in Figure 2. The boundaries of the filament can be defined in terms of either gas density or the region where the line-of-sight velocity is constant to within the thermal line width. The gas in the tube is predominantly molecular hydrogen, with trace amounts of water vapor (about one part in $`10^5`$) and other constituents. If this maser medium is saturated, then the luminosity is given by $$L=h\nu n\mathrm{\Delta }PV,$$ (1) where $`h`$ is Planck’s constant, $`\nu `$ is the frequency, $`n`$ is the population density in pumping level, $`\mathrm{\Delta }P`$ is the differential pump rate per molecule, and $`V`$ is the volume of the masing cloud. This is the most luminosity a maser of given pump rate and volume can produce. The emission will be beamed along the major axis of the filament into an angle $$\beta \frac{d}{L},$$ (2) where $`d`$ is the cross-sectional diameter, and $`L`$ is the length of the filament. The beam angle and the length of the maser are not directly observable. Since the volume of the maser is approximately $`d^2L`$, the observed flux density from a maser beamed toward the earth is $$F_\nu =\frac{1}{2}h\nu n\frac{\mathrm{\Delta }P}{\mathrm{\Delta }\nu }\frac{V}{D^2\mathrm{\Omega }}=\frac{1}{2}h\nu n\frac{\mathrm{\Delta }P}{\mathrm{\Delta }\nu }\frac{L^3}{D^2},$$ (3) where $`\mathrm{\Delta }\nu `$ is the line width, $`D`$ is the distance between the maser and the observer, and $`\mathrm{\Omega }`$ is the beam angle of the emission, $`\beta ^2`$. The maximum allowable hydrogen number density is about $`10^{10}`$ molecules per cubic centimeter, above which the maser levels become thermalized by collisions. This maximum allowable density ($`\rho _c=3\times 10^{13}`$ g cm<sup>-3</sup>, an important parameter in much of the following discussion) along with the maximum allowable pump rate, which equals the Einstein A-coefficient for infrared transitions linking the maser levels, limit the luminosity of a maser of given volume. Water masers outside our Galaxy were first discovered in the spiral arms of the nearby galaxy M33 by Churchwell et al. (1977). Their properties were found to be similar to masers found in Galactic star-forming regions. Much more luminous water masers were found in the AGN associated with NGC 4945 and the Circinus galaxy by dos Santos & Lepine (1979) and Gardner & Whiteoak (1982), respectively. Claussen, Heiligman, & Lo (1984) and Claussen & Lo (1986) conducted surveys and found five additional masers associated with AGN, including the one in NGC 4258 (see Figure 3). They suggested that these masers might arise in gas associated with dust-laden molecular tori that had been proposed to surround black holes by Antonucci & Miller (1985). Nakai, Inoue, & Miyoshi (1993), with a powerful new spectrometer of 16,000 channels spanning a velocity range of 3000 km s<sup>-1</sup>, observed NGC 4258 and discovered satellite line clusters offset from the systemic velocity by about $`\pm `$ 1000 km s<sup>-1</sup>, which are shown in Figure 3 (see also Miyoshi 1999). The compelling reason that the radio emission from the water vapor transition arises from the maser process is straightforward. A typical example of a maser line from a small part of the spectrum of NGC 4258 is shown in Figure 4. The line width is about 1 km s<sup>-1</sup>, the thermal broadening expected for a gas cloud at 300 K. However, the angular size determined by radio interferometry is less than 100 $`\mu `$as, implying that the equivalent blackbody temperature must exceed $`10^{14}`$ K. The exceedingly high brightness of the radiation is the principal evidence for the maser process. (Typical molecular lines from molecular clouds have velocities of several tens of km s<sup>-1</sup> — due to thermal and turbulent broadening — and brightness temperatures of less than 100 K.) Cosmic masers produce very bright spots of radiation but have little else in common with terrestrial masers. It is difficult to use masers to determine physical conditions (e.g., temperature, density) in molecular clouds because of the complexity of the maser process. However, as compact sources of narrowband radiation, masers are ideal probes of the dynamics of their environment. ## 3 VLBI The most important tool for the study of the angular structure of masers is very long baseline interferometry (VLBI). Signals from a maser, or from other bright compact radio sources, are converted to a low-frequency baseband and recorded in digital format on magnetic tape at Nyquist sampling rates of up to about $`10^8`$ samples per second at widely separated telescopes that operate independently. They form a radio version of the classical Michelson stellar interferometer, whose coherence is maintained by the use of atomic frequency standards to preserve the signal phase and timing (see Figure 5). The received signals (which are proportional to the incident electric fields) from an array of two or more telescopes are cross-correlated pairwise to form cross-correlation functions. Taking advantage of the earth’s rotation, the spatial cross-correlation function of the incident electromagnetic field, or visibility, can be measured over a wide range of projected baseline vectors. The image and fringe visibility functions are related through a Fourier transform (see Thompson, Moran, & Swenson 1986). The temporal Fourier transform of the cross-correlation function gives the cross-power spectrum of the radiation, or visibility as a function of frequency, so that images at different frequencies can be obtained. The intrinsic angular resolution, $`\theta `$, of a multielement interferometer is $`0.7\lambda /B`$, where $`\lambda `$ is the wavelength, and $`B`$ is the longest baseline length. For water vapor, $`\lambda `$ = 1.35 cm, and $`B`$ is typically 6000 km, which gives a resolution of 200 $`\mu `$as. The spectral resolution available is typically about 15 KHz, or about a fifth of the line widths (see Figure 4). In maser sources, one spectral feature at a particular frequency or velocity can be used as a phase reference for the interferometer, and all other phases referred to it. With this technique, the coherence time of the interferometer can be extended indefinitely, and the relative positions of masers with respect to the reference feature can be measured to a small fraction of the fringe spacing, or intrinsic resolution. The relative position of an unresolved maser component can be measured to an accuracy of about $$\mathrm{\Delta }\theta =\frac{1}{2}\frac{\theta }{SNR},$$ (4) where $`SNR`$ is the signal-to-noise ratio. ## 4 THE STUDY OF NGC 4258 The imaging of the maser in NGC 4258 was one of the first projects undertaken by a dedicated VLBI system known as the Very Long Baseline Array (VLBA) (see Figure 6) in the spring of 1994. Previous VLBI measurements of the systemic features had shown that they arose from an elongated structure with a velocity gradient along the major axis, highly suggestive of a rotating disk seen edge-on (Greenhill et al. 1995a). The VLBA measurements of all the maser components provide convincing evidence for a rotating disk around a massive central object. The basic observational results on NGC 4258 obtained over the past few years can be summarized as follows (see Table 1 for a list of parameters): 1. The masers appear to trace a highly elongated, although slightly curved, structure (Figure 7). The high-velocity, redshifted and blueshifted features are offset in position on the left and right sides of the systemic features, respectively. The velocities of the high-velocity features as a function of impact parameter (position along the major axis of the distribution) follow the prediction of Kepler’s third law of orbital motion. The systemic features show linear dependence with impact parameter (Miyoshi et al. 1995). 2. The distribution of maser features in the direction normal to the major axis is too small to be measured at present (see Figure 7). The upper limit on the ratio of thickness to radius of the disk is 0.0025 (Moran et al. 1995). 3. The upper limit of any toroidal component of the magnetic field in the masers, derived from searches for Zeeman splitting in the line at 1306 km s<sup>-1</sup>, is less than 300 mG (Herrnstein et al. 1998a). 4. The accelerations (i.e., the linear drift in the line-of-sight velocity with time) of the systemic features are about 9 km s<sup>-1</sup> yr<sup>-1</sup> (Haschick, Baan, & Peng 1994; Greenhill et al. 1995b; Nakai et al. 1995). The high-velocity features that have been tracked have accelerations in the range $`\pm 0.8`$ km s<sup>-1</sup> yr<sup>-1</sup> (Bragg et al. 1999). 5. The high-velocity features show no proper motions with respect to a fixed-velocity component in the systemic range (Herrnstein 1996). The systemic features show proper motions of about 32 $`\mu `$as yr<sup>-1</sup> (Herrnstein et al. 1999). 6. There is an elongated continuum radio source, which appears to be a jet emanating from the black hole position, parallel to the axis of rotation (Herrnstein et al. 1997). There is no 1.35 cm wavelength emission from the position of the black hole (Herrnstein et al. 1998b). There is virtually no doubt that the masers trace a very thin disk in nearly perfect Keplerian motion. Five of six phase-space parameters have been measured for each maser spot, two spatial coordinates and three velocity coordinates. The missing coordinate is the position along the line of sight, which must be inferred from the constraint provided by Kepler’s third law. The approximate placement of the masers in the disk can be understood by considering a simple thin, flat disk viewed edge-on. In this case the line-of-sight velocity, $`v_z`$, of a maser will be given by $$v_zv_0=\sqrt{\frac{GM}{R}}\mathrm{sin}\varphi ,$$ (5) where $`v_0`$ is the line-of-sight velocity of the central object (i.e., the systemic velocity), $`G`$ is the gravitational constant, $`R`$ is the distance of a maser component from the black hole, and $`\varphi `$ is the azimuth angle in the disk, measured from the line between the black hole and the observer. If the disk were randomly filled with observable masers, one might expect to see a velocity position diagram as shown in Figure 8. The linear boundaries of the distribution are populated by masers at the inner and outer edges of the annular disk. The masers on the curved boundaries lie on the midline, where $`\varphi `$ = 90. Hence, the masers in NGC 4258 have a very specific distribution: the high-velocity masers lie close to the midline, and the systemic masers lie within a narrow range of radii. From Equation 5, the radius of a particular maser can be determined as $$R=(GM)^{\frac{1}{3}}\left[\frac{b}{v_zv_0}\right]^{\frac{2}{3}},$$ (6) where $`b`$ is the projected distance on the sky along the major axis from the center of the disk ($`\mathrm{sin}\varphi =b/R`$). Similarly, positional offsets from the midline, $`z`$, of the high-velocity features can be determined by deviations from a Keplerian curve; that is, $$z=\sqrt{R^2b^2}.$$ (7) There is a two-fold ambiguity in the $`z`$ component of the position for the edge-on disk case. Unambiguous estimates of the positions of the high-velocity features have been derived from the accelerations by Bragg et al. (1999), who showed that the masers lie within 15 degrees of the midline. An expanded plot of the Keplerian part of the velocity curve is shown in Figure 9. The data fit a Keplerian curve to an accuracy of about 3 km s<sup>-1</sup>, or less than 1 percent of the rotation speed. However, there are noticeable deviations from a perfect fit. The estimate of the central mass of the disk derived from this data depends on the distance to the maser, and has a value of $`3.9\times 10^7`$ M for a distance of 7.2 Mpc. This mass corresponds to an Eddington luminosity (where radiation pressure from Thomson scattering would balance gravity) of $`5\times 10^{45}`$ erg s<sup>-1</sup>. Since the total electromagnetic emission appears to be less than $`10^{42}`$ erg s<sup>-1</sup>, the system is highly sub-Eddington. Since this binding mass must lie inside the inner radius of the maser disk, the mass density, assuming a spherical mass distribution, is at least $`4\times 10^9`$ M$`_{}`$pc<sup>-3</sup> ($`3\times 10^{13}`$ g cm<sup>-3</sup>). It is unlikely that this mass is in the form of a dense star cluster (Maoz 1995). The average density of stars in the solar neighborhood is about 1 M$`_{}`$pc<sup>-3</sup>, and the density of the densest known star cluster is about $`10^5`$ M$`_{}`$pc<sup>-3</sup>. A star cluster will have a mass distribution that decreases monotonically with radius. In order not to disrupt the Keplerian curve, the core mass for a reasonable distribution must have a peak density of at least $`1\times 10^{12}`$ M$`_{}`$pc<sup>-3</sup>. A cluster of massive stars at this density would evaporate from gravitational interactions on a timescale short with respect to the age of the galaxy, while a cluster of low-mass stars would destroy itself from collisions over a similar timescale. Hence, it is unlikely that the central mass is in the form of a star cluster (see also Begelman & Rees 1978). The best explanation is that the central object is a supermassive black hole, with a Schwarzschild radius ($`R_S`$) of about $`1.2\times 10^{13}`$ cm. Hence, the masers are distributed in a zone between 40,000 and 80,000 R<sub>S</sub>. Because the maser clouds are so far from the event horizon, deviations of their motions from the predictions of Newtonian mechanics are small. The gravitational redshift and transverse Doppler shift are about 4 km s<sup>-1</sup> (detectable), the expected Lense-Thirring precession (see below) is less than about 3over the maser annulus (possibly detectable), and the apparent shift of the maser positions due to gravitational bending is about 0.1 $`\mu `$as (undetectable). The disk is remarkably thin. In a disk supported against gravity by pressure (hydrostatic equilibrium), the density distribution is expected to have a Gaussian profile with a thickness, $`H`$, given by the relation $$H/R=(c_s^2+v_a^2)^{\frac{1}{2}}/v_\varphi ,$$ (8) where $`c_s`$ is the sound speed and $`v_a`$ is the Alfvén speed, which characterize thermal and magnetic support pressure, respectively, and $`v_\varphi `$ is the Keplerian rotational speed. Since $`H/R<0.0025`$, the quadrature sum of the sound speed and Alfvén speed is less than 2.5 km s<sup>-1</sup>. The upper limit on the magnetic field of 300 mG suggests that the Alfvén speed, $`B/\sqrt{4\pi \rho }`$, where $`\rho `$ is the density, is less than 3 km s<sup>-1</sup> for $`\rho =\rho _c`$ (the critical density for quenching maser emission). If the support were completely due to thermal pressure, the temperature would be less than 1000 K. A proper determination of the positions of the masers on the disk requires that the warp and the inclination of the disk to the line of sight be taken into account. An example of such a disk, slightly warped (in position angle only) and slightly inclined to the line of sight, that fits the maser distribution in position and velocity is shown in Figure 10. The distance to the maser of 7.2 $`\pm `$ 0.3 Mpc was determined from analysis of the proper motions and accelerations of the systemic features (Herrnstein et al. 1999). Fifteen features were tracked over a period of two years to an accuracy of 0.5–10 $`\mu `$as in relative position and 0.4 km s<sup>-1</sup> in velocity. The distance estimate is based on simple geometric considerations. The Keplerian curve of the high-velocity masers gives the mass function $`GM\mathrm{sin}^2i/D`$, where $`i`$ is the inclination of the disk to the line of sight. The radius, $`R`$, of the systemic masers (in angular units) is determined from Equation (6), based on the slope of the velocity versus impact parameter curve shown in Figure 7. This fixes the angular velocity, $`v_\varphi `$, of the systemic masers under the assumption that the orbits are circular. The expected accelerations and proper motions of the systemic features are $`v_\varphi ^2/R`$ and $`v_\varphi /D`$, respectively. The assumption that the orbits are circular is reasonable on theoretical grounds because of viscous relaxation and on observational grounds because the continuum emission arises close to the center of symmetry of the maser distribution. The distance to the 15 Cepheid variables in NGC 4258 has been estimated to be $`8.1\pm 0.8`$ Mpc (Maoz et al. 1999). The statistical component of the error is 0.4 Mpc and the systematic error associated with the calibration of the Cepheid distance scale is 0.7 Mpc. The discrepancy between the two distance measurements to NGC 4258 may have cosmological implications (Paczynski 1999). ## 5 INTERESTING UNANSWERED QUESTIONS 1. What is the rate of radial inflow of material through the disk? The accretion rate of material onto the black hole is an important parameter that affects our understanding of radiation processes around the black hole. The maser data provide some information about the accretion rate. Key issues are the long timescale needed for material to flow from the disk to the black hole and the assumption that the masers trace all the disk material. It is useful to first estimate the mass of the disk. If we account for systematic effects in addition to the random scatter of 3 km s<sup>-1</sup>, the deviation from Keplerian motion due to the finite mass of the disk, $`\mathrm{\Delta }v_\varphi `$, is less than about 10 km s<sup>-1</sup> over the radius of the disk. This limits the mass of the disk to less than about $`2M\mathrm{\Delta }v_\varphi /v_\varphi `$, or about $`10^6`$ M. The density of the molecular gas must be less than $`\rho _c`$ ($`10^{10}`$ hydrogen molecules per cubic centimeter). Since $`H/R<0.0025`$, the upper limit on mass is $`10^5`$ M. If, in addition, the disk is stable against the effects of self-gravity (Toomre 1964; Binney & Tremaine 1987), then the mass of the disk must be less than $`M(H/R)`$, or about $`10^5`$ M. The mass accretion rate of a disk in steady state is given by $$\dot{M}=2\pi R\mathrm{\Sigma }v_R,$$ (9) where $`\mathrm{\Sigma }`$ is the surface density of the disk, and $`v_R`$ is radial drift velocity, which depends on the viscosity of the disk. Unfortunately, $`v_R`$ is only weakly constrained by the observations (i.e., the possible difference between the optical and radio systemic velocities) to be $`<10`$ km s<sup>-1</sup>. This provides a crude limit on the accretion rate of 100 M$`_{}`$yr<sup>-1</sup>. To further constrain the mass accretion rate requires an estimate of the viscosity of the disk. In the standard model of a thin, viscous accretion disk, as formulated by Shakura and Sunyaev (1973), $`v_R`$ can be written (see Frank, King, & Raine 1992) as $$v_R=\alpha v_\varphi \left(\frac{H}{R}\right)^2,$$ (10) where $`\alpha `$ is the dimensionless viscosity parameter ($`0\alpha 1`$). The observational limit on the ratio $`H/R`$ implies that $`v_R<0.006\alpha `$ km s<sup>-1</sup>. With the limit on mass given by the deviation from Keplerian motion, the accretion rate is less that $`10^1\alpha `$ M$`_{}`$yr<sup>-1</sup>. The infall time from the masing region is $`R/v_R`$, which from Equation (10) can be written as $$T\frac{1}{\alpha }\left(\frac{c}{c_S}\right)^2\left(\frac{R_S}{c}\right)\left(\frac{R}{R_S}\right)^{\frac{1}{2}}.$$ (11) For NGC 4258, with $`\alpha =0.1`$ and c$`{}_{S}{}^{}=2.5`$ km s<sup>-1</sup>, $`T=10^{16}`$ s, or about $`3\times 10^8`$ yrs. From the magnetic field limit and the assumption of equipartition of magnetic and thermal energy, the upper limit on $`\dot{M}`$ is also $`10^1\alpha `$ M$`_{}`$yr<sup>-1</sup>. If the maser density is the maximum allowable value, $`\rho _c`$, and the maser traces all the material in the disk, then the limit on disk thickness leads to an upper limit on $`\dot{M}`$ of $`10^2\alpha `$ M$`_{}`$yr<sup>-1</sup>. Detailed theoretical modeling can give estimates for the accretion rate. For example, a model in which the cause of the outer radial cutoff in maser emission is attributed to the transition from molecular to atomic gas leads to an estimate of $`10^4\alpha `$ M$`_{}`$yr<sup>-1</sup> (Neufeld & Maloney 1995). Gammie, Narayan, & Blandford (1999) favor an accretion rate of $`10^1\alpha `$ M$`_{}`$yr<sup>-1</sup>, based on an analysis of the continuum radiation spectrum. If the accretion rate is high, then the relative weakness of the continuum radiation may be due to the process of advection (Gammie, Narayan, & Blandford 1999). On the other hand, if the accretion rate is low, then the weak emission is due to the dearth of infalling material. In this case the gravitational power in the accretion flow may be insufficient to power the jets. 2. What is the form and origin of the warp? The form of the warp is difficult to determine precisely, because the filling factor of the masers in the disk is so small. Better measurements of the positions and directions of motion of the high-velocity features are key to defining the warp more accurately. The cause of the warp is unknown, but several suggestions have been put forward. Papaloizou, Terquem, & Lin (1998) show that the warp could be produced by a binary companion orbiting outside the maser disk. Its mass would need to be comparable to the mass of the disk ($`<10^6`$ M). Alternatively, radiation pressure from the central source will produce torques on a slightly warped disk and will cause the warp to grow (Maloney, Begelman, & Pringle 1996). Finally, it is conceivable that in the absence of other torques, the observed warp is due to the Lense-Thirring effect. A maximally rotating black hole will cause a precession of a nonaligned orbit (weak field limit) of $$\mathrm{\Omega }_{LT}=\frac{2G^2M^2}{c^3R^3},$$ (12) which can be rewritten in terms of the Schwarzschild radius as $$\mathrm{\Omega }_{LT}=\frac{1}{2}\frac{c}{R_S}\left(\frac{R_S}{R}\right)^3.$$ (13) At the inner radius of the disk ($`R/R_S=40,000`$), the precession amounts to $`3\times 10^{17}`$ s<sup>-1</sup>. This precession is very small but might be significant over the lifetime of the disk. Equation (11) suggests that the lifetime might be $`10^{16}`$ s, which would produce a differential precession of about 10 across the radius of the disk. If the axis of the disk is inclined to the axis of the black hole, then the viscosity of the disk is expected to twist the plane of the innermost part of the disk to the equatorial plane of the black hole (Bardeen & Petterson 1975; Kumar & Pringle, 1985). 3. Do the water masers trace the whole disk? The inner and outer radii of the observed masers are undoubtedly due to excitation conditions in the maser. In the vertical direction it is also possible that the masers form in a thin region within a thicker disk with atomic and ionized components. It has also been proposed that the high-velocity features are not indicative of a warped disk but trace material that has been blown off a flat disk (Kartje, Konigl, & Elitzur 1999). These proposals are difficult to test. 4. What are the physical properties of the maser spots? The spectrum of the maser has many discrete peaks which correspond to spots of maser emission on the sky. The success of measuring proper motions and accelerations of these masers suggests that they correspond to discrete condensations or density-enhanced regions in the disk. A cartoon of the blobs in a disk is shown in Figure 11. The blobs in front of the black hole may be visible because they amplify emission from the central region. No masers have been seen on the backside of the disk. On the other hand, the high-velocity masers have no continuum emission to amplify, and we may only see the ones near the midline, where the gradient in the line-of-sight velocity is small. Blobs in the rest of the disk may be radiating in directions away from the earth. The clumpiness of the medium allows us to track the individual masers. If the appearance of spots is due to blobs, or density enhancements in the disk, then the minimum gradient condition would not seem to be necessary. However, intense high-velocity maser spots may occur when blobs at the same velocity line up to form two-stage masers (Deguchi & Watson 1989). This situation forms a highly beamed maser, like the filamentary maser described in Section 2. The probability of realizing this situation is greatest along the midline, where the velocity gradient is smallest. All evidence suggests that the masers arise from discrete physical condensations. There have been several suggestions that the apparent motions of the maser spots may be due to a phase effect (e.g., a spiral density wave moving through the disk, Maoz & McKee 1998), but there is no observational evidence for this. ## 6 MASERS IN OTHER AGN At this time (early 1999), 22 masers have been detected among about 700 galaxies searched (e.g., Braatz, Wilson, & Henkel 1997). A list of these galaxies with masers is given in Table 2. The yield rate of detections is only about 3 percent. The major reason for this paucity is probably that the maser disks can only be seen if they are edge-on to the line of sight. If the typical beam angle, $`\beta `$, is 8, as in NGC 4258, then the probability of seeing a maser is about equal to $`\mathrm{sin}\beta ,`$ or 8 percent. Braatz, Wilson, & Henkel (1997) have shown that most of the known masers are associated with Seyfert II galaxies or LINERs where the accretion disks are thought to be edge-on to the earth. It is difficult to make VLBI measurements on masers weaker than about 0.5 Jy because of the need to detect the maser within the coherence time of the interferometer. Nine masers have been studied with VLBI. Four of these show strong evidence of disk structure, and two more show probable disk structure. The properties of these masers are listed in Table 3. Unfortunately, none of these masers show the simple, well-defined structure that would make them useful for precise study of the physical properties of accretion disks around black holes. ## 7 SUMMARY The measurements of the positions and velocities of the masers in the nucleus of NGC 4258 offer compelling evidence for the existence of a supermassive black hole and provide the first direct image of an accretion disk within $`10^5`$ R<sub>S</sub> of the black hole. Much more can be learned from this system. 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J., 139, 1217–1238. Trotter, A. S., Greenhill, L. J., Moran, J. M., Reid, M. J., Irwin, J. A., & Lo, K. Y. 1998, Astrophys. J., 495, 740–748. Weaver, H., Williams, D. R. W., Dieter, N. H., & Lum, W. T. 1965, Nature, 208, 29–31. QUESTIONS AND ANSWERS * What other measurements can one do with masers to try to get sensitivity to the general relativistic effects or to get information from distances closer to the mass concentration? * We have searched for masers in NGC 4258 that are closer to the central mass than 40,000 Schwarzschild radii and have found none. The reason for the absence of masers closer to the center is unclear. Masers are very sensitive to local conditions such as temperature, density, and pump power, which may depend on the geometry of the warp. Water masers arise in neutral media at temperatures below a few thousand degrees, so it is unrealistic to expect to find masers very close to the Schwarzschild radius. One possible general relativity effect might be the deflection of radiation from masers on the far side of the accretion disk (i.e., a distortion of the image of the back side of the accretion disk). However, no such masers have been found.
warning/0002/astro-ph0002527.html	ar5iv	text	# Prompt and afterglow emission from the X-ray rich GRB981226 observed with BeppoSAX ## 1 Introduction Follow-up observations of arcminute positions of Gamma-Ray Bursts (GRBs) provided by provided by the Wide Field Cameras on the BeppoSAX satellite Jager et al. (1997); Boella et al. 1997a have shown that in most cases a fading X-ray counterpart identified as the GRB X-ray afterglow is detected Frontera et al. (1998). The fading law is generally a smooth power law (e.g., Frontera et al. 1999) except in two cases: GRB970508 Piro et al. (1998), in which a late-time outburst of about 10<sup>5</sup> s duration started about 6$`\times 10^4`$s after the main event, and GRB970828 Yoshida et al. (1999), in which a peak structure of 4000 s duration appeared 1.25 $`\times 10^5`$ s after the main event. Generally the afterglow light curves extrapolated back to the time of the bursts, are in agreement with the tail of the GRB time profiles Costa et al. (1997); Piro et al. (1998); Frontera et al. (1999). This fact is considered as evidence that late afterglow emission and tail of the prompt GRB emission have the same origin Frontera et al. (1999). In the fireball model scenario (see, e.g., the recent review by Piran 1999 ), this means that both the tail of the prompt emission and the late afterglow emission can be due to an external shock produced by the interaction of a relativistically expanding fireball with the Interstellar Medium (ISM). Of the promptly localized GRBs 80% show X-ray afterglow, about 50% exhibit also optical emission, and 30% show radio emission. Radio emission is generally accompanied by optical emission, except in two cases: GRB990506 Taylor, Frail and Kulkarni (1999) from which also X-ray afterglow emission was observed with the RXTE/PCA experiment (Bacodine trigger 7549, Hurley 1999 ), and GRB981226. For this event, in spite of several attempts to find the optical counterpart Galama et al. 1999a ; Rhoads et al. (1998); Bloom et al. (1998); Schaefer et al. (1998); Castro-Tirado et al. (1998); Wozniak (1998, 1998); Lindgren et al. (1998) none was identified. The best upper limit (R$``$23 mag) to the optical flux of the GRB counterpart was reported by Lindgren et al. (1999) . Frail et al. (1999a) reported the detection of a radio counterpart (VLA 232937.2–235553). The radio source peaked to 173$`\pm 27`$ $`\mu `$Jy at 8.46 GHz at about 10 days after the burst. This time delay is tipical of all previously studied radio afterglows. However, the source declined relatively fast, following a power law decay ($`t^{\delta _R}`$) with $`\delta _R=\mathrm{\hspace{0.17em}2.0}\pm 0.4`$ Frail et al. 1999b . Two interpretations of this rapid decay have been proposed by Frail et al. (1999b) : it is either the consequence of a jetted GRB or the result of a fireball shock in an ambient medium with variable density. An optical galaxy (R = 24.85 mag), consistent with the radio transient position, has been proposed as the host galaxy of GRB981226 Frail et al. 1999b . GRB981226 was also followed up with the BeppoSAX Narrow Field Instruments. A transient X-ray source was observed, which was proposed as the X-ray afterglow of GRB 981226 Frontera et al. (1998). Here we report the properties of this afterglow emission along with those of the prompt X– and gamma–ray emission. We will discuss these properties in the light of the current models of GRBs. ## 2 Observations GRB981226 was detected with the BeppoSAX Gamma–Ray Burst Monitor (GRBM, 40–700 keV; Frontera et al. 1997, Amati et al. 1997) and WFC unit 1 (1.5–26.1 keV; Jager et al. 1997) on December 26 starting at 09:47:25 UT Di Ciolo et al. (1998). Its position was determined with an error radius of $`6^{}`$ (99% confidence level) and was centered at $`\alpha _{2000}=23^\mathrm{h}29^\mathrm{m}40^\mathrm{s}`$, $`\delta _{2000}=23^{}55^{}30^{\prime \prime }`$. A precursor is detected only in the WFC on 09:44:20 UT. About eleven hours after the burst, the Narrow Field Instruments on-board BeppoSAX were pointed at the burst location for a first target of opportunity (TOO1) observation, from December 26.8785 UT to December 28.2986 UT. A new X–ray source was detected Frontera et al. (1998) in the GRB error box with the Low Energy (LECS, 0.1-10 keV, Parmar et al. 1997 ) and Medium Energy (MECS, 2–10 keV, Boella et al. 1997b) Concentrators/Spectrometers. The net exposure time on the source for MECS and LECS was 58.4 ks and 25.6 ks, respectively. The same field was again observed about 7 days after the main event (TOO2) from 1999 January 2.7604 UT to January 3.5590 UT (total net exposure time of 25.1 ks for MECS and 8.0 ks for LECS). During this observation, the source was no more detected. Data available from GRBM include two 1 s ratemeters in two energy channels (40–700 keV and $`>`$100 keV), 128 s count spectra (40–700 keV, 225 channels) and high time resolution data (up to 0.5 ms) in the 40–700 keV energy band. WFCs (energy resolution $``$ 20% at 6 keV) were operated in normal mode with 31 channels in 1.5–26 keV and 0.5 ms time resolution Jager et al. (1997). The burst direction was offset by 7 with respect to the WFC axis. With this offset, the effective area exposed to the GRB was $``$ 420 cm<sup>2</sup> in the 40-700 keV band and 91 cm<sup>2</sup> in the 2–26 keV energy band. The background in the WFC and GRBM energy bands was fairly stable during the burst, with a slight increase with time in the $`>`$100 keV channel (about 2% in 350 s). The GRBM background was estimated by linear interpolation using the 250 s count rate data before and after the burst. The WFC spectra were extracted through the Iterative Removal Of Sources procedure (IROS <sup>1</sup><sup>1</sup>1WFC software version 105.108, e.g. Jager et al. 1997 ) which implicitly subtracts the contribution of the background and of other point sources in the field of view. The MECS source count rates and spectra for TOO1 were extracted, using the XSELECT package, from a $`3^{}`$ radius region around the source centroid, while the background level was estimated from an annulus centered on the source with inner and outer radii of 4 and 8.5 arcmin, respectively. The spectra from MECS 2 and 3 were equalized and co-added. Given the much lower exposure time, the source was much less visible in the LECS. We used for the source extraction from the LECS image the XIMAGE package Giommi et al. (1991), that permits a more refined choice of the background region. The source counts were extracted from a square box centered on the source centroid consistent with the MECS centroid position and with side of 3.5 arcmin, while the corresponding background was extracted from a square annulus of inner side 3.5 arcmin and outer side 12 arcmin, centered on the source. The uncertainties will be given as single parameter errors at 90% confidence level. ## 3 Results ### 3.1 Prompt emission Figure 1 shows the measured time profiles of GRB981226 in three energy channels after the background subtraction. In the $`\gamma `$–ray band (40–700 keV; fig. 1, middle panel), the GRB shows a single peak of about 5 s duration. Some marginal evidence of the peak appears in the high energy range ($`>`$100 keV; fig. 1, bottom panel). In the X–ray energy band (2–26 keV; fig 1, top panel), the prompt emission starts about 180 s earlier with a precursor-like event of about 50 s duration. The X–ray main event exhibits two peaks, the first of which is coincident with that detected by the GRBM. The total duration of X-ray main event is about 80 s. From the WFC images both precursor and main event are consistent with the same direction in the sky. The spectral evolution of precursor and main event was studied by subdividing the GRB time profile into five temporal slices and performing an analysis on the average spectrum of each slice (see Fig. 1). We fit the spectra with a power law (N(E)$`E^\alpha `$) and a smoothly broken power law Band et al. (1993), both photoelectrically absorbed by a neutral hydrogen column density N<sub>H</sub> Morrison and McCammon (1983). The count statistics do not permit to constrain N<sub>H</sub>, that was thus fixed to the Galactic value along the GRB direction (1.8 $`\times 10^{20}\mathrm{cm}^2`$). In Table 1 we show the results. Both laws fit the data. The Band law permits to determine the value of the peak energy E<sub>p</sub> of the logarithmic power per photon energy decade (the $`\nu F_\nu `$ spectrum). For the time slices A, C, D, where only upper limits to the gamma-ray flux were available, the upper limits of E<sub>p</sub> were derived by freezing the value of the high energy index $`\beta `$ to $`2`$. In Fig. 2 we show the $`\nu F(\nu )`$ spectra of the chosen temporal slices. The continuous line is the best fit of the Band law for slice B and of the power law for the other slices. A spectral softening is observed from slice B (onset of the main event) to the following slices. Also the precursor spectrum appears softer than that in the slice B. The spectral evolution is better shown by the behavior of the peak energy E<sub>p</sub> (see Table 1): it is low of our energy passband during the precurson, it achieves a value of about 60 keV at the onset of the main event and then goes down below our passband at the end of the burst. The $`\gamma `$–ray (40–700 keV) fluence of the burst is S$`{}_{\gamma }{}^{}=(4\pm 1)\times 10^7`$ erg cm<sup>-2</sup>, while the corresponding value found in the 2–10 keV band is S$`{}_{X}{}^{}=(5.7\pm 1.0)\times 10^7`$ erg cm<sup>-2</sup>, with a ratio S<sub>X</sub>/S<sub>γ</sub> = (1.4$`\pm `$0.4). The $`\gamma `$–ray peak flux, derived from the 1 s ratemeters, is P$`{}_{\gamma }{}^{}=\mathrm{\hspace{0.17em}0.33}`$$`\pm `$0.13 photons/cm<sup>2</sup> s corresponding to $`(6.5\pm 2.6)\times 10^8`$ erg cm<sup>-2</sup> s<sup>-1</sup>, while the corresponding 2–10 keV peak flux is P$`{}_{X}{}^{}=\mathrm{\hspace{0.17em}2.7}`$$`\pm `$0.3 photons/cm<sup>2</sup> s, corresponding to $`(1.7\pm 0.2)\times 10^8\text{ erg cm}\text{-2}\text{ s}\text{-1}`$. The peak flux in gamma-rays is the lowest observed thus far with BeppoSAX . ### 3.2 Afterglow emission During TOO1, a previously unknown X-ray source, 1SAX J2329.6-2356, was detected in the MECS, almost at the center of the GRB error box, at celestial coordinates $`\alpha _{2000}=23^\mathrm{h}29^\mathrm{m}36.1^\mathrm{s}`$, $`\delta _{2000}=23^{}55^{}58.3^{\prime \prime }`$, with an error radius of 1 arcmin Frontera et al. (1998). The MECS image obtained in the first part of TOO1 (exposure time of 26950 s), when the source was stronger, is shown in fig.Prompt and afterglow emission from the X-ray rich GRB981226 observed with BeppoSAX, left. The source is also visible in the corresponding 0.1–2 keV LECS image (6880 exposure time). The source was not detected in the TOO2, when the MECS exposure time was 25080 s (see fig. 3, right). Other count excesses, compatible with very weak celestial sources, are present in the TOO1 image. They are likely field sources, the number of which is in agreement with the log N-log S distribution of the 5-10 keV X-ray sources found by Fiore et al. (1999b) with BeppoSAX . Given its transient behavior, 1SAX J2329.6-2356 is likely the X-ray aftergow of GRB981226. From the above log N-log S distribution, the chance probability for its coincidence with a background source is about 5$`\times 10^3`$. #### 3.2.1 Spectrum We derived the average 0.1-10 keV count spectrum of 1SAX J2329.6-2356 over the first 15 hrs of TOO1, when the source was brightest. We fit it both with a power law (N(E)$`E^\alpha `$) and a blackbody, photoelectrically absorbed by the Galactic column density along the GRB direction (see Section 3.1). In the fits a normalization of a factor 0.8 was applied to the LECS spectra following the cross-calibration tests between the LECS and MECS Fiore et al. 1999a . Both laws are acceptable descriptions of the data: $`\chi _\nu ^2=\mathrm{\hspace{0.17em}0.5}`$ (5 degrees of freedom, dof) for a power-law and $`\chi _\nu ^2=\mathrm{\hspace{0.17em}1.3}`$ (5 dof) for a blackbody. However, in the case of the blackbody, for energies above 5 keV, the best fit curve is constantly below the measured bins. This suggests that the power law provides a better description of the data. The best fit power-law index is $`\alpha =1.92\pm 0.47`$. We do not find evidence of a spectral evolution of the emission: the time behavior of the C(4–10 keV)/ C(1.4–4 keV) hardness ratio is statistically consistent with a constant. No evidence of the source is found in the 15–300 keV energy range: the BeppoSAX PDS instrument Frontera et al. (1997) during TOO1 does not show any statistically significant count excess over the background level. Assuming the above power-law index, the 2 $`\sigma `$ upper limit in the 15–60 keV energy band is $`5.5\times 10^{12}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$, which is a factor $``$30 higher than the extrapolated flux from the LECS+MECS spectrum. As above mentioned, no source excess is apparent in the MECS + LECS data during TOO2. Assuming the best fit power law index obtained from the TOO1 spectrum and the Galactic column density, we derived the following 2 $`\sigma `$ upper limits to the source flux during TOO2: $`1.3\times 10^{13}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$ and $`8.2\times 10^{14}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$ in the 0.1–2 keV and 2–10 keV ranges, respectively. #### 3.2.2 Light curve The 2–10 keV MECS light curve of the afterglow in bins of 7000 s elapsed time is shown in fig. Prompt and afterglow emission from the X-ray rich GRB981226 observed with BeppoSAX (top). Its main feature is the weakness of the source in the first 7000 s bin, where its flux is below the MECS sensitivity limit (2$`\sigma `$ upper limit of $`1.5\times 10^{13}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$ in the 2-10 keV energy band). Checks were done to verify whether this non detection could be due to attitude malfunctions of the satellite. However we found that the source was correctly pointed since the beginning of the NFI TOO1. Afterwards the source flux increases by more than a factor 3 in about 10000 s ((5$`\pm 1)\times 10^{13}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$). After this peak flux the source starts fading. The light curve is also reported in the bottom panel of fig. Prompt and afterglow emission from the X-ray rich GRB981226 observed with BeppoSAX along with the WFC data points. The later fading of 1SAX J2329.6-2356 is apparent, that is well described by a power law, F(t)$``$ t, with index $`\delta =\mathrm{\hspace{0.17em}1.31}_{0.39}^{+0.44}`$. In fig. Prompt and afterglow emission from the X-ray rich GRB981226 observed with BeppoSAX, bottom we show the best fit power law along with the slope uncertainty region (90% confidence level). From the best fit light curve, the 2–10 keV afterglow fluence integrated over the time interval from the GRB end (80 s) to 10<sup>6</sup> s is $`S_a=(4.3\pm 3.8)\times 10^7\text{ erg cm}\text{-2}`$, with a ratio between X–ray afterglow fluence and the prompt $`\gamma `$–ray fluence of 1.1$`\pm `$0.9, vs. a corresponding value of (1.4$`\pm `$0.4) for the prompt emission (see Sect. 3.1). ## 4 Discussion GRB981226 shows the lowest gamma-ray peak flux among the bursts localized thus far with BeppoSAX . In the log N–log P distribution of the GRBs observed with BATSE Paciesas et al. (1999), it is located near the faint end of the distribution. The burst is marked by three peculiarities, two of which have been seen for the first time: 1. Of the bursts localized by BeppoSAX, this burst is the richest in the X-ray band: the X-ray to gamma-ray fluence ratio of $`1.4\pm 0.4`$ is the highest of all SAX bursts Frontera et al. (1999). The peak energy $`E_p60`$ keV is softer than that measured in other BeppoSAX bursts. 2. An isolated X-ray precursor occurs about 180 s before the main event (Figures 1–2). Onset of X-ray emission before the gamma-rays has been observed in other GRBs Laros et al. (1984); Murakami et al. (1991); in ’t Zand at al. (1999); Feroci et al. (1999), but only another isolated X-ray precursor, started about 25 s before the start of $`\gamma `$–rays, has been reported thus far Laros et al. (1984). 3. The X-ray afterglow light curve is peculiar. As can be seen from Figure 4 (top), the afterglow emission is undetectable during first 2 hours of the start of the NFI observations (2$`\sigma `$ upper limit of $`1.5\times 10^{13}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$ in 2-10 keV), after which it rises rapidly to (5$`\pm 1)\times 10^{13}\text{ erg cm}\text{-2}\text{ s}\text{-1}`$, then undergoes a decline in the typical power law fashion (index $`\delta =\mathrm{\hspace{0.17em}1.31}_{0.39}^{+0.44}`$). Thus at least during epoch 11–13 hrs the afterglow is undetectable. Item (2) is an uncommon feature of GRBs. Assuming the synchrotron shock model (e.g., Piran 1999 ), it implies a an initial fireball expansion Lorentz factor smaller than that found in other GRBs (e.g., Frontera et al. 1999). Item (3) is the most interesting and mysterious aspect of this burst. In all SAX observed bursts to date, the afterglow emission at X-ray wavelengths begins more or less as soon as the main gamma-ray event ends. Also in the case of GRB981226 we have evidence that the afterglow starts during the main event. Indeed the peak flux of the X-ray prompt emission ($`300`$ $`\mu `$Jy) has the same order of magnitude as that (173$`\pm 27`$ $`\mu `$Jy at 8.46 GHz) observed in the radio band $``$10 days after the burst Frail et al. 1999b . On the contrary, the X-ray flux measured by us about 13 hrs after the main event is about 4 orders of magnitude lower than the radio peak flux. Now similar values of peak fluxes in X-rays and in the radio band are expected in quasi-adiabatic cooling shocks of relativistically expanding fireballs Sari, Piran and Narayan (1998). Thus a simplest conclusion is that the afterglow emission began immediately after the burst (phase E in fig. 4), rather than 13 hrs later, and then it was strongly reduced or completely ceased, and eventually restarted 13 hours after the burst. If we accept this explanation, then we must explain the physical cause for this rather extended gap. The afterglow emission is attributed to the forward shock, i.e. shock of the ambient gas particles swept up by the advancing blast wave. A cessation of X-ray afterglow would require that there be no ambient gas (or very reduced density) and the resurgence of the afterglow at 13 hr would then require increased density – in short, a cavity surrounding the explosion. The rapid decline of the X-ray and radio afterglow require that the ambient density not be a constant but decrease with the radius Chevalier and Li (1999). Indeed, one expects such a circumburst medium around massive stars to have a complicated geometry. For example, a star which exploded as a blue supergiant would have suffered two episodes of mass loss. During the first phase, the red supergiant phase, the wind speed is low leading to a rich circumstellar medium. During the next phase, the blue supergiant phase, the fast blue supergiant wind sweeps up the circumburst medium shaped by the red supergiant wind with the net result of a low-density cavity surrounded by a dense shell at the outskirts. If one accepts this explanation then we have found an additional evidence linking GRBs to massive stars. In this scenario, we find a satisfying explanation also for item (1) and (2). The X-ray richness is because the blast wave picks up matter (baryons) as it tunnels through the envelope of the massive star. The resulting large baryon content leads to lower $`E_p`$ and hence more X-ray emission. We attribute the percursor emission to emission from shock breakout that is a natural consequence of models in which GRBs arise from the death of massive stars MacFadyen et al. (1999). The rapid decline in the radio and X-ray afterglow is because of the radial density gradient in the circumburst medium. Finally, the absence of an optical afterglow could well be due to extinction towards the GRB. If this GRB arises from death of a massive star then most likely the progenitor was in a dusty region and hence the extinction. We tend to exclude that the X-ray increase of the afterglow observed $``$13 hrs after the burst is a signature of a supernova emission. Given that the X-ray flux is similar to that measured for the X-ray counterpart of SN1998bw Pian et al. (1999), one would expect a moderate distance for GRB981226 (z$``$ 0.01) and therefore a significant detection in the optical band also in the case of a large extinction. This research is supported by the Italian Space Agency (ASI). We thank the BeppoSAX mission director R.C. Butler and the teams of the BeppoSAX Operative Control Center and Scientific Data Center for their efficient and enthusiastic support to the GRB alert program.
warning/0002/astro-ph0002009.html	ar5iv	text	# The giant radio galaxy 8C 0821+695 and its environment ## 1 Introduction Giant Radio Galaxies (GRGs) constitute an unusual class of radio sources with projected linear sizes larger than 1 Mpc<sup>1</sup><sup>1</sup>1We assume H<sub>0</sub>=75 $`\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$ and q$`{}_{0}{}^{}=0.5`$ throughout this paper.. These objects, although giant in size, do not stand out in luminosity, most of them being of low surface-brightness (Subrahmanyan et al. subra2 (1996)). Moreover, the enormous size of GRGs sometimes hampers the identification of emission from two distant lobes as part of an individual object (e.g. B2 1358+305; Parma et al. parma (1996)). Probably, these are the reasons why to date there are only about 50 known GRGs (Ishwara-Chandra & Saikia ishwara (1999); Schoenmakers et al. 2000a ). Only recently, radio surveys like the Northern VLA Sky Survey (NVSS; Condon et al. nvss (1998)) or the Westerbork Northern Sky Survey (WENSS; Rengelink et al. wenss (1997)) allow sensitive searches of GRGs with adequate angular resolution. Two basic scenarios have been envisaged to explain the outstanding sizes of GRGs. First, their lobes could be fed by extremely powerful central engines which would endow the jets with the necessary thrust to bore their long way through the ambient medium. Second, GRGs could be normal radio sources evolving in very low-density environments offering little resistance to the expansion of the jets. While the first possibility requires the existence of prominent cores and hot-spots, which are not always observed (Ishwara-Chandra & Saikia ishwara (1999)), the second possibility is supported by the high degree of polarization found in GRGs also at low frequencies (Willis & O’Dea willis (1990)) and seems to be the most plausible scenario in most cases (Mack et al. 1997b ). It seems that the expansion of a radio galaxy in a low density environment during a time long enough to allow reaching Mpc sizes, rather than higher than usual radio powers or expansion velocities, are the two basic ingredients to build up the GRG population (Schoenmakers et al. 2000b ). GRGs are located in regions hardly accessible via direct observations: they do not reside in rich galaxy clusters (Subrahmanyan et al. subra2 (1996)) and the X-ray emission around their host galaxies is usually weak (Mack et al. 1997a ). However, information about the ambient medium at very large distances from the host galaxies can still be gained through the study of their radio properties. Most GRGs probe the intergalactic medium (IGM) providing information about the density of matter outside galactic halos, adding important observational constraints to current cosmological models (Begelman & Cioffi begelman (1989); Nath nath (1995)). We discuss in this paper new observations of the GRG 8C 0821+695, a Fanaroff-Riley type II (Fanaroff & Riley fanaroff (1974)). It is optically identified with a faint $`M_R22.2`$ galaxy at a redshift of 0.538 (Lacy et al. lacy (1993)). Its high redshift, compared to other GRGs, renders 8C 0821+695 a very interesting object since it provides information on the external environment at large cosmological distances. No X-ray source coincident with the radio source is found in the Bright Point Source catalogue from the ROSAT All Sky Survey. At the distance of 8C 0821+695, one arcsecond corresponds to 4.9 kpc. ## 2 Observations and data analysis We observed 8C 0821+695 with the VLA in the framework of a complete sample of large angular size radio sources selected from the NVSS (see Lara et al. lara (1999) for a sample description), and with the 100-m Effelsberg telescope. We also incorporate for the discussion maps from the NVSS and WENSS, and maps presented by Lacy et al. (lacy (1993)). The data have been calibrated according to the scale of Baars et al. (baars (1977)). ### 2.1 WENSS and NVSS maps 8C 0821+695 appears as a straight $``$7′ long FR II radio source in the NVSS and the WENSS maps, with its main axis at a position angle (P.A.) of 11°, measured north through east. The WENSS map (Fig. 1a), at a frequency of 327 MHz and an angular resolution of 57$`\stackrel{}{.}`$7$`\times `$54$`\stackrel{}{.}`$0, presents two prominent lobes (N and S) connected by a continuous bridge of emission, although the position of the core is not evident at all. The NVSS map (Fig. 1b), made at a frequency of 1400 MHz and an angular resolution of 45″, shows a prominent central core straddling the two radio lobes. The N-lobe has higher surface brightness than the S-lobe. The mean fractional polarization ($`p_m`$) at 1400 MHz is $`p_m`$=23% in the N-lobe and $`p_m`$=26% in the S-lobe, while the core is unpolarized. The E-vectors have a similar and rather uniform orientation in the N- and S-lobes, oblique to the source main axis. ### 2.2 Effelsberg observations We observed 8C 0821+695 with the 100-m Effelsberg telescope at 10.6 GHz (Tab. 1) in order to obtain information about the morphology and polarization properties at high frequencies. The observational and data reduction procedures were those detailed by Gregorini et al. (gregorini (1992)). The maps were CLEANed as described by Klein & Mack (kleinmack (1995)). We made 40 coverages resulting –after combining– in a final noise level of 0.5 mJy/beam in total power and 0.1 mJy/beam in the polarized channels. The polarization maps were corrected for the non-Gaussian noise distribution of the polarized intensity, as described by Wardle & Kronberg (wardle (1974)). This is of particular importance in case of polarized low-brightness regions, e.g. in extended radio lobes. The 10.6 GHz map (Fig. 1c), with an angular resolution of $`69\mathrm{}`$, shows three components corresponding to the core and the two lobes. The superimposed vectors represent the electrical fields. Since at this high frequency Faraday effects are most probably negligible, a rotation by 90° immediately yields the direction of the projected magnetic field. It is oriented predominantly parallel to the source main axis. The degree of polarization at this frequency is 26% in both lobes. We also observed 8C 0821+695 with the 100-m telescope at 4850 MHz. After combination of 10 coverages we reached the confusion limit of 0.6 mJy/beam in total intensity, while the noise level was 0.1 mJy/beam in the polarization channels. Because of the large beam size (143″), this map does not reveal any additional morphological details, so we do not need to display it here. The degree of polarization derived from the Effelsberg map at 4850 MHz is 19% in the N- and S-lobes. ### 2.3 VLA observations We made continuum observations of 8C 0821+695 with the VLA in the B- and C-configurations at 1425 and 4860 MHz in dual polarization (see Tab. 1). The interferometric phases were calibrated with the nearby source J0903+679, except during 4860 MHz observations with the B-array, for which J0841+708 was used as phase calibrator. The radio sources 3C286 and 3C48 served as primary flux density calibrators. Data from the B and C arrays were combined in order to take advantage of the higher B-array resolution and of the higher C-array sensitivity to extended emission. The calibration and mapping of the data were carried out with the NRAO AIPS package. Maps at 4860 MHz had to be corrected from primary beam attenuation. In addition, correction for the non-Gaussian noise distribution in the polarized intensity map was applied. The VLA map at 1425 MHz shows a well defined compact core, and two lobes of extended emission (Fig. 2a). The N-lobe is dominated by a compact component, suggesting the existence of a faint hot-spot at the end of the jet, while a similar feature is not observed in the S-lobe. The polarized emission of 8C 0821+695 at 1425 MHz comes predominantly from this hot-spot in the N-lobe. At this position, the degree of polarization is 30%. It is 20% in the rest of the N-lobe and 25% in the S-lobe. As in the NVSS map, the electric vectors are inclined with respect to the source main axis, probably due to Faraday rotation. The VLA map of 8C 0821+695 at 4860 MHz is displayed in Fig. 2b. At this frequency, the core is the most prominent feature. The hot-spot in the N-lobe appears at the end of an elongated structure aligned with the radio axis. The N-lobe shows clear oscillations in its ridge line that are reminiscent of the “dentist drill” model which assumes a jitter of the jet-head with time (Scheuer scheuer (1982)). On the other hand, the S-lobe appears much more “relaxed” and without a strong hot-spot. We measure a total source length of 6$`\stackrel{}{.}`$95 from Fig. 2b, which corresponds to a projected linear size of 2 Mpc. The north-to-south arm ratio is 0.89. The polarized emission of 8C 0821+695 at 4860 MHz is again dominated by the N-lobe hot-spot, where we measure a mean degree of polarization of 22%, reaching 26% in the brightest peak. E-vectors are predominantly oriented perpendicularly to the source main axis. We do not detect significant polarization in the S lobe at this frequency and resolution. Our highest resolution observations (VLA B-array at 4860 GHz) provide the following coordinates for the compact core of 8C 0821+695 (J2000.0): RA = 08<sup>h</sup>25<sup>m</sup> 59$`\stackrel{s}{.}`$770, DEC = +69°20′38$`\stackrel{}{.}`$59, fully consistent with the position of the host galaxy. ## 3 Results ### 3.1 Rotation Measure and depolarization We have estimated the Rotation Measure ($`RM`$) over the two radio lobes of 8C 0821+695. To do that, we convolved the NVSS 1400 MHz and the VLA 4860 MHz polarization maps to the beam of the Effelsberg 10.6 GHz map (a circular Gaussian beam of $`69\mathrm{}`$ FWHM), and applied the AIPS task RM. Although the convolution of the VLA map with such a large beam is in general not advisable, in this case the orientation of the polarization vectors was not seriously affected. We did not use the Effelsberg 4850 MHz map because of its too low angular resolution, which would prevent us from finding any possible structure in the $`RM`$ distribution. Even so, we obtain a rather uniform $`RM`$ over the two lobes of 8C 0821+695, with an average value of $`20`$ rad m<sup>-2</sup>. The uniform and similar distributions of the electric field P.A. and of the $`RM`$ over the two lobes suggest that the contribution of the radio source to the observed $`RM`$ is negligible. Moreover, the mean galactic $`RM`$ at the position of 8C 0821+695 (galactic coordinates $`l=145\stackrel{}{.}7`$ and $`b=33\stackrel{}{.}54`$) lies between $`30`$ rad m<sup>-2</sup> and 0 rad m<sup>-2</sup> (Simard-Normandin & Kronberg simard (1980)), consistent with our observed value. We thus conclude that the observed Faraday rotation is most plausibly of galactic origin, although a small local halo contribution between $`20`$ to 10 rad m<sup>-2</sup> cannot be excluded. The fractional polarizations derived at different frequencies are essentially consistent with the lack of depolarization at low frequencies. This is clearly derived from the comparison of the low resolution data (10.6 and 5 GHz Effelsberg data and 1.4 GHz data from NVSS), taking into account that the 5 GHz data have a much larger beam and therefore are likely to suffer from beam depolarization. The higher resolution VLA data show that the two lobes are still highly polarized at 1.4 GHz, in agreement with the presence of a magnetic field ordered on the restoring beam scale. We ascribe the lack of detected polarization in the S-lobe at 5 GHz, with 9″ resolution, to the lower sensitivity of this image to extended structure. ### 3.2 The broad-band radio-spectrum of 8C 0821+695 We have compiled all available flux density measurements on 8C 0821+695 and have plotted them as a function of frequency giving values, when possible, for the whole source, the two lobes and the core separately (Tab. 2; Fig. 3). We have obtained the spectral index $`\alpha `$ (defined so that the flux density $`S\nu ^\alpha `$) of the different components from linear fits to the data: $`\alpha _N=1.15\pm 0.02`$; $`\alpha _S=1.10\pm 0.04`$. On the other hand, the core shows a flat spectrum with a mean $`\alpha =0.30\pm 0.07`$. In order to study the dependence of the spectral properties of 8C 0821+695 with frequency and distance from the core, we have made three low resolution spectral index maps using the 10.6 GHz Effelsberg, the 1400 MHz NVSS, the 327 MHz WENSS and the 151 MHz CLFST map by Lacy et al. (lacy (1993)). To construct the $`\alpha `$-maps, all total intensity maps were convolved to a circular beam of $`69\mathrm{}`$, and then were registered to the nominal position of the core. Figure 4 displays spectral index profiles along the main radio axis for the three frequency intervals, with plotted 1-$`\sigma `$ errors deduced taking into account the rms noise level of each image and the uncertainties in their flux density scales. The core shows up as a flattening of the spectrum at the center of the higher frequency profiles. In both lobes we find an overall steepening of the spectrum from the extremes towards the compact core, a behaviour typical of FR II radio sources. In addition, the steepening of the spectrum with increasing frequency is also evident from this plot. The spectrum at low frequencies is rather flat in the source extremes. This result is expected since these regions are dominated by the flat spectrum hot-spot like regions and is consistent with the injection spectral index derived in Sect. 3.4. To study in detail the spectral index distribution over 8C 0821+695 we have made a high-resolution $`\alpha `$-map registering our VLA maps at 1425 and 4860 MHz to the position of the core and performing a pixel-by-pixel evaluation of the spectral index between these two frequencies (Fig. 5). We find the same general trends observed at lower resolutions, but convolved with a now evident structure in the $`\alpha `$-distribution. The core shows a flat spectrum between these two frequencies. The lobes present similar spectral indices, although the N-lobe hints at a slightly steeper spectrum which, if real, would indicate a more efficient energy dissipation in this region. ### 3.3 Physical parameters of 8C 0821+695 We have measured the FWHM and the surface brightness at several positions on the ridge line of 8C 0821+695 by fitting Gaussians to surface brightness profiles taken perpendicular to the source main axis in the high-sensitivity 1400 MHz map by Lacy et al. (lacy (1993)). The deconvolution of the width and surface brightness was done following Appendix A in Killeen et al. (killeen (1986)). We used the standard formulae of synchrotron radiation (e.g. Miley miley (1980)) to calculate the minimum energy density $`u_{me}`$ at these positions and the corresponding magnetic field $`B_{me}`$, which is approximately the equipartition field. The total pressure is assumed to be that of equipartition between particles and magnetic field, $`P_{eq}=0.62u_{me}`$. Besides, the following assumptions were made in the calculations: i) the magnetic field is assumed to be random; ii) the energy of particles is equally stored in the form of relativistic electrons and heavy particles; iii) lower and upper frequency cutoffs were set to 10 MHz and 100 GHz, respectively; iv) the spectral index is 1.1; and v) the line-of-sight depth is equal to the deconvolved FWHM. The results are listed in Tab. 3 as a function of the filling factor of the emitting regions $`\eta `$. ### 3.4 Spectral aging Based on the maps at 151 MHz (Lacy et al. lacy (1993)), 327 MHz (WENSS), 1.4 GHz (NVSS), 4.8 GHz (VLA) and 10.4 GHz (Effelsberg), we have performed a spectral aging analysis of 8C 0821+695. We have determined the spectrum at different positions of the source averaging the flux density over selected regions with a beam equivalent area ($`69\mathrm{}`$) in order to insure independent measurements. Three selected regions were centered on the N-lobe (at the position of the hot-spot and at 711 and 364 kpc from the core) and other three on the S-lobe (at the southern extreme and at 903 and 553 kpc from the core). We show in Fig. 6 the spectra of the different selected regions and the fits to the data after the application of synchrotron loss models. Left and right panels refer to the N- and S-lobe, respectively. The spectra at the source extremes (top panels) are best fitted by the continuous injection model (CI; Pacholczyk pacholczyk (1970)), giving an injection spectral index of $`\alpha _{inj}=0.4`$ for the northern hot-spot and $`\alpha _{inj}=0.6`$ for the southern extreme of the source, in agreement with the low frequency spectral index in Fig. 4. The spectra in these regions show a break at low frequencies ($`1`$ GHz) with a moderate steepening afterwards. However, the flat spectrum of the hot-spots at low frequencies might indicate that the source here is optically thick and any spectral fit in these regions must be taken with caution. Moreover, Meisenheimer et al. (meisenheimer (1989)) find that low frequency breaks at hot-spots are more likely related to the ratio between the outflow distance and the outflow velocity of the post-shocked material after the Mach disk rather than to synchrotron aging. Thus, we will not use the information of the break frequency at the source extremes to derive synchrotron ages. Middle panels in Fig. 6 correspond to the spectra taken at 711 kpc (N-Lobe) and 903 kpc (S-Lobe) from the core, respectively. Our data do not allow us to distinguish between the Jaffe-Perola (JP; Jaffe & Perola jaffe (1973)) or the Kardashev-Pacholczyk (KP; Kardashev kardashev (1962), Pacholczyk pacholczyk (1970)) synchrotron loss models, both producing equivalent results. An injection spectral index $`\alpha _{inj}=0.7`$ has been found and kept fixed in these fits. The discrepancy between the injection spectral index in the lobes and in the hot-spots is similar to that found in Cygnus A (Carilli et al. carilli (1991)), although a physical interpretation of this fact remains unclear. The derived break frequencies are 15 GHz at 711 kpc in the N-lobe and 18 GHz at 903 kpc in the S-lobe. Bottom panels in Fig. 6 refer to the two regions nearest to the core, at 363 kpc (N-Lobe) and 553 kpc (S-Lobe). Again, fits using KP and JP models are indistinguishable. The break frequency is 1.6 GHz in the N-lobe and 2.2 GHz in the S-lobe. Due to the sharp cut-off, the flux densities at 4.8 and 10 GHz are below the noise level. In the fit we kept $`\alpha _{inj}=0.7`$ as a fixed parameter. Synchrotron ages were derived from the break frequencies in the lobes using the equation (Carilli et al. carilli (1991)): $$t_{syn}=1.61\times 10^3\frac{\sqrt{\mathrm{B}_{\mathrm{eq}}}}{\mathrm{B}_{\mathrm{eq}}^2+\mathrm{B}_{\mathrm{IC}}^2}\frac{1}{\sqrt{\nu _\mathrm{B}(1+\mathrm{z})}}.$$ (1) The synchrotron age $`t_{syn}`$ is given in Myr, the break frequency $`\nu _\mathrm{B}`$ in GHz and magnetic fields in $`\mu `$G. For the strength of the equipartition magnetic field B<sub>eq</sub> we took the corresponding values from Tab. 3, and a magnetic field equivalent to the Inverse Compton microwave background of $`\mathrm{B}_{\mathrm{IC}}=7.7\mu `$G ($`\mathrm{B}_{\mathrm{IC}}=3.25(1+\mathrm{z})^2`$). The spectral ages derived from the data are plotted in Fig. 7. The errors in the spectral ages can be derived from the uncertainties of $`\mathrm{B}_{\mathrm{eq}}`$ and $`\nu _\mathrm{B}`$; however, when $`\mathrm{B}_{\mathrm{eq}}\frac{\mathrm{B}_{\mathrm{IC}}}{\sqrt{3}}`$ (as it is approximately given in our case), the total error is dominated by the uncertainties of the break frequencies. Therefore errors in Fig. 7 depend directly on the errors in the break frequencies which we have estimated considering the 1-$`\sigma `$ region of allowance in the space of free-parameters in the fits. A weighted least-square fit yields a mean expansion velocity of 0.08 c for both lobes, which represents a measure of the rate of separation of the jet-head from the lobe material. The source spectral age, obtained by extrapolation of the age profiles up to the core, is 42 Myr. We note that the derived expansion velocity is consistent with that obtained assuming a zero age at the hot-spots, supporting the reliability of our spectral fit argument. ## 4 Discussion ### 4.1 The “conical” lobe In this section we propose a simple scenario of a jet nourishing a radio lobe with the aim of constraining the physical parameters of the jet itself and of the ambient medium surrounding the radio emission. We consider the geometry outlined in Fig. 8 for the jet and the lobe of a radio galaxy, and assume that the external medium is uniform and steady. Relativistic corrections are not taken into account here, since lobe propagation velocities are small compared with the speed of light (eg. Begelman et al. begelman2 (1984)). In general, the ratio of the advance speed of the emitting body (the jet-head velocity $`v_h`$) to the sound speed of the unperturbed external region ($`v_s`$), is related to the angle between the shock front and the velocity direction (the Mach angle $`\mathrm{\Phi }`$), through the equation: $`\mathrm{sin}\mathrm{\Phi }`$ $`=`$ $`{\displaystyle \frac{v_s}{v_h}}`$ (2) $`=`$ $`{\displaystyle \frac{1}{v_h}}\sqrt{{\displaystyle \frac{\gamma P_a}{\rho _a}}},`$ where $`\gamma `$, $`P_a`$ and $`\rho _a`$ are the adiabatic index, the pressure and the mass density of the ambient medium, respectively. The mass density may also be written $`\rho _a=\sigma n_a`$, where $`\sigma `$ is the mean atomic weight and $`n_a`$ is the number of particles per unit volume. Assuming an ideal gas, the pressure, density and temperature are related through the equation of state $`P_a=n_akT_a`$, where $`k`$ is the Boltzmann’s constant and $`T_a`$ is the temperature of the ambient medium. Equation 2 provides a relationship between $`v_h`$, the lobe geometry and the properties of the external medium. We note that $`v_h`$ as given in Eq. 2 corresponds to the advance velocity of the jet-head with respect to the external medium at the time of the observations since it depends only on local conditions measured close to the hot-spot. On the other hand, determinations based on aging arguments (Sect. 3.4) correspond to the jet-head velocity measured with respect to the lobe emitting material and averaged over the whole life of the radio source ($`v_h^{}`$). Even if all our assumptions are correct, these two estimations of the advance velocity, $`v_h`$ and $`v_h^{}`$, may be different i) if the backflow $`v_{bf}`$ of the lobe material is not negligible or ii) if the jet-head is accelerated, so that comparing mean and instantaneous velocities is meaningless. Besides, incorrect assumptions (e.g. deviations from minimum energy conditions, a break frequency distribution not related to the separation velocity, wrong estimations of the external physical parameters, etc.) may lead to differing results (see Carilli et al. carilli (1991) for a detailed discussion). In the following we will consider that all our assumptions are reasonably good approximations to the real physical situation. The absence or presence of diffuse extended tails of emission perpendicular to the source axis provide information of how important the backflow of lobe material is. In general, we will consider that the mean backflow velocity is a factor $`ϵ`$ of the mean advance velocity of the jet-head with respect to the external medium, i.e. $`v_{bf}=ϵv_h`$, so that the velocity of the head with respect to the lobe $`v_h^{}=(1+ϵ)v_h`$. Similarly, the spectral age is related to the dynamical age of the source as $`t=(1+ϵ)t_{syn}`$. Moreover, we can assume a constant jet-head deceleration and determine the initial head velocity $`v_{hi}`$ as $$v_{hi}=2(1+ϵ)^1v_h^{}v_h.$$ (3) On the other hand, the balance of the ram pressure of the ambient medium and the thrust of the jet, spread over the cross-sectional area $`A_h`$ of the bow shock at the end of the jet is (Begelman & Cioffi 1989): $$v_h\sqrt{\frac{L_j}{\sigma n_av_jA_h}},$$ (4) where $`L_j`$ is the total jet power and $`v_j`$ is the jet bulk velocity. In an uniform medium and if these two jet parameters were constant with time, a varying $`v_h`$ could be obtained through a variation in the contact surface $`A_h`$. We note that a constant $`A_h`$ and $`n_a`$ decreasing with core distance would lead to $`v_h`$ increasing with time, a situation highly unplausible (Loken et al. loken (1992)). Although a combination of both situations might occur leading to a given $`v_h`$ evolution, we will neglect density variations for simplicity. Therefore, from Eq. 4, we find a relation between $`v_h`$ and $`A_h`$ at present and initial conditions: $$A_h=A_{hi}\left(\frac{v_{hi}}{v_h}\right)^2.$$ (5) This relation translates to an opening angle ($`\theta `$), defined by the time evolution of the contact surface, given by $$\mathrm{tan}\theta =\frac{\sqrt{\frac{A_h}{\pi }}\left(1\frac{v_h}{v_{hi}}\right)}{l},$$ (6) where $`l`$ is the total jet length. ### 4.2 Application to 8C 0821+695 The N-lobe of 8C 0821+695 clearly resembles the simple geometry outlined in Fig. 8, so it seems reasonable to apply the previous analysis to this lobe. Even if we need to make several assumptions and the uncertainties of our results are high, we obtain at least an idea of the order of magnitude of the different parameters, which is the aim of these calculations. From the VLA maps at 4860 and 1425 MHz (Fig. 2), we measure an angle $`\mathrm{\Phi }=12^{}`$ at the N-lobe (see Fig. 8). However, we note that this value is a lower limit to the true Mach angle since we implicitly assume that the edge of the synchrotron emitting region defines the lobe contact discontinuity, and that this discontinuity is coincident with the bow shock produced by the jet-head advance in the external medium. Such assumption underestimates the true Mach angle, since it implies that the region of shocked external medium between the contact discontinuity and the bow shock is negligible. Moreover, we assume $`\sigma =1.4`$ amu, $`\gamma =\frac{5}{3}`$ and an upper limit of $`T_a=10^7`$ K for the external medium temperature (Barcons et al. barcons (1991)). We have accepted the minimum energy conditions and used the NVSS and WENSS maps to derive the cocoon pressure at low brightness regions in the lobes of 8C 0821+695, far away from the hot-spot; we obtain $`P_c1.5\times 10^{14}`$ Nm<sup>-2</sup>. Strictly speaking, this value constitutes an upper limit to the external pressure $`P_a`$ since the bridges in FR II radio sources are most probably overpressured with respect to the surrounding medium (e.g. Subrahmanyan & Saripalli subra1 (1993); Nath nath (1995)). However, observations (Subrahmanyan et al. subra2 (1996)) and models (Kaiser & Alexander kaiser (1997)) indicate that the lobes of FR II radio galaxies grow in a self-similar way, so that GRGs are expected to have the lowest pressures, being closer to a situation of pressure equilibrium with the outer medium than other FR II radio galaxies. Thus, we can reasonably assume for our calculations that the low-brightness regions in the lobes of 8C 0821+695 provide a good approximation of the ambient gas pressure $`P_a`$ (see also Schoenmakers et al. 2000b ). In fact, the pressure we obtain is as low as that found at the very periphery of galaxy clusters, confirming that the environment of this giant source is tenuous. Considering the upper limit temperature of $`10^7`$ K for the ambient gas, the particle density resulting from the equation of state is about $`10^2`$ m<sup>-3</sup>. This density is consistent with the limit on $`RM`$ ($`\|RM\|20`$ rad m<sup>-2</sup>) if the ambient uniform magnetic field is about 0.5 $`\mu `$G and the Faraday depth less than 1 Mpc, which are reasonable values since we do not expect a strong field to be distributed over large regions. From Eq. 2, we obtain $`v_h5\times 10^3`$ c which, according to our assumptions is an upper limit to the true jet-head velocity. Besides, from spectral aging arguments (Sect. 3.4) we obtained a mean expansion velocity of the head with respect to the lobe material $`v_h^{}=0.08`$ c. On the other hand, Lacy et al. (lacy (1993)) find evidence of a tail at the base of the N-lobe and suggest a backflow speed close to the advance speed of the head. According to that, we will assume $`ϵ1`$, so that $`v_hv_{bf}0.04`$ c. Using Eq. 3, we estimate an initial jet-head velocity $`v_{hi}0.075`$ c. The area of the head contact surface $`A_h`$ can be derived by fitting an elliptical Gaussian to the northern hot-spot at 5 GHz. We obtain a deconvolved angular diameter of $`5\mathrm{}`$ ($`25`$ kpc). We can then calculate the total power of the jet from Eq. 4, assuming a jet bulk velocity $`v_j`$c (Fernini et al. fernini (1997)). We obtain $`L_j=7.1\times 10^{37}`$ W. Alternatively, we can estimate the total power from the average minimum energy density ($`\overline{u}_{me}=2.7\times 10^{13}`$ J m<sup>-3</sup>; Tab. 3), the volume of the source (simplified to a cylinder of 2 Mpc$`\times `$200 kpc) and its dynamical age $`t=2t_{syn}=84`$ Myr, giving $`L=9\times 10^{37}`$W, fully consistent with the previous result. From Eq. 5, the angular diameter of the jet contact surface at the initial stages of the source is about 0$`\stackrel{}{.}`$33. The increase in the contact surface from this value to the measured one ($`5\mathrm{}`$) requires an opening angle $`\theta `$, defined by the hot-spot size evolution, of 0$`\stackrel{}{.}`$7 (Eq. 6). Thus, a small increase of the hot-spot size with time can naturally explain the deceleration of the jet-head and the difference in the velocity determinations from local and global conditions. Finally, the radio power of the northern lobe $`L_r`$ can be derived from the observed flux density and spectral index: $$L_r=4\pi D_L^2_{\nu _1}^{\nu _2}S(\nu )𝑑\nu ,$$ (7) where $`D_L`$ is the luminosity distance, and $`\nu _1=10`$ MHz and $`\nu _2=100`$ GHz are the assumed lower and upper frequency cutoffs. Since $`S(\nu )\nu ^\alpha `$ and $`\alpha =1.1`$, we obtain $`L_r=6.9\times 10^{35}`$ W. $`L_r`$ being about two orders of magnitude lower than $`L_j`$ would indicate that most of the jet power is devoted to the expansion of the lobe against the external medium. Results are summarized in Table 4. The S-lobe does not present such a suitable morphology for the application of the previous simple model since the head of the lobe does not have a clear cone-like appearance (see Sect. 2.3). From the observed arm-ratio, we might deduce that the external medium here could be more tenuous than in the northern lobe region. ## 5 Conclusions We present new radio observations made with the VLA and the 100-m Effelsberg radio telescope, of the GRG 8C 0821+695 at different frequencies and angular resolutions. Our data have been analyzed together with survey and literature data in order to study the details of a high-redshift GRG, and obtain information about the external medium surrounding the radio source. 8C 0821+695 is a straight 2 Mpc long FR II radio source, with a north-to-south arm-ratio of 0.89. The N-lobe contains a hot-spot, responsible for most of the source polarized emission. The S-lobe presents a more relaxed structure, without a well defined hot-spot. At high frequencies ($`\nu 1400`$ MHz) 8C 0821+695 shows a prominent compact core. We have not found any trace of the jets nourishing the lobes. The spectral index distribution over the lobes of 8C 0821+685 is typical of FR II-type radio galaxies, showing a gradual steepening from the outer ends towards the core. The mean lobe spectral index is $`\alpha =1.1`$. The core has a flat spectrum with $`\alpha =0.3`$. Using the available data, we have made a spectral-aging analysis of the source lobes, providing the dependence of the spectral break frequency and the synchrotron age with the distance from the core. We obtain a mean expansion velocity of the jet head with respect to the lobe material of $`0.08`$ c, and a spectral age of 42 Myr. This age determination might be affected by the possible existence of backflow of material in the lobe, being a lower limit to the true age. The age we derive for 8C 0821+695 is of the order of ages estimated for other GRGs (Schoenmakers et al. 2000b ). We have studied the $`RM`$ over the lobes of 8C 0821+695, obtaining a smooth and uniform distribution, which we ascribe to Faraday rotation mostly produced by the Galactic medium. Equipartition conditions have been assumed in order to derive physical parameters of the lobes at different positions, yielding magnetic fields, pressures and energy densities that are consistent with estimated values in other GRGs. Under very simple assumptions we have estimated physical parameters of the jet and the external medium of 8C 0821+695. We find that the present expansion velocity is significantly lower than the mean expansion velocity even if backflow is allowed, implying the existence of deceleration. We explain this deceleration by an increase of the cross-sectional area of the bow shock at the end of the jet with time. We note that external density estimates in the literature for other GRGs usually consider the contact surface measured at the time of the observations together with mean quantities (like the expansion velocity derived from aging arguments), resulting in external densities lower than the density we obtain. However, our results are still consistent with GRGs evolving in poor density regions. ###### Acknowledgements. We thank the referee Dr. R. Perley for helpful and constructive comments to the paper. The National Radio Astronomy Observatory is a facility of the National Science Foundation operated under cooperative agreement by Associated Universities, Inc. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. This research is supported in part by the Spanish DGICYT (PB97-1164). KHM was supported by the European Commission, TMR Programme, Research Network Contract ERBFMRXCT96-0034 “CERES”. LF, GG and MM acknowledge a partial support by the Italian Ministry for University and Research (MURST) under grant Cofin98-02-32.
warning/0002/math-ph0002047.html	ar5iv	text	# Untitled Document SUPERSINGULAR SCATTERING T. Dolinszky KFKI-RMKI, H-1525 Budapest 114, P.O.B. 49, Hungary Abstract: In ’supersingular’ scattering the potential $`g^2U_A(r)`$ involves a variable nonlinear parameter $`A`$ upon the increase of which the potential also increases beyond all limits everywhere off the origin and develops a uniquely high level of singularity in the origin. The problem of singular scattering is shown here to be solvable by iteration in terms of a smooth version of the semiclassical approach to quantum mechanics. Smoothness is achieved by working with a pair of centrifugal strengths within each channel. In both of the exponential and trigonometric regions, integral equations are set up the solutions of which when matched smoothly may recover the exact scattering wave function. The conditions for convergence of the iterations involved are derived for both fixed and increasing parameters. In getting regular scattering solutions, the proposed procedure is, in fact, supplementary to the Born series by widening its scope and extending applicability from nonsingular to singular potentials and from fixed to asymptotically increasing, linear and nonlinear, dynamical parameters. r: 24.10 Fr… 03.65 Nk I. Introduction We are going to consider scattering by singular potentials $`g^2U_A(r)`$ at fixed as well as asymptotical values of the linear and nonlinear parameters $`g^2`$ and $`A`$, respectively. The pioneering work in treating singular scattering by exact means has been done by Calogero<sup>1</sup> in terms of the phase approach to quantum particle dynamics. In particular, he calculated the high energy limit of scattering by inverse power potentials<sup>2</sup>. A further correct approach is due to Froemann and Thylwe<sup>3</sup>, who analytically proved the exactness of the first order WKB (Wentzel-Kramers-Brillouin) approximation for the case of pure inverse power potentials in the short wavelength limit. Esposito<sup>4</sup> worked in singular scattering by the wave function polydromy method to get a recursion formula between solutions with potentials of different inverse powers in the radial distance. Notice all the above approaches relied on nonlinear first order or linear second order differential equations solvable analytically, in terms of special functions, for certain simple potential shapes. Nevertheless, it is the method of integral equations that is the first rank candidate for solving two body scattering problems involving a $`general`$ shape and stage of the potential singularity. The pertinent Volterra type equations can be, perhaps, best classified in terms of the reference scattering problems implied , as follows. The most evident reference basis seems to be the unperturbed case $`g^2=0`$ taken at the physical energy $`k^2`$ and the physical orbital angular momentum $`l`$. The relevant integral equations set up for the regular wave function $`u^+(r)`$ work in the case of nonsingular potentials, exclusively (Born series). It is true that for finding irregular solutions $`u^{}(r)`$, there are analogous integral equations available, for cases of nonsingular and singular potentials alike, as shown by Newton<sup>5</sup>. However, all the above mentioned procedures can only be employed at fixed sets of the scattering parameters. Namely, the residual potential contained is there invariably given by the expression $`\mathrm{\Delta }(r)=g^2U_A(r)`$ , whence the iterated series blows up term-by-term in the limit $`g^2\mathrm{}`$ (strong coupling). The same is true of the supersingular limit. Namely, the form factor $`U_A(r)`$, besides being singular for $`r0`$ at fixed parameters, increases, by definition, beyond all limits at each fixed point off the origin for $`A\mathrm{}`$. One has therefore to seek alternative approaches outside of the Born series for solving these asymptotical cases. An interesting choice of reference basis, as proposed by Newton<sup>5</sup>, offers itself in the exceptional cases of the solution being known in closed form at zero energy and zero orbital angular momentum. The residual potential is then $`k^2l(l+1)/r^2`$. It is worthwhile to note that the iteration of the integral equation converges exclusively for singular potentials. Moreover, convergence obviously holds in both limits $`g^2\mathrm{}`$ and $`A\mathrm{}`$. Quite wide is the scope of applicability of the third type of reference problems which we are going to consider in detail. It is, in fact, the semiclassical approximation. The Langer version of this method, see e.g. Newton’s monograph<sup>5</sup>, is specified by the reference centrifugal strength $`(l+\frac{1}{2})^2`$. Owing to this very choice, the model wave function reproduces the exact regular solution of quantum mechanics (QM) near the origin in the case of nonsingular potentials. As to singular potentials, the small distance QM wave function and its WKB approximation do coincide, even independently of the selection of the reference centrifugal strength. The problem of a possible convergent iteration and that of asymptotical values taken by the potential parameters had remained thereby still open. Recently, Dolinszky<sup>6</sup> proposed a smooth version of the semiclassical approach which is free of the highly inconvenient turning point singularity inherent in both the Langer’s and the standard WKB approximation. The new procedure has been applied for developing convergent series expansion of the wave function describing scattering by singular potentials at fixed<sup>7</sup> as well as increasing<sup>8</sup> linear dynamical parameters. In the present paper an improved and generalized version of that approach will be proposed, which also includes the asymptotical case of $`increasingnonlinear`$ parameters present in the Schroedinger equation. II. Physical and reference problem The physical problem to be solved first is two-body scattering by a central singular potential say, $`g^2U_A(r)`$ at the energy $`k^2`$ in the channel of index $`l`$, with $`A`$ representing a nonlinear parameter. The potential should fulfill the following requirements: $$\begin{array}{cc}& [r^4U_A(r)]\mathrm{},[r0];U_A(r)>0,U_A^{}(r)<0,[r0];\hfill \\ & [r^3U_A(r)]0,[r\mathrm{}];U_A(r)\mathrm{},[r0,A\mathrm{}].\hfill \end{array}$$ $`(2.1)`$ In order to construct a suitable reference system for treating the problem, we introduce a triad $`(\lambda _ϵ^2,\lambda _\tau ^2,\lambda ^2)`$ of auxiliary centrifugal strengths, subject to the restrictions $$\lambda ^2(l)=\frac{1}{2}[\lambda _ϵ^2(l)+\lambda _\tau ^2(l)],\lambda _ϵ^2(l)>\lambda _\tau ^2(l).$$ $`(2.2)`$ The constants $`\lambda _ϵ^2`$ and $`\lambda _\tau ^2`$ are for the present freely chosen. The concept ’matching distance’ is defined as the positive root of the ’master equation’ $$k^2R^2g^2R^2U_A(R)\lambda ^2=0;R=R(k^2,g^2,\lambda ^2;A)>0.$$ $`(2.3)`$ The radial distance $`r`$ will be in general substituted for by the dimensionless radial coordinate $`t`$. We put $$t=\frac{r}{R},[r0].$$ $`(2.4)`$ Different regions of the space are distinguished as follows: $$\begin{array}{cc}\hfill t& <1,\mathrm{region}ϵ,(\mathrm{exponential}\mathrm{region});\hfill \\ \hfill t& >1,\mathrm{region}\tau ,(\mathrm{trigonometric}\mathrm{region});\hfill \\ \hfill t& =1,(\mathrm{matching}\mathrm{point}).\hfill \end{array}$$ $`(2.5)`$ The scattering process is governed by the radial Schroedinger equation $$\{\frac{\mathrm{d}^2}{\mathrm{d}t^2}+k^2R^2g^2R^2U_A(Rt)\frac{l(l+1)}{t^2}\}u^\pm (t)=0.$$ $`(2.6)`$ A regular-irregular pair $`u^\pm (t)`$ of its solutions behaves near the origin as $$u^\pm (t)[g^2R^2U_A(Rt)]^{\frac{1}{4}}\mathrm{exp}[\pm gR_1^tdt^{}U_A(Rt^{})^{\frac{1}{2}}],[t0].$$ $`(2.7)`$ The reference problem will be a special smooth version of the zero-order semi-classical approximation. The entire argument hinges upon a pair of wavenumber function squares for the exponential and trigonometric regions, respectively, as follows: $$\begin{array}{cc}\hfill K^2(t)& =K_ϵ^2(t)=k^2+g^2U_A(Rt)+\frac{\lambda _ϵ^2}{R^2t^2},[t<1];\hfill \\ & =K_\tau ^2(t)=k^2g^2U_A(Rt)\frac{\lambda _\tau ^2}{R^2t^2},[t>1].\hfill \end{array}$$ $`(2.8)`$ These definitions imply a number of properties concerning behavior around the matching point. In particular, one extracts from the master equation $`(2.3)`$ that $$\left[K_ϵ^2(t)\right]_{t=1}=\left[K_\tau ^2(t)\right]_{t=1};[K^2(t)]_{t=1}=\frac{\lambda _ϵ^2\lambda _\tau ^2}{2R^2}.$$ $`(2.9)`$ Overall properties also follow from Eq.(2.8). Indeed, $$\begin{array}{cc}\hfill K^2(t)& >0,\mathrm{min}\{K^2(t)\}=\left[K^2(t)\right]_{t=1},[0t];\hfill \\ \hfill K_ϵ^2(t_1)& >K_ϵ^2(t_2)\mathrm{if}t_1<t_2;K_\tau ^2(t_3)<K_\tau ^2(t_4)\mathrm{if}t_3<t_4.\hfill \end{array}$$ $`(2.10)`$ These statements issue from the restrictions (2.1) imposed on $`U_A(r)`$. Small and large distance behavior can also be thence extracted such as $$K_ϵ^2(t)\mathrm{}\mathrm{if}t0;K_\tau ^2(t)k^2\mathrm{if}t\mathrm{}.$$ $`(2.11)`$ For the derivatives of $`K^2(t)`$ we shall apply the notation: $$D_\gamma ^{(s)}(t)=\frac{1}{K_\gamma ^2(t)}\frac{\mathrm{d}^sK_\gamma ^2(t)}{\mathrm{d}t^s},[s=1,2;\gamma =ϵ,\tau ].$$ $`(2.12)`$ Hence one obtains by Eqs.(2.1)-(2.2) and (2.8) $$\left[D_\tau ^{(1)}(t)D_ϵ^{(1)}(t)\right]_{t=1}=2^4\{2\lambda ^2g^2R^3U_A^{}(R)\}>0.$$ $`(2.13)`$ A regular-irregular pair of reference wave functions $`w_\gamma ^\pm (t)`$ is defined as $$w_ϵ^\pm (t)=\eta _ϵ(t)\mathrm{exp}[\pm \omega _ϵ(1,t)],[t<1];$$ $`(2.14)`$ $$w_\tau ^\pm (t)=\eta _\tau (t)[C^\pm \mathrm{cos}\omega _\tau (1,t))+S^\pm \mathrm{sin}\omega _\tau (1,t)],[t>1],$$ $`(2.15)`$ where we introduced the ’amplitude function’ and the ’phase function’ such as $$\eta _\gamma (t)\left(\frac{k^2}{K_\gamma ^2(t)}\right)^{\frac{1}{4}},[\gamma =ϵ,\tau ];$$ $`(2.16)`$ $$\omega _\gamma (t_1,t_2)R_{t_1}^{t_2}dt^{}\|K_\gamma (t^{})\|,[\gamma =ϵ,\tau ].$$ $`(2.17)`$ As to the constants $`C^\pm `$ and $`S^\pm `$ , these parameters can be, for the moment, freely chosen. Out of them, $`C^+`$ and $`S^+`$ will be specified later by smoothness requirements. The choice $`C^{}`$ and $`S^{}`$ will, in turn, remain once for all free but the only restriction $$C^+S^{}S^+C^{}0,$$ $`(2.18)`$ which warrants independence of the functions $`w_\tau ^+(t)`$ and $`w_\tau ^{}(t))`$. Yet, there exists a sophisticated definition of the constants $`C^{}`$ and $`S^{}`$ in terms of $`C^+`$ and $`S^+`$, namely the one implied in the identity $$w_\gamma ^{}(t)w_\gamma ^+(t)\{1_1^t\frac{\mathrm{d}t^{}}{w_\gamma ^+(t^{})^2}\},[\gamma =ϵ,\tau ].$$ $`(2.19)`$ This relationship is automatically satisfied in the region $`ϵ`$ by the definition (2.14). As regards the point $`t=1\pm 0`$, Eq.(2.19) guarantees there smooth matching of $`w_\tau ^{}(t)`$ to $`w_ϵ^{}(t)`$ whenever $`w_\tau ^+(t)`$ matches there $`w_ϵ^+(t)`$ smoothly. The particular choice (2.19) is quite irrelevant from the viewpoint of the present argument. Nevertheless, it justifies the use of the same superscripts $`(^{})`$ over $`w_{\tau }^{}{}_{}{}^{}(t)`$ and $`w_{ϵ}^{}{}_{}{}^{}(t)`$. The reference wave functions $`w_\gamma ^\pm (t),[\gamma =ϵ,\tau ]`$ , of Eqs.(2.14)-(2.15) solve the pair $`(ϵ,\tau )`$ of differential equations $$\{\frac{\mathrm{d}^2}{\mathrm{d}t^2}+k^2R^2W_\gamma ^\pm (t)\frac{l(l+1)}{t^2}\}w_\gamma ^\pm (t)=0,[\gamma =ϵ,\tau ].$$ $`(2.20)`$ Notice the differential equation (2.20 is common for $`w_\gamma ^+(t)`$ and $`w_\gamma ^{}(t)`$, in the exponential and trigonometric regions alike. Namely, calculation yields for the reference potential $$W_\gamma ^\pm (t)W_\gamma (t),[\gamma =ϵ,\tau ].$$ $`(2.21)`$ with the notation $$W_\gamma (t)=g^2R^2U_A(Rt)+\mathrm{\Delta }_\gamma (t),[\gamma =ϵ,\tau ].$$ $`(2.22)`$ The expressions of the residual potential $`\mathrm{\Delta }_\gamma (t)`$ introduced here are extracted from Eqs.(2.14)-(2.15) and (2.20) as $$\mathrm{\Delta }_ϵ(t)\frac{5}{16}[D_ϵ^{(1)}(t)]^2+\frac{1}{4}D_ϵ^{(2)}(t)\frac{\lambda _ϵ^2(t)l(l+1)}{t^2},[t<1],$$ $`(2.23)`$ $$\mathrm{\Delta }_\tau (t)\frac{5}{16}[D_\tau ^{(1)}(t)]^2+\frac{1}{4}D_\tau ^{(2)}(t)\frac{\lambda _\tau ^2l(l+1)}{t^2},[t>1].$$ $`(2.24)`$ Recall that $`K^2(t)`$ of Eq.(2.8) is by Eqs. (2.9) continuous at $`t=1`$ but, due to (2.12)-(2.13), not smooth. Therefore, the residual potential develops at the matching point a jump. Indeed, $$\mathrm{\Delta }_\tau (t+0)\mathrm{\Delta }_ϵ(t0)0.$$ $`(2.25)`$ For simplicity, see expression (2.24), we shall work hence forward with the centrifugal strengths $$\lambda _ϵ^2=(l+\frac{1}{2})^2,\lambda _\tau ^2=l(l+1).$$ $`(2.26)`$ Observe that this choice is compatible with the inequality contained in the postulates (2.2). III. A pair of convergent expansions A comparison of the Schroedinger equations (2.6), set up for the exact wave function $`u^+(r)`$, and Eq.(2.20), solved by the semiclassical wave function $`w^+(r)`$, suggests construction of a pair of integral equations, $$v_ϵ^+(t)=w_ϵ^+(t)+_0^tdt^{}\mathrm{\Delta }_ϵ(t^{})G_ϵ^+(t,t^{})v_ϵ^+(t^{}),[t<1],$$ $`(3.1)`$ $$v_\tau ^+(t)=w_\tau ^+(t)+_1^tdt^{}\mathrm{\Delta }_\tau (t^{})G_\tau ^+(t,t^{})v_\tau ^+(t^{}),[t>1].$$ $`(3.2)`$ The solutions $`v_ϵ^+(t)`$ and $`v_\tau ^+(t)`$, if exist, are solutions of the exact Schroedinger equation within the respective regions. In particular, $`v_ϵ^+(t)`$ is uniquely defined by Eq.(3.1) and furnishes a regular solution of the differential equation (2.6) in the exponential region. The solution $`v_\tau ^+(t)`$ of Eq.(3.2), in turn, while solving the Schroedinger equation in the trigonometric region, still involves two free constants, $`C^+`$ and $`S^+`$, as noticed following Eq.(2.15). Smoothness requirement for the overall regular solution $`v^+(t)`$ at $`t=1`$ is just sufficient to unequivocally specify these coefficients. As to the notation in (3.1)-(3.2), the residual potentials $`\mathrm{\Delta }_\gamma (t)`$ have been defined by (2.22)-(2.24). The resolvents involved in the integral equations are formally given as $$G_\gamma ^+(t,t^{})=\frac{1}{d_\gamma ^+}[w_\gamma ^+(t)w_\gamma ^{}(t^{})w_\gamma ^{}(t)w_\gamma ^+(t^{})],[\gamma =ϵ,\tau ],$$ $`(3.3)`$ where the Wronskians contained are by Eq.(2.21) constant and read in general $$d_\gamma ^+=\mathrm{W}_\gamma \{w_\gamma ^+(t);w_\gamma ^{}(t)\}=\mathrm{const}.,[\gamma =ϵ,\tau ].$$ $`(3.4a)`$ In particular, one obtains after some calculations $$d_ϵ^+=2kR,d_\tau ^+=kR(C^+S^{}C^{}S^+).$$ $`(3.4b)`$ The general expressions (3.3) can be recast in terms of the local wave numbers as $$G_ϵ^+(t,t^{})=2\frac{\mathrm{sinh}[\omega _ϵ(t,t^{})]}{R[K_ϵ^2(t)K_ϵ^2(t^{})]^{\frac{1}{4}}},[0t^{}t<1],$$ $`(3.5)`$ $$G_\tau ^+(t,t^{})=\frac{\mathrm{sin}[\omega _\tau (t,t^{})]}{R[K_\tau ^2(t)K_\tau ^2(t^{})]^{\frac{1}{4}}},[1t^{}t].$$ $`(3.6)`$ Irrelevance of any special choice of the basis $`w_\tau ^\pm (t)`$ for inclusion in the formula (3.3) manifests itself in the absence of the constants $`C^\pm ,S^\pm `$ from the last expression. The solution of the integral equations (3.1)-(3.2) rests upon the recursion schemes $$w_{ϵn}^+(t)_0^t\mathrm{d}t^{}\mathrm{\Delta }_ϵ(t^{})G_ϵ^+(t,t^{})w_{ϵn1}^+(t^{}),[n=1,2,3..];w_{ϵ0}^+(t)w_ϵ^+(t),$$ $`(3.7)`$ $$w_{\tau m}^+(t)_1^t\mathrm{d}t^{}\mathrm{\Delta }_\tau (t^{})G_\tau ^+(t,t^{})w_{\tau m1}^+(t^{}),[m=1,2..];w_{\tau 0}^+(t)w_\tau ^+(t).$$ $`(3.8)`$ A necessary condition for getting the solution sought for by iteration is convergence of each of the infinite series $$v_ϵ^+(t)=\underset{n=0}{\overset{\mathrm{}}{}}w_{ϵn}^+(t),[t<1];$$ $`(3.9a)`$ $$v_\tau ^+(t)=\underset{m=0}{\overset{\mathrm{}}{}}w_{\tau m}^+(t),[t>1].$$ $`(3.9b)`$ The overall solution $`v^+(t)=\{v_ϵ^+(t);v_\tau ^+(t)\}`$ should be throughout smooth. This requirement is realized off the matching point spontaneously. At the matching point, it is the zero order term of the series (3.9b) that exclusively contributes to both the solution $`v_\tau ^+(t)`$ and its first derivative. If the series (3.9a) is in the region $`ϵ`$ convergent, the smoothness postulate for $`t=1`$ can be recast in a simple form as $$[w_\tau ^+(t)]_{t=1}=[v_ϵ^+(t)]_{t=1};[w_\tau ^+(t)^{}]_{t=1}=[v_ϵ^+(t)^{}]_{t=1}.$$ $`(3.10)`$ These conditions simultaneously fix the trigonometric constants $`C^+`$ and $`S^+`$ thus completing the solution in the region $`\tau `$. The convergence proof for the above series becomes more transparent by using the following notation in both regions $`\gamma =ϵ,\tau `$ $$q_{\gamma s}^+(t)\frac{w_{\gamma s}^+(t)}{w_{\gamma 0}^+(t)};q_{\gamma 0}^+(t)=1,$$ $`(3.11)`$ $$p_\gamma (t)\frac{\mathrm{\Delta }_\gamma (t)}{RK_\gamma (t)},$$ $`(3.12)`$ $$P_\gamma (t_1,t_2)_{t_1}^{t_2}dt^{}\|p_\gamma (t^{})\|.$$ $`(3.13)`$ . In the exponential region, the formula (3.7) can thus be rewritten as $$q_{ϵn}^+(t)=_0^tdt^{}p_ϵ(t^{})\{1\mathrm{exp}[2\omega _ϵ(t,t^{})]\}q_{ϵn1}^+(t^{}).$$ $`(3.14)`$ Hence we get by Eq.(2.17) the inequality $$\|q_{ϵn}^+(t)\|_0^tdt^{}\|p_ϵ(t^{})q_{ϵn1}^+(t^{})\|,[t<1].$$ $`(3.15)`$ Iteration yields then by Eq.(3.13) $$\|q_{ϵn}(t)\|\frac{1}{n!}[P_ϵ(0,t)]^n,[t<1],$$ $`(3.16)`$ whence one extracts by analysis $$\underset{n=0}{\overset{\mathrm{}}{}}\|w_{ϵn}^+(t)\|<\|w_ϵ(t)\|\mathrm{exp}[P_ϵ(0,t)],[t<1].$$ $`(3.17)`$ The series $`v_ϵ^+(t)`$ of Eq.(3.9a) is thus absolutely convergent if and only if the integral $`P_ϵ(0,t)`$ exists and is bounded in $`t=(0,1)`$. As to the trigonometric region, the resolvent formula (3.6) combines with the notation (3.12) to an equivalent form of the recursion relationship (3.8). Accordingly, we get $$w_{\tau m}^+(t)=_1^tdt^{}p_\tau (t^{})\mathrm{sin}\omega _\tau (t,t^{})\frac{K_\tau ^{\frac{1}{2}}(t^{})}{K_\tau ^{\frac{1}{2}}(t)}w_{\tau m1}^+(t^{}),[t>1].$$ $`(3.18)`$ The monotonicity relationship (2.10) implies then the inequality $$\|w_{\tau m}^+(t)\|<_1^tdt^{}\|p_\tau (t^{})\|\|w_{\tau m1}^+(t^{})\|,[t>1,m=1,2,3\mathrm{}].$$ $`(3.19)`$ Recall now the definition (2.15) and conclude by the zero order, $`m=0`$, identity in the relationships (3.8) that $$\|w_{\tau 0}^+(t)\|<(kR)^{\frac{1}{2}}(\|C^+\|+\|S^+\|),[t1].$$ $`(3.20)`$ Iteration furnishes then by the inequalities (3.19)-(3.20) $$\|w_{\tau m}^+(t)\|<(kR)^{\frac{1}{2}}(\|C^+\|+\|S^+\|)\frac{1}{m!}[P_\tau (1,t)]^m,[t>1;m=1,2,3\mathrm{}].$$ $`(3.21)`$ Summation over $`m`$ yields thus $$\underset{m=0}{\overset{\mathrm{}}{}}\|w_{\tau m}^+(t)\|<(kR)^{\frac{1}{2}}(\|C^+\|+\|S^+\|)\mathrm{exp}P_\tau (1,t),[t>1].$$ $`(3.22)`$ So the series $`v_\tau ^+(t)`$ of the definition (3.9b) is absolutely convergent whenever the integral $`P_\tau (1,t)`$ of Eq.(3.13) does exist. Suppose the regional existence conditions $$P_ϵ(0,t)<\mathrm{},[t<1],P_\tau (1,t)<\mathrm{},[t>1],$$ $`(3.23)`$ and, in addition, the smoothness postulate (3.10) are satisfied. If so, then solutions of the integral equations (3.1)-(3.2) together recover the regular solution of the differential equation (2.6). In particular, $$u^+(t)=v_ϵ^+(t),[t<1];u^+(t)=v_\tau ^+(t),[t>1].$$ $`(3.24)`$ Virtually, one always works , instead of the infinite expansions (3.9a)-(3.9b), with cut-off series such as $$v_ϵ^{+(N)}(t)=\underset{n=0}{\overset{N}{}}w_{ϵn}^+(t),[t<1];$$ $`(3.25)`$ $$v_\tau ^{+(NM)}(t)=\underset{m=0}{\overset{M}{}}w_{\tau m}^{+(N)}(t),[t>1].$$ $`(3.26)`$ Smooth matching of these functions at $`t=1`$ in terms of the relevant analog of Eqs. (3.10) fixes the trigonometric coefficients $`C^+`$ and $`S^+`$ in terms of the cut-off ’length’ $`N`$ in the exponential region but independently of its pair $`M`$ in the trigonometric region. These constants therefore should carry the superscripts $`(N)`$ only. The postulates (3.10), rewritten for the cut-off approach, read $$2^{\frac{3}{4}}(kR)^{\frac{1}{2}}C^{+(N)}=[v_ϵ^{+(N)}(t)]_{t=1},$$ $`(3.27)`$ $$2^{\frac{5}{4}}[D_\tau ^{(1)}(t)]_{t=1}(kR)^{\frac{1}{2}}C^{+(N)}+2^{\frac{3}{4}}S^{+(N)}=[v_ϵ^{+(N)}(t)^{}]_{t=1}.$$ $`(3.28)`$ Notice the cut-off approach must not be used unless the conditions of convergence are fulfilled in both regions $`ϵ`$ and $`\tau `$. IV. The supersingularity limit The term ’supersingularity’ implies in our terminology a singular potential which is subject to the requirements (2.1). This involves, among others, that it contains a nonlinear parameter $`A`$ which may increase beyond all limits. The relevant asymptotical form of the master equation (2.3) reads $$k^2\frac{g^2}{r_0^2}U_A(R_A)\frac{\lambda ^2}{R_A^2}0,[A\mathrm{}].$$ $`(4.1)`$ Bounded values such as $`R_A0`$ and $`R_A\mathrm{const}.`$ are by the properties (2.1) obviously excluded from the large-$`A`$ solutions of the Eq.(4.1). One is thus left with the only possibility that $`R_A\mathrm{}`$ for $`A\mathrm{}`$. The master equation itself reduces therefore in the supersingularity limit to $$U_A(R_A)\frac{k^2r_0^2}{g^2},[A\mathrm{}].$$ $`(4.2)`$ It is perhaps worth recalling that $`U_A(r)0\mathrm{if}A=\mathrm{fixed},r\mathrm{}`$ while $`U_A(r)\mathrm{}\mathrm{if}r=\mathrm{fixed},A\mathrm{}`$, as implied in the set of postulates (2.1). In Eq. (4.2), in turn, the singularity parameter $`A`$ and the matching radius $`R_A`$ vary simultaneously and both increase to infinity so as to keep the left hand side of the equation constant. Four classes of strongly singular potentials will be introduced with each potential being specified by a variable core parameter $`A`$ and a fixed tail parameter $`B`$. The dimensionless form factor is in each case a product of an, for $`r0`$, exponentially or powerlaw increasing core factor and an, for $`r\mathrm{}`$, exponentially or powerlike decreasing tail factor. Concerning the subscripts, we use an obvious notation when writing $$\begin{array}{cc}\hfill U_{\alpha \beta }(r)& =\mathrm{exp}(\frac{\alpha \rho _1}{r}\frac{\beta r}{\rho _2}),\hfill \\ \hfill U_{a\beta }(r)& =(1+\frac{r_1}{r})^a\mathrm{exp}(\frac{\beta r}{\rho _2}),\hfill \\ \hfill U_{\alpha b}(r)& =\mathrm{exp}(\frac{\alpha \rho _1}{r})(\frac{r_2}{r_2+r})^b,\hfill \\ \hfill U_{ab}(r)& =(1+\frac{r_1}{r})^a(\frac{r_2}{r_2+r})^b,\hfill \\ \hfill [0r,0<\alpha ,\beta ,& r_1,r_2,\rho _1,\rho _2;a>4;b>3].\hfill \end{array}$$ $`(4.3)`$ In the last section, we established general criteria for the convergence of the expansions (3.9a)-(3.9b). In the present section, fulfillment of those requirements will be checked for increasing nonlinear parameters. In the supersingularity limit $`A\mathrm{}`$,the asymptotical form of the master equation (4.2) is explicitly solvable in cases of exponential tail for both types of core singularity. In fact, one concludes from the definitions (4.3) that $$\begin{array}{cc}\hfill R_{\alpha \beta }& (\rho _1\rho _2)^{\frac{1}{2}}(\frac{\alpha }{\beta })^{1/2},[\alpha \mathrm{}];\hfill \\ \hfill R_{a\beta }& (r_1\rho _2)^{\frac{1}{2}}(\frac{a}{\beta })^{\frac{1}{2}},[a\mathrm{}].\hfill \end{array}$$ $`(4.4)`$ In the powerlaw tail cases, in turn, only implicit expressions can be extracted from the formulas (4.2)-(4.3). Indeed, one gets $$\begin{array}{cc}\hfill R_{\alpha b}& r_2\mathrm{exp}(\frac{\alpha \rho _1}{bR_{\alpha b}}),[\alpha \mathrm{}],\hfill \\ \hfill R_{ab}& r_2\mathrm{exp}(\frac{ar_1}{bR_{ab}}),[a\mathrm{}].\hfill \end{array}$$ $`(4.5)`$ The matching distance increases in all of our examples slower than the respective variable parameter $`A=\alpha \mathrm{or}a`$. Indeed, one finds by the relationships (4.4)-(4.5) that $$\frac{R_{AB}}{A}0,[A\mathrm{};A=\alpha ,a;B=\beta ,b].$$ $`(4.6)`$ The notation $`A\mathrm{}`$ and $`R_{AB}\mathrm{}`$ compete in representing the supersingular limit by a single symbol. We prefer the use of the latter alternative where possible. Accordingly, we shall throughout eliminate from the formulas the parameter $`A`$ in terms of $`R_{AB}`$. It is also obvious by analysis that the present scattering formalism should, in the supersingularity limit, become, mutatis mutandis, identical for exponentially and powerlaw increasing cores. The equivalence is realized at the following correspondence of potential parameters: $$r_1a\rho _1\alpha ,[a,\alpha \mathrm{},B=\beta ,b].$$ $`(4.7)`$ In the limit considered, one has thus to treat, out of the four cases in (4.3), only two essentially different ones, namely ($`A,\beta `$) and ($`A,b)`$. We now rewrite the wave number squares of Eqs.(2.8), in terms of the variable $`t`$ of the definition (2.4), for $`R_{AB}\mathrm{}`$. In the cases ($`A,\beta `$), this transformation is, owing to the explicit relationships (4.4), straightforward. Indeed, one obtains by the definitions (4.3) in the exponential region $$\begin{array}{cc}& K_{ϵA\beta }^2(t)k^2\{\frac{g^2}{k^2r_0^2}\mathrm{exp}[\frac{\beta R_{A\beta }}{\rho _2}(\frac{1}{t}t)]+\frac{\lambda _ϵ^2}{k^2R_{A\beta }^2t^2}1\},\hfill \\ & [t<1,R_{A\beta }\mathrm{},A=\alpha ,a].\hfill \end{array}$$ $`(4.8)`$ A little bit more complicated is the incorporation of formula (4.5) into Eq.(2.8) in power tail cases for which one gets still in the region $`ϵ`$ $$\begin{array}{cc}& K_{ϵAb}^2(t)k^2\{[\frac{k^2r_0^2}{g^2}(\frac{R_{Ab}}{r_2})^b]^{\frac{1}{t}1}\frac{1}{t^b}1+\frac{\lambda _ϵ^2}{k^2R_{Ab}^2t^2}\},\hfill \\ & [t<1,R_{Ab}\mathrm{},A=\alpha ,a].\hfill \end{array}$$ $`(4.9)`$ In the trigonometric region, distinction should be made between S-wave and higher partial waves. Indeed, one extracts from the definitions (2.8) and (4.3) for the potential classes $`A=\alpha `$ and $`a`$ equally, that $$K_{\tau A\beta }^2(t)k^2\{1\frac{g^2}{k^2r_0^2}\mathrm{exp}[\frac{\beta R_{A\beta }}{\rho _2}(t\frac{1}{t})]\},[t>1,l=0,R_{A\beta }\mathrm{}],$$ $`(4.10)`$ $$K_{\tau Ab}^2(t)k^2\{1[\frac{g^2}{k^2r_0^2}(\frac{r_2}{R_{Ab}})^b]^{1\frac{1}{t}}\frac{1}{t^b}\},[t>1,l=0,R_{Ab}\mathrm{}],$$ $`(4.11)`$ $$K_{\tau AB}^2(t)k^2\{1\frac{\lambda _\tau ^2}{k^2R_{AB}^2t^2}\},[t>1,l>0,R_{AB}\mathrm{}].$$ $`(4.12)`$ Notice for the highest values of the parameter $`A`$ the wave number function becomes independent of the potential: $$K_{\tau AB}^2(t)k^2,[t>1,l0,R_{AB}\mathrm{}].$$ $`(4.13)`$ The quantities $`D_\gamma ^{(s)}(t)`$ of the definition (2.12 will be calculated below from the set of supersingularity expressions (4.8)-(4.12). In the exponential region one finds for both cases $`A=\alpha ,a`$ in the limit $`R_{AB}\mathrm{}`$ $$D_{ϵA\beta }^{(1)}(t)\frac{\beta R_{A\beta }}{\rho _2}(\frac{1}{t^2}+1),D_{ϵA\beta }^{(2)}(t)(\frac{\beta R_{A\beta }}{\rho _2})^2(\frac{1}{t^2}+1)^2,[t>1,l0],$$ $`(4.14)`$ $$D_{ϵAb}^{(1)}(t)\frac{b}{t^2}\mathrm{ln}(\frac{R_{Ab}}{r_2}),D_{ϵAb}^{(2)}(t)\frac{b^2}{t^4}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2,[t>1,l0].$$ $`(4.15)`$ As regards the trigonometric region, the formulae are again sensitive to the orbital angular momentum. Indeed, for the potential classes $`A=\alpha ,a`$ alike, one gets $$D_{\tau A\beta }^{(1)}(t)\frac{\beta R_{A\beta }}{\rho _2}(\frac{1}{t^2}+1),D_{\tau A\beta }^{(2)}(t)(\frac{\beta R_{A\beta }}{\rho _2})^2(\frac{1}{t^2}+1)^2,[t>1,l=0],$$ $`(4.16)`$ $$D_{\tau Ab}^{(1)}(t)\frac{b}{t^2}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2,D_{\tau Ab}^{(2)}(t)\frac{b^2}{t^4}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2,[t>1,l=0],$$ $`(4.17)`$ $$D_{\tau AB}^{(1)}(t)\frac{2\lambda _\tau ^2}{k^2R_{AB}^2t^3},D_{\tau AB}^{(2)}(t)\frac{6\lambda _\tau ^2}{R_{AB}^2t^4},[t>1,l>0].$$ $`(4.18)`$ The residual potentials (2.23)-(2.24) are quadratic and linear expressions of the quantities (4.14)-(4.18). They read in the exponential region, for exponential and powerlaw cores in like manner, $$\mathrm{\Delta }_{ϵA\beta }(t)\frac{1}{16}(\frac{\beta R_{A\beta }}{\rho _2})^2(\frac{1}{t^2}+1)^2,[t<1,R_{A\beta }\mathrm{}],$$ $`(4.19)`$ $$\mathrm{\Delta }_{ϵAb}(t)\frac{1}{16}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2\frac{b^2}{t^4},[t<1,R_{Ab}\mathrm{}].$$ $`(4.20)`$ The analogous formulae of the trigonometric region are sensitive to the orbital angular momentum as $$\mathrm{\Delta }_{\tau A\beta }(t)\frac{9}{16}(\frac{\beta R_{A\beta }}{\rho _2})^2(\frac{1}{t^2}+1)^2,[t>1,l=0,R_{A\beta }\mathrm{}],$$ $`(4.21)`$ $$\mathrm{\Delta }_{\tau Ab}(t)\frac{9}{16}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2\frac{b^2}{t^4},[t>1,l=0,R_{Ab}\mathrm{}],$$ $`(4.22)`$ $$\mathrm{\Delta }_{\tau AB}(t)\frac{3\lambda _\tau ^2}{2k^2R_{AB}^2t^4},[t>1,l>0,R_{AB}\mathrm{},B=\beta ,b].$$ $`(4.23)`$ The next step towards checking expansions (3.9) for the realization of convergence criteria is calculation of the quantities $`p_\gamma (t)`$ introduced by Eq.(3.12). The asymptotical expressions (4.8)-(4.13) combine with the ones of (4.19)-(4.23) to yield, first for the exponential region, $$p_{ϵA\beta }(t)\frac{1}{16}\frac{\beta ^2R_{A\beta }r_0}{g\rho _2^2}(\frac{1}{t^2}+1)^2\mathrm{exp}[\frac{\beta R_{A\beta }}{\rho _2}(\frac{1}{t}t)],[t<1,R_{A\beta }\mathrm{}],$$ $`(4.24)`$ $$p_{ϵAb}(t)\frac{1}{16}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2\frac{b^2}{kR_{Ab}}t^{\frac{b}{2}2}[\frac{g^2}{k^2r_0^2}(\frac{r_2}{R_{Ab}})^b]^{\frac{1}{t}1},[t<1,R_{Ab}\mathrm{}],$$ $`(4.25)`$ as well as in the trigonometric region $$p_{\tau A\beta }(t)\frac{9}{16}(\frac{\beta R_{A\beta }}{\rho _2})^2\frac{1}{kR_{A\beta }}(\frac{1}{t^2}+1)^2,[t>1,l=0,R_{A\beta }\mathrm{}],$$ $`(4.26)`$ $$p_{\tau Ab}(t)\frac{9}{16}[\mathrm{ln}(\frac{R_{Ab}}{r_2})]^2\frac{b^2}{t^4}\frac{1}{kR_{Ab}},[t>1,l=0,R_{Ab}\mathrm{}],$$ $`(4.27)`$ $$p_{\tau AB}(t)\frac{3\lambda _\tau ^2}{2k^2R_{AB}^2t^4}\frac{1}{kR_{AB}},[t>1,l>0,R_{AB}\mathrm{}].$$ $`(4.28)`$ The supersingularity forms of $`P_ϵ(0,t)`$ and $`P_\tau (1,t)`$ are obtained by first integrating the expressions (4.24)-(4.28) over the relevant intervals and subsequently going over to the limit $`R_{AB}\mathrm{}`$. Within the exponential region, the limits $`t0`$ and $`R_{AB}\mathrm{}`$ mutually strengthen the rate of vanishing. The factor $`(\frac{r_2}{R_{Ab}})^{\frac{1}{t}}`$ in the expression (4.25) vanishes at $`R_{Ab}>r_2`$ for $`t0`$ . The functions $`p_{ϵAB}(t)`$ are integrable near the point $`t=0`$ even at finite values of $`R_{AB}`$. Moreover, its integral vanishes in the limit $`R_{AB}\mathrm{}`$. Thus $$P_{ϵAB}(0,t)0,[t<1;A=\alpha ,a;R_{AB}\mathrm{}].$$ $`(4.29)`$ By virtue of the inequality (3.17), the condition for the absolute convergence of the series (3.9a) is thus fulfilled in the supersingularity limit along the exponential region for each of the potentials $`U_{AB}(r)`$ of the set (4.3). As regards the trigonometric region, there is a competition between the potential and the centrifugal term and this is the point that governs convergence.In the absence of the latter, one extracts from Eq.(4.26) for both cases $`A=\alpha ,a`$ that $$P_{\tau A\beta }(1,t)\mathrm{},[t>1;l=0;R_{A\beta }\mathrm{}].$$ $`(4.30)`$ The convergence of the series (3.9b) is thereby frustrated for the S-wave whenever the physical potential decreases exponentially. Not so for the higher partial waves or the cases of powerlaw tails. Equations (4.27)-(4.28) imply namely that $$P_{\tau Ab}(1,t)0,[t>1;l=0;R_{Ab}\mathrm{}],$$ $`(4.31)`$ $$P_{\tau AB}(1,t)0,[t>1;l>0;R_{AB}\mathrm{}].$$ $`(4.32)`$ The asymptotical relationships (4.29), (4.31) and (4.32) ensure, by the the inequalities (3.17) and (3.22), for the respective potential classes and partial waves, fulfillment of the convergence conditions (3.23), simultaneously at fixed and increasing values of the nonlinear parameter $`A`$, involved in the scattering potential $`U_{AB}(r)`$. One can thus write in the limit $`R\mathrm{}`$ that $$u^+(t)v_ϵ^+(t),[t<1;B=\beta ,b;l0],$$ $`(4.33)`$ $$u^+(t)v_\tau ^+(t),[t>1;B=b;l=0],$$ $`(4.34)`$ $$u^+(t)v_\tau ^+(t),[t>1;B=b,\beta ;l>0].$$ $`(4.35)`$ V. Asymptotical exactness In the last section we studied the question whether a semiclassical expansion that is convergent at a fixed set of dynamical parameters preserves this property invariably in the supersingular limit. A further point we are left with to clear is how the structure of the series changes upon going with the nonlinear parameter to infinity. A characteristic quantity of the argument for finding the answer is the ratio of two neighboring general terms. Recall first the definition (3.11) of $`q_{\gamma s}^+(t)`$ along with the relevant recursion relationships (3.14) and (3.18). For simplicity, in the present section we are going to suppress in the formulas the matching distance $`R`$. Remember it is a functional of the scattering potential $`g^2U_A(Rt)`$ and is, in fact, invisibly present in each quantity and expression below. In the exponential region, the exact formula (3.14) reduces in the supersingular limit, owing to the definition (2.17), to $$q_{ϵn}^+(t)_0^tdt^{}p_ϵ(t^{})q_{ϵn1}^+(t^{}),[t<1,R\mathrm{}].$$ $`(5.1)`$ Iteration yields then in the limit considered a familiar expression $$q_{ϵn}^+(t)\frac{1}{n!}[P_ϵ(0,t)]^n,[t<1,R\mathrm{}].$$ $`(5.2).`$ The scattering wave function reads thus by the definitions (3.9a) and (3.13) in the limit discussed $$v_ϵ^+(t)w_{ϵ0}^+(t)\mathrm{exp}P_ϵ(0,t),[t<1,R\mathrm{}].$$ $`(5.3)`$ For treating the trigonometric region , it is useful to introduce higher order trigonometric ’coefficients’, which are, as a matter of fact, no more constant. The definition is meant to hold for both fixed or variable dynamical parameters and reads $$w_{\tau m}^+(t)=\eta _\tau (t)[C_m^+(t)\mathrm{cos}\omega _\tau (1,t)+S_m^+(t)\mathrm{sin}\omega _\tau (1,t)],[t>1,m=0,1..].$$ $`(5.4)`$ Recognize the identity (5.4) reproduces at $`m=0`$ the definition (2.15), on account of which one finds that $`C_0^+=C^+`$ and $`S_0^+=S^+`$ . Combination of this identity with the recursion formula (3.18) furnishes in our limit, after separating the sine- and cosine- contributions, the following system of equations $$C_m^+(t)_1^tdt^{}p_\tau (t^{})\mathrm{sin}[kR(t^{}1)]w_{\tau m1}^+(t^{}),[t>1,R\mathrm{}],$$ $`(5.5)`$ $$S_m^+(t)_1^tdt^{}p_\tau (t^{})\mathrm{cos}[kR(t^{}1)]w_{\tau m1}^+(t^{}),[t>1,R\mathrm{}].$$ $`(5.6)`$ Insertion of the definition (5.4) into the right hand sides of the last two formulas yields by analysis, owing to the infinitely rapid oscillations in the integrands, a pair of coupled systems of integral equations such as $$C_m^+(t)\frac{1}{2}_1^tdt^{}p_\tau (t^{})S_{m1}^+(t^{}),[t>1,R\mathrm{}],$$ $`(5.7)`$ $$S_m^+(t)\frac{1}{2}_1^tdt^{}p_\tau (t^{})C_{m1}^+(t^{}),[t>1,R\mathrm{}].$$ $`(5.8)`$ At the end of the iteration, at $`m=0`$, one encounters the constant coefficients $`C_0^+,S_0^+`$ . The system of integral equations becomes thereby explicitly solvable. The solution reads, in terms of the notation $$T_m(t)\frac{1}{m!}[\frac{1}{2}P_\tau (1,t)]^m,[t>1,m=0,1,2,..],$$ $`(5.9)`$ as follows: $$C_{4\mu }^+(t)+T_{4\mu }(t)C_0^+,S_{4\mu }^+(t)+T_{4\mu }(t)S_0^+,$$ $`(5.10)`$ $$C_{4\mu +1}^+(t)+T_{4\mu +1}(t)S_0^+,S_{4\mu +1}^+(t)T_{4\mu +1}(t)C_0^+,$$ $`(5.11)`$ $$C_{4\mu +2}^+(t)T_{4\mu +2}(t)C_0^+,S_{4\mu +2}^+(t)T_{4\mu +2}(t)S_0^+,$$ $`(5.12)`$ $$C_{4\mu +3}^+(t)T_{4\mu +3}(t)S_0^+,S_{4\mu +3}^+(t)+T_{4\mu +3}(t)C_0^+,$$ $`(5.13)`$ where $`\mu =0,1,2,..`$. Incorporation of the formulas (5.10)-(5.13) into the definition (3.9b) furnishes $$\begin{array}{cc}\hfill v_\tau ^+(t)\underset{\mu =0}{\overset{\mathrm{}}{}}\{& [T_{4\mu }(t)T_{4\mu +2}(t)]w_{\tau 0}^+(t)+\hfill \\ & [T_{4\mu +1}(t)T_{4\mu +3}(t)]w_{\tau 0}^{}(t)\},[t>1,R\mathrm{}],\hfill \end{array}$$ $`(5.14)`$ where we used the ad hoc yet not inconsistent notation $$w_{\tau 0}^{}(t)=\eta _\tau (t)\{S_0^+\mathrm{cos}[kR(t1)]C_0^+\mathrm{sin}[kR(t1)]\},[t>1],$$ $`(5.15)`$ Observe there is only a single quadrature involved in the asymptotical formula (5.14), namely the one implicitly contained in the definition (5.9). Knowledge of the higher order $`t`$-dependent coefficients thus rests upon the knowledge of the $`R\mathrm{}`$ form of the zero order constants $`C_0^+`$ and $`S_0^+`$. These constants are fixed by the claim for the smoothness of the exact but supersingularity solution at $`t=1`$ as follows. At the matching point itself, approaching it from either of the regions $`ϵ`$ or $`\tau `$, one obtains the respective amplitude functions, see Eq.(2.16), along with the relevant derivatives as follows: $$[\eta _\gamma (t)]_{t=1}2^{\frac{3}{4}}(kR)^{\frac{1}{2}},[\gamma =ϵ,\tau ;R\mathrm{}],$$ $`(5.16)`$ $$[\frac{\mathrm{d}\eta _\gamma (t)}{\mathrm{d}t}]_{t=1}\pm 2^{\frac{7}{4}}(kR)^{\frac{1}{2}}[2\lambda _\gamma ^2g^2R^3U^{}(R)],[\gamma =(_\tau ^ϵ);R\mathrm{}].$$ $`(5.17)`$ The second term in the brackets vanishes in this limit by the restrictions imposed upon the potentials in virtue of the asymptotical relationships (2.1). As to the exponential region, the supersingularity formula (5.3) can be identically recast by means of the definition (2.14) as $$v_ϵ^+(t)\eta _ϵ(t)\mathrm{exp}[\omega _ϵ(1,t)+P_ϵ(0,t)],[t<1,R\mathrm{}].$$ $`(5.18)`$ Approaching the matching point from this region , one concludes by Eqs. (5.16)-(5.18) that $$[v_ϵ^+(t)]_{t=1}2^{\frac{3}{4}}(kR)^{\frac{1}{2}}\mathrm{exp}P_ϵ(0,1),[R\mathrm{}],$$ $`(5.19)`$ $$[\frac{\mathrm{d}v_ϵ^+(t)}{\mathrm{d}t}]_{t=1}2^{\frac{3}{4}}(kR)^{\frac{1}{2}}\mathrm{exp}P_ϵ(0,1)\{4\lambda _ϵ^2+2^{\frac{3}{2}}+[p_ϵ(t)]_{t=1}\},[R\mathrm{}].$$ $`(5.20)`$ Recall now the expressions (4.24)-(4.25) and extract from them $$[p_ϵ(t)]_{t=1}=𝒪\{\frac{1}{kR}\}0,[B=b;R\mathrm{}],$$ $`(5.21)`$ $$[p_ϵ(t)]_{t=1}=𝒪\{\frac{R}{\rho _2}\}\mathrm{},[B=\beta ;R\mathrm{}].$$ $`(5.22)`$ Therefore, we have to restrict further discussion to scattering by potentials of the class $`U_{Ab}(r),[A=\alpha ,a]`$ of Eqs. (4.3). As regards the trigonometric region, the definition (2.15) yields $$[w_{\tau 0}^+]_{t=1}2^{\frac{3}{4}}(kR)^{\frac{1}{2}}C_0^+,[B=b,R\mathrm{}],$$ $`(5.23)`$ $$[\frac{\mathrm{d}w_{\tau 0}^+(t)}{\mathrm{d}t}]_{t=1}(kR)^{\frac{1}{2}}[2^{\frac{11}{4}}\lambda _\tau ^2C_0^++S_0^+],[B=b,R\mathrm{}].$$ $`(5.24)`$ Smooth matching is realized by equating Eqs.(5.19) and (5.23) as well as (5.20) and (5.24), respectively. In doing so, one obtains $$C_0^+\mathrm{exp}P_ϵ(0,1),S_0^+\mathrm{exp}P_ϵ(0,1)[2^{\frac{15}{4}}\lambda ^2+2^{\frac{3}{4}}],[B=b,R\mathrm{}].$$ $`(5.25)`$ As to the region $`ϵ`$, Eqs. (4.29), (4.33) and (5.3) , as well as concerning the region $`\tau `$, Eqs. (4.31), (4.32), (5.9) and (5.14), combine to the statement that , for potentials $`U_{AB}(r)`$ with powerlaw tails, \[B=b\], the pair of zero order terms in the semiclassical expansions of the respective regions themselves recover in the supersingular limit the exact QM solution. In particular, $$u^+(t)w_{ϵ0}^+(t),[t<1;B=\beta ,b;R\mathrm{}],$$ $`(5.26)`$ $$u^+(t)w_{\tau 0}^+(t),[t>1;B=b;l0];R\mathrm{}].$$ $`(5.27)`$ In cases of exponential potential tails, the S-partial waves should be excluded from the present approach. VI. Discussion Usefulness of a modified semiclassical approach in treating $`singular`$ $`scattering`$ has been checked in three steps. First, the conditions were reconsidered at which a smooth WKB method, proposed recently,produces convergent expansion of the wave function at fixed set of the potential parameters. An inherent new point is that, instead of a single one, a $`pair`$ of $`integral`$ $`equations`$ should be set up, one for each of the exponential and the trigonometric regions. A regular solution working in the exponential region selects by virtue of the smoothness postulate a particular one out of the solutions prepared for the trigonometric region. Another new item is the extension of the argument over potentials involving $`variable`$ $`nonlinear`$ $`parameters`$ through the variation of which one may increase the core singularity to asymptotically high levels. Analyzed are four classes of interactions each of which is a product of a core factor implying exponential or powerlaw singularity and of a tail factor that decays either exponentially or powerlike. Independently of the stage of the singularity, the powerlike decaying potentials invariably develop absolute convergent expansions, along both regions, also in the supersingularity limit. In scattering by exponentially decaying potentials, the criteria of convergence seemingly fail to work within our argument of treating supersingularity limit. A further new feature of the present approach is a discussion of the $`quality`$ of $`convergence`$. The conditions are found which the shape of the potential has to show up so that the infinite series shrink, in the supersingularity limit, to a single term. It is, perhaps, worth mentioning that, by varying the length of the cut-off series in the region beyond the matching point, one can obtain, at the expense of a single quadrature, a solution that becomes correct to any prescribed order in the reciprocal supersingularity parameter at its asymptotically large values. Recall the Born series furnishes the physical scattering wave function for nonsingular potentials at fixed linear and nonlinear dynamical parameters, exclusively. The proposed smooth WKB approach should work for singular potentials at fixed and asymptotical values of linear and nonlinear parameters involved in the Schrödinger equation. Acknowledgement Many thanks are due to Dr. G.Bencze for useful critical remarks. The author is grateful to Dr. G.Kluge and Dr. I. Racz for very valuable discussions. The work was partly supported by the Hungarian NSF under Grant No. OTKA 00157. References <sup>1</sup>F. Calogero: Variable Phase Approach to Potential Scattering, (Academic Press, New York, 1967) <sup>2</sup>F. Calogero: Phys. Rev. B 135, 693 (1964) <sup>3</sup>G. Esposito: J.Phys A 31, 9493 (1998) <sup>4</sup>N. Froeman and K.-F. Thylwe: J. Math. Phys 20, 1716 (1979) <sup>5</sup>R.G. Newton: Scattering Theory of Waves and Particles, (Springer Verlag, 1981) <sup>6</sup>T. Dolinszky: Physics Letters A 132, 69 (1988) <sup>7</sup>T. Dolinszky: J.Math.Phys. 36, 1621 (1995) <sup>8</sup>T. Dolinszky: J.Math.Phys. 38, 16 (1997).
warning/0002/math0002195.html	ar5iv	text	# Untitled Document ADDENDUM to THEOREM 10.4 in “BOUNDARIES OF ANALYTIC VARIETIES” by F. Reese Harvey and H. Blaine Lawson, Jr. The main result of \[HL\], put simply, provides a characterization of the boundaries of complex subvarieties in $`C^n`$. One of the minor applications of this result, namely Theorem 10.4, requires clarification because of the note \[LY\] of Luk-Yau.(See \[E\].) The intent of \[LY\] is to provide a counterexample to the boundary regularity assertion of Theorem 10.4. However, Theorem 10.4 is fundamentally correct. Furthermore, the authors of \[LY\] seem not to have realized that their example already appears in \[HL\] (Example 9.1). This example is simply an immersion $`C^2C^3`$ which folds back on itself, and therefore when restricted to balls gives rise to crossing singularities both in the interior and at the boundary. It shows that for boundaries $`M`$ which are embedded and strongly pseudoconvex, the fill-in variety may not be embedded. The boundary regularity stated in Theorem 10.4 conforms to this example. The possibility of crossing singularities is explicitly stated in the last line of the theorem where multiple local components of $`V`$ at the boundary are discussed. There was a minor error in the exposition of Theorem 10.4. Since this mis-statement was internally contradictory, the correct version may have been evident to the reader. Nevertheless, in this note we amend the error. We also give an alternative version of the result which we thought was obvious, but perhaps was not. In addition this note corrects a result of Stephen Yau \[Y\], demonstrates that Lempert’s use of Theorem 10.4 carries through, and shows how the unproven theorem in \[LY\] follows trivially from our paper. Incidentally, Theorem 10.4 was not a new result for $`\text{dim}(V)3`$. In the sentence preceding the theorem we pointed out that the result follows from the classical Lewy extendibility of CR functions (as described in Theorem 10.3) combined with the work of Rossi \[R\]. The really new work in \[HL\] and its sequel \[HL2\] are the global results characterizing boundaries of varieties without mention of the Levi form. Theorem 10.4. was intended to assert the existence of a variety with smooth boundary and a finite number of isolated interior singularities holomorphically immersed into $`C^n`$. In the statement, the word “immersed” was erroneously omitted. Its intention is implicit in a serious reading of the result and the material prior to it. (For instance see Example 9.1, the sentence prior to Theorem 9.2, and the first paragraph of Theorem 10.3.) However, to completely clarify Theorem 10.4 we shall correct the wording of the result and then in the Lemma below we shall explicitly establish an equivalent formulation in terms of immersions. The subsequent applications of Theorem 10.4 appear in two papers: \[Y\] and \[L\]. In fact \[Y\] presents an alternate proof of the Theorem 10.4 which overlooks the possibility of immersions shown in the example above. Curiously, no reference to this appears in \[LY\]. Nevertheless, as we shall show below, the results in \[Y\] and the arguments in \[L\] are easily amended. To correct the error in exposition in Theorem 10.4 we recall some elementary facts. Let $`V`$ be a variety with $`d[V]=[M]`$ as in the main theorem 8.1 of \[HL\]. Fix $`pM`$ and suppose that in a neighborhood $`𝒰`$ of $`p`$ there is a local component $`W`$ of $`V`$ which is a $`C^k`$-submanifold with boundary $`M`$. Then $`\overline{VW}`$ is an analytic subvariety of $`𝒰`$ and therefore has a finite number of irreducible components at $`p`$. (See \[K\] or \[H\].) Definition. Suppose now that every point $`pM`$ has the property above (as is the case when $`M`$ is strictly pseudoconvex). Then a point $`p\overline{V}`$ is defined to be an intrinsic singular point if it is a singular point of some local irreducible component of $`V`$ at $`p`$ if $`pV`$, or of $`\overline{VW}`$ if $`pM`$. Theorem 10.4 should be amended in line 4 by replacing the word “isolated” with the word “intrinsic”. Theorem 10.4. (amended): Let $`MC^n`$ be a connected $`C^k`$ manifold satisfying the hypothesis of Theorem 8.1, and suppose $`M`$ is pseudoconvex. Then there exists an irreducible, $`p`$-dimensional complex analytic subvariety $`VC^n\backslash M`$ with $`\overline{V}`$ having at most finitely many intrinsic singularities, such that $`[M]=d[V]`$, with $`C^k`$ boundary regularity for each local component of $`V`$ near $`M`$. The proof should be amended in line 4 to read: Theorem 10.3a now shows that the intrinsic singularities of $`\overline{V}`$ form a compact subvariety of $`C^n`$ which must have dimension 0. As mentioned above the example proclaimed in the title of \[LY\] appears explicitly in \[HL\] in Example 9.1. It is the simplest holomorphic immersion $`C^2C^3`$ with self-intersections. In fact the example $`F`$ in \[LY\] differs from Example 9.1 in \[HL\] by a linear change of variables. More precisely, if we define $`L(x,y,z)=(4x+1,z,8y4x)`$ and $`\lambda (t,z)=(\frac{1}{2}(t+1),z),`$ then $`\mathrm{\Phi }=LF\lambda `$ is exactly Example 9.1. Example 9.1 \[HL\] considers the variety $`V`$ given by the $`F`$-image of a ball whose radius $`r_0`$ is chosen to be the first $`r`$ for which the image has a self intersection. This value of $`r`$ was considered particularly interesting because the boundary $`M`$ of $`V`$ is a (strictly pseudoconvex) real analytic submanifold of $`C^3`$ and $`V`$ is a complex submanifold of $`C^3M`$ but the pair is not a topological submanifold-with-boundary. The apparent content of \[LY\] is to mention that one can also consider $`r>r_0`$ in this example. Incidentally, the theorem announced without proof in \[LY\] follows immediately from \[HL\]. Luk-Yau assume the additional hypothesis that $`M`$ is contained in the boundary of a bounded strictly pseudoconvex domain $`D`$ in $`C^N`$. In a neighborhood of $`M`$, the subvariety $`V`$ obtained from Theorem 10.4 has a component $`W`$ (a “strip”) which is a smooth submanifold with boundary $`MN`$ where $`N`$ is a nearby “parallel” manifold. Thus, $`VW`$ has boundary $`N`$. Since $`N`$ is contained in a smaller strictly pseudoconvex domain $`D(ϵ)D`$, the Stein manifold version of the main result (Theorem 8.6) in \[HL\], gives a subvariety $`Z`$ of $`D(ϵ)N`$ with $`d[Z]=[N]`$. By uniqueness $`VW`$ and $`Z`$ must agree. Hence $`VW`$ misses a neighborhood of $`M`$. This rules out singular points of $`V`$ near $`M`$. Hence the entire singular subvariety of $`V`$ reduces to a finite set. As noted above, the authors of \[LY\] neglected to mention that the example they present contradicts a result of their own, namely \[Y; Thm 5.14 (Thm C in the introduction)\]. In proving Theorem 5.12 in \[Y\], from which 5.14 is stated to be an “easy consequence”, Yau constructs a normal variety over $`C^N`$, and he carefully points out (on page 89) that self- intersections may occur after projecting to $`C^N`$. However, this point is completely ignored in the statement of Theorem 5.14 which should be amended to read: “Then $`M`$ is the boundary of an immersed complex submanifold …”. To prove this amended statement we use the following. Theorem 10.4. Suppose $`MC^n`$ is a compact, connected, oriented, maximally complex submanifold of class $`C^k`$ and dimension $`2p1>1`$. Assume $`M`$ is strictly pseudoconvex. Let $`VC^nM`$ be the analytic subvariety of dimension $`p`$ and of finite volume with $`d[V]=[M]`$ given by \[HL, Thm. 8.1\]. Then there exists: (i) A compact space $`\overline{X}=XX`$ with where $`X`$ is a normal Stein variety having at most a finite number of singular points, and such that $`(X,X)`$ is a $`C^k`$-manifold-with-boundary away from the singular points, and (ii) A map $`\rho :\overline{X}C^n`$, which is holomorphic on $`X`$ and of class $`C^k`$ up to the boundary, inducing a $`C^k`$-diffeomorphism from $`X`$ to $`M`$ and having $`\rho (\overline{X})=\overline{V}`$. Furthermore, $`\rho `$ is an immersion outside a finite subset of $`X`$ which contains the singularities of $`X`$ and is contained in the preimage of the intrinsic singularities of $`\overline{V}`$. Finally, when $`V`$ is a hypersurface, $`\rho `$ is a local holomorphic embedding. Proof. Let $`\rho _0:\stackrel{~}{V}V`$ be the normalization of $`V`$. Since $`V`$ has a finite number of intrinsic singularities, the singular set of $`\stackrel{~}{V}`$ is finite. We complete $`\stackrel{~}{V}`$ to $`\overline{X}`$ as follows. Each $`pM`$ has a neighborhood $`𝒰`$ such that $`\overline{V}𝒰=WV_1\mathrm{}V_m`$ where $`W`$ is a $`C^k`$-submanifold with boundary and where $`V_1,\mathrm{},V_m`$ are irreducible subvarieties of $`𝒰`$ each of which has a finite singular set (again because the intrinsic singularities are finite). Let $`\rho _j:\stackrel{~}{V_j}V_j`$ be the normalization of $`V_j`$. Note that $`\stackrel{~}{V_j}`$ has a finite singular set and $`\rho _j`$ is a holomorphic homeomorphism. These maps induce a map $$\rho _𝒰:W\stackrel{~}{V_1}\mathrm{}\stackrel{~}{V_m}WV_1\mathrm{}V_m=\overline{V}𝒰,$$ which is canonically isomorphic to $`\rho _0`$ on the preimage of $`V𝒰`$ by the uniqueness of normalization. Gluing these pieces to $`\stackrel{~}{V}`$ and adding the boundary in the obvious way produces $`X`$. Note that $`X`$ contains no compact subvarieties of positive dimension since $`\rho `$ has discrete fibres, but would be constant on connected components of such subvarieties. Since $`X`$ is strictly pseudoconvex we therefore conclude that $`X`$ is a Stein space by \[G\]. The last statement is a consequence of the fact that isolated hypersurface singularities are normal. When $`p=n13`$ the arguments in \[Y\] apply to show that $`X`$ is non-singular if and only if the Kohn-Rossi coholomogy groups of the boundary complex are 0. This gives the amendment to \[Y\] discussed above. Since $`X`$ is Stein, the arguments on page 13 of \[L\] which use \[HL; 10.4\] carry through unchanged. References \[E\] Editor’s Note on Papers by Harvey-Lawson and by Luk-Yau, Annals of Math. (to appear). \[G\] H. Grauert, Über Modifikationen und exzeptionelle analytische Mengen Math. Ann. 146 (1962), 331-368. \[H\] F. R. Harvey, Holomorphic chains and their boundaries in Several Complex Variables, vol.1, Proc. Symp. Pure Math. (1977), 309-382. \[HL\] F. R. Harvey and H. B. Lawson, Jr., On boundaries of complex analytic varieties, I, Annals of Math. 102 (1975), 223- 290. \[HL2\] F. R. Harvey and H. B. Lawson, Jr., On boundaries of complex analytic varieties, II, Annals of Math. 106 (1977), 213- 238. \[K\] J. King, The currents defined by analytic varieties, Acta Math. 127 (1971), 185-220. \[L\] L. Lempert, Embeddings of three-dimensional Cauchy-Riemann manifolds, Math. Ann. 300 (1994), 1-15. \[Y\] Stephen S.-T. Yau, Kohn-Rossi cohomology and its application to the complex Plateau problem, I, Annals of Math. 113 (1981), 67-110. \[LY\] H. S. Luk and Stephen S.-T. Yau, Counterexample to boundary regularity of a strongly pseudoconvex CR submanifold: An addendum to the paper of Harvey-Lawson, Annals of Math. 148 (1998), 1153-1154. \[R\] H. Rossi, Attaching analytic spaces to an analytic space along a pseudoconcave boundary, Proc. Conf. on Complex Analysis, Minneapolis, Springer Verlag, 1964.
warning/0002/cond-mat0002456.html	ar5iv	text	# Supercooled confined water and the Mode Coupling crossover temperature ## Abstract We present a Molecular Dynamics study of the single particle dynamics of supercooled water confined in a silica pore. Two dynamical regimes are found: close to the hydrophilic substrate molecules are below the Mode Coupling crossover temperature, $`T_C`$, already at ambient temperature. The water closer to the center of the pore (free water) approaches upon supercooling $`T_C`$ as predicted by Mode Coupling Theories. For free water the crossover temperature and crossover exponent $`\gamma `$ are extracted from power-law fits to both the diffusion coefficient and the relaxation time of the late $`\alpha `$ region. The effect of supercooling on the dynamics of liquids in confined environments is a research field that has become more and more popular over the last few years. Out of an extremely rich phenomenology that shows diversification of specific behavior depending on the size of the particles, the confining geometry and the specific interaction with the substrate, some general trends can none the less be extracted . In fact two competing effects seem to be the main contributions to the modification of the dynamics of the confined liquid with respect to the bulk phase: the bare geometric confinement and the interaction with the substrate. In particular there is evidence from experiments that liquid-wall interactions can lead to a layering and a decrease of mobility close to the substrate, with a substantial increase of the glass transition temperature. This effect is stronger for attractive interactions between the substrate and the liquid. In some liquids experimental evidences show two distinct dynamical regimes . Among liquids water plays a most fundamental role on earth. The study of the dynamics of water at interfaces or confined in nanopores as a function of temperature and hydration level is relevant in understanding important effects in systems of interest to biology, chemistry and geophysics . In particular the single particle dynamics of water confined in nanopores have been studied by different experimental techniques such as neutron diffraction and nuclear magnetic resonance . A slowing down of the dynamics with respect to the bulk phase is observed. Nevertheless the details of the microscopic dynamic behavior of confined water upon supercooling are still unclear. Below $`235`$ K bulk liquid water enters the so called no man’s land , where nucleation processes, most likely triggered by the presence of impurities, take place and drive the liquid to the solid crystalline phase, preventing the experimental approach to the glass phase . The dynamical behavior of the bulk water simulated upon supercooling with the use of the Simple Point Charge/Extended (SPC/E) site model potential , fits in the framework of the idealized version of Mode Coupling Theory (MCT) , predicting a temperature of ideal structural arrest, or crossover temperature $`T_C`$, that coincides with the so called singular temperature of water . This behaviour has been observed several years ago by computer simulation and substantiated by experimental signatures and further simulation and theoretical works . When a liquid approaches the crossover temperature $`T_C`$ MCT predicts that the dynamics is dominated by the cage effect. After an initial ballistic motion, the particle is trapped in the transient cage formed by its nearest neighbours. Once the cage relaxes the particle enters the Brownian diffusive regime. Below $`T_C`$, according to the idealized version of MCT, the system becomes non-ergodic. In real structural glasses hopping processes restore ergodicity. $`T_C`$ is therefore a crossover temperature from a liquid-like to a solid-like regime. These activated processes are not relevant above $`T_C`$ for most liquids. Until now, no systematic computer simulation studies of the microscopic dynamics of confined water upon supercooling have been attempted. Thermometric studies , NMR spectroscopy , neutron diffraction and X-ray diffraction show evidence that two types of water are present in the confining pores, free water which is in the middle of the pore and bound water which resides close to the surface. Free water is observed to freeze abruptly in the cubic ice structure. Bound water freezes gradually but it does not make any transition to an ice phase . Layering effects of water close to the substrates have been observed in all the simulations for different geometry and water-substrate interaction . In this letter we present evidence from Molecular Dynamics (MD) simulations that water confined in a hydrophilic nanopore exhibits two distinct dynamical regimes. In particular the fraction of free water molecules behaves similarly to the bulk phase, and its dynamics is consistent with several MCT predictions. Power laws fits based on MCT yield the crossover temperature of the fraction of free water molecules; simultaneously, we show that the bound water molecules are already below $`T_C`$ at room temperature. In our simulations water has been confined in a silica cavity modeled to represent the average properties of the pores of Vycor glass . Water-in-Vycor is a system of particular interest , since the porous silica glass is characterized by a quite narrow pore size distribution with an average diameter of $`40`$ Å. The pore size does not depend on the hydration level and the surface of the pore is strongly hydrophilic. Moreover the water-in-Vycor system can be considered as a prototype representing more complex environments of interfacial water. Different from previous work on the dynamics of water close to a planar regular silica surface we use a cylindrical geometry with a corrugated surface. We have constructed a cubic cell of silica glass by the usual quenching procedure. As described in detail in a previous work inside the cube of length $`L=71.29`$ Å we cut out a cylindrical cavity of diameter $`40`$ Å and height $`L`$ by eliminating all the atoms lying within a distance $`R=20`$ Å from the axis of the cylinder, which is taken as the $`z`$-axis. After elimination of silicon atoms with less than four oxygen neighbours, we saturate the dangling bonds of oxygen atoms with hydrogen atoms, in analogy to the experimental situation . Water molecules described by the SPC/E model are introduced into the cavity. The water sites interact with the atoms of the rigid matrix by means of an empirical potential model, where different fractional charges are assigned to the atomic sites of the silica glass and where the oxygen sites of water additionally interact with the oxygen atoms of the substrate via Lennard-Jones potentials . The MD calculations are performed with periodic boundary conditions along the $`z`$-direction and the temperature is controlled via coupling to a Berendsen thermostat ; the shifted force method is used with a cutoff at $`9`$ Å to truncate long-range interactions . In this work we consider water at a density corresponding to the experimentally determined level of full hydration . In the chosen geometry this corresponds to $`N_w=2600`$ water molecules and to a density $`\rho =0.867`$ $`g/cm^3`$. We investigate the dynamical behavior of the confined water for five temperatures, namely $`T=298,270,240,220`$ and $`210`$ K. In the following we focus on the single particle dynamics of the water molecules contained in the pore and test some of the main predictions of MCT. The radial density profile of the water oxygen atoms, normalized to bulk water density, is displayed for $`T=298`$ K and $`T=210`$ K in the inset at the upper left corner of Fig. 1. It shows the high hydrophilicity of the pore in the form of density oscillations close to the substrate. These oscillations are not sensitive to supercooling. The average density of water molecules at $`T=298`$ K for $`0<R<15`$ Å is $`\rho =0.897`$ $`g/cm^3`$, for $`15<R<18`$ Å is $`\rho =1.079`$ $`g/cm^3`$ and for $`R>18`$ Å (a depletion layer) is $`\rho =0.493`$ $`g/cm^3`$, where $`R=\sqrt{(x^2+y^2)}`$. In a preliminary analysis done at ambient temperature as a function of hydration level we found that, due to the presence of strong inhomogeneities in our system, a fit of the total correlators to an analytic shape can be carried out only by excluding the subset of molecules in the double layer close to the substrate ($`R>15`$ Å) which we now identify with the so called bound water. In the following we show that this shell analysis keeps its validity upon supercooling. We analyze the mean square displacement (MSD) and the self intermediate scattering function (ISF) separately for bound water and for the remaining inner layers ($`R>15`$ Å), which we identify with the free water. Since no asymptotic free motion is possible in the $`xy`$ plane, we separately analyze the dynamics within this plane and along the pore $`z`$-axis. In the main frame of Fig. 1 the MSD of free water is displayed along the non-confined $`z`$-direction for the investigated temperatures. As supercooling progresses a plateau region due to the cage effect develops after the initial ballistic motion, starting around $`t=0.3`$ ps. The MSD in the $`xy`$-direction (not shown) is characterized by a similar trend, but the dynamics is slower. In the low right corner the inset shows the behavior of the diffusion coefficient $`D`$ extracted from the fit of the MSD in the Brownian regime region for both the $`z`$ and the $`xy`$ direction together with fits of the MCT predicted power law $`D(TT_C)^\gamma `$ to the $`5`$ data points. In the $`z`$-direction, we obtain $`T_C185.3`$ K and $`\gamma 2.21`$, which are similar to the values found for the SPC/E bulk water for ambient pressure, namely $`T_C186.3`$ K and $`\gamma 2.29`$ . In the $`xy`$-direction, where the dynamics is slower, the critical temperature of free water increases to $`T_C194.5`$ K and $`\gamma 1.90`$. In Fig. 2 we show the ISF of free water at the peak of the oxygen-oxygen structure factor along the $`xy`$-direction, $`Q_{MAX}=2.25`$ $`\mathrm{\AA }^1`$, as a function of temperature. The free water molecules inside the pore show, similar as in SPC/E bulk water, a diversification of the relaxation times as supercooling proceeds. The plateau region stretches as $`T_C`$ is approached. The long time region, the so-called late $`\alpha `$ region, is expected to have a stretched exponential decay for a liquid approaching $`T_C`$. In the same figure, the fit of the function $$F_S(Q,t)=\left[1A(Q)\right]e^{\left(t/\tau _s\right)^2}+A(Q)e^{\left(t/\tau _l\right)^\beta }$$ (1) to the data points is shown, where $`A(Q)=e^{a^2Q^2/3}`$ is the Debye-Waller factor arising from the cage effect with $`a`$ the effective cage radius. $`\tau _s`$ and $`\tau _l`$ are, respectively, the short and the long relaxation times, and $`\beta `$ is the Kohlrausch exponent. Obviously the Gaussian form of the fast relaxation can be only an approximate one. In the inset of Fig. 2 we show the full layer analysis for $`T=240`$ K as a representative case. The topmost curve shows the behavior of bound water, while the lower curve is the ISF of free water (identical to the curve in the main picture). It is clearly seen that bound water is below $`T_C`$, since the correlation function does not decay to zero on the nanosecond time scale. The central curve is the total ISF of confined water. It displays a strong non-exponential tail, which cannot be fitted by a stretched exponential function. Our layer analysis shows clearly that the contribution of free water can be separated from the one of bound water and that the stretched exponential function is able to give a very good fit to the late part of the $`\alpha `$ region in the free water subsystem as we supercool. The values of the Kohlrausch exponents $`\beta `$ and the relaxation times $`\tau _l`$ extracted from the fits of Eq. (1) to the ISF data (Fig. 2) are reported in Fig. 3 for free water as a function of temperature for the $`z`$ and the $`xy`$ directions. The $`T`$ dependence of $`\beta `$ and $`\tau _l`$ is in agreement with MCT. MCT also predicts that the inverse of the $`\alpha `$-relaxation time $`\tau _l`$ vanishes with the same power law as the diffusion coefficient $`1/\tau _l(TT_C)^\gamma `$. In the upper part of Fig. 3 we show log-log plots of the inverse relaxation times as function of $`(TT_C)`$ as points together with power law fits as continuous lines. From the fits we obtain values very close to the ones obtained above from the power law fit of the diffusion coefficients (see inset in Fig. 1). This result seems therefore in agreement with the MCT prediction that the $`\gamma `$ exponent should be independent of the quantity investigated. In pure SPC/E water the quantity $`D\tau `$ was found relatively constant along isochores. None-the-less a slight increase of the product was found on cooling most likely due to the progressive breakdown of MCT on approaching $`T_C`$, see Fig. 5c of Ref.. In our case a definitive conclusion on the coupling of $`D`$ and $`\tau `$ cannot be made at the present stage due to the fluctuations in the data. The cage radii extracted from the value of the Debye-Waller factor along both the $`xy`$ and $`z`$ directions range from $`a=0.51`$ Å for $`T=298`$ K to $`a=0.45`$ Å for $`T=210`$ K. The $`\tau _s`$ values are $`\tau _s0.2`$ ps. These values are very similar to those of SPC/E bulk water . In summary, we have presented evidence that the dynamical behavior of SPC/E water confined in a silica nanopore upon supercooling can be analyzed in terms of two subsets of water molecules with clearly distinct dynamical regimes, in agreement with signatures found in experimental studies on confined liquids . Due to the presence of a strongly attractive surface we do not find, at variance with MD studies on other confined liquids, density oscillations all over the confining space , or a continuous behaviour in going from the center to the substrate . The bound water being close to a strong hydrophilic surface suffers a severe slowing-down already at room temperature. Power laws fit based on MCT predict that for free water ideal structural arrest would occur at $`T_C`$ $`186.3K`$ and $`194.5`$ respectively along $`z`$ and $`xy`$ directions for the full hydration level of the pore. For the investigated quantities dynamics appears well accounted for by the idealized MCT of supercooled liquids. In this respect the predictions of the idealized MCT appear to be robust and able to describe also confined molecular liquids provided that the effects of the interaction with the substrate are properly taken into account. Experimental evidence of a possible MCT behavior of water-in-Vycor have been observed . Therefore the analysis presented here represents an important step towards the understanding of slow structural relaxation of highly non trivial glass forming systems. M. R. and P. G. acknowledge the financial support of the G Section of the INFM.
warning/0002/math-ph0002026.html	ar5iv	text	# Untitled Document RGGR-94/1 22 February 1994 Relationships between various characterisations of wave tails Luca Bombelli and Sebastiano Sonego RGGR, Université Libre de Bruxelles, Campus Plaine CP 231, 1050 Brussels, Belgium E-mail: lbombell@ulb.ac.be (Current e-mail: luca@phy.olemiss.edu) E-mail: ssonego@ulb.ac.be (Current e-mail: Sebastiano.Sonego@Dic.Uniud.It) Abstract One can define several properties of wave equations that correspond to the absence of tails in their solutions, the most common one by far being Huygens’ principle. Not all of these definitions are equivalent, although they are sometimes assumed to be. We analyse this issue in detail for linear scalar waves, establishing some relationships between the various properties. Huygens’ principle is almost always equivalent to the characteristic propagation property, and in two spacetime dimensions the latter is equivalent to the zeroth order progressing wave propagation property. Higher order progressing waves in general do have tails, and do not seem to admit a simple physical characterisation, but they are nevertheless useful because of their close association with exactly solvable two-dimensional equations. PACS numbers: 03.40.Kf; 02.30.Jr. Short title: Characterisations of wave tails. 1.Introduction The question whether a propagating wave leaves a tail behind it is one that has interested many physicists and mathematicians for a long time, and has found applications from the first studies in light propagation to the theory behind the proposed experiments to detect gravitational waves . Intuitively, a wave tail can be described as follows. Suppose that, at some time $`t`$, an instantaneous pulse of a field $`\varphi `$ is produced at a point labelled by spatial coordinates $`\{x^i\}`$. This pulse originates a wave front propagating outward with some speed, that will be detected by an observer sitting at some point $`\{x^i\}`$ at a time $`t^{}>t`$. We say that the wave has a tail or wake if such an observer continues to detect a nonvanishing field even after the wave front has passed, i.e., at times greater than $`t^{}`$. In this case we speak of diffusive, as opposed to sharp, propagation. One can roughly identify three distinct possible reasons for the occurrence of tails: (i) Mass-like terms in the wave equation. For example, the Klein-Gordon equation $$\text{ }\text{ }\text{ }\text{ }\text{ }\varphi \mu ^2\varphi =0$$ $`(1.1)`$ in four-dimensional Minkowski spacetime has tails when $`\mu 0`$, but not when $`\mu =0`$. This feature can be seen as the wave-mechanical counterpart of the fact that massive particles move slower than the speed of light. (ii) Dimensionality of spacetime. For example, the massless Klein-Gordon equation in $`m`$-dimensional Minkowski spacetime has tails for odd values of $`m`$, but not when $`m`$ is even, with the exception of $`m=2`$ \[4–7\]. (iii) Backscattering off potentials and/or spacetime curvature . This is clearly the most interesting tail production mechanism from a physical viewpoint. The study of the occurrence of tails, as well as of their implications, is certainly of great interest and importance, but relies on the possibility of using precise definitions. The heuristic characterisation given above is evidently too loose for this purpose, but it can be easily formalised to obtain a tail-free property. For linear wave equations, however, a no-tails condition is better expressed in terms of the Green function, and corresponds to what is usually called the Huygens principle. Furthermore, still other definitions (e.g., characteristic propagation property, progressing waves propagation property) can be found in the literature. These characterisations are not all equivalent, although they are sometimes (explicitly or implicitly) assumed to be. It is the purpose of this paper to explore this issue in detail, by establishing the relationships between them and clarifying their meaning and domain of applicability. If the wave equation is inhomogeneous (which corresponds, physically, to the presence of sources), the formulation of no-tails properties is a delicate task, because of the need to capture the notion of sharp propagation when the field due to the sources is superposed to the one arising purely from the boundary conditions. The situation is even worse for nonlinear waves. For these reasons, we shall restrict ourselves here to considering an arbitrary linear, second order homogeneous hyperbolic partial differential equation, which can always be written in the form $$g^{\mu \nu }_\mu _\nu \varphi +h^\mu _\mu \varphi +K\varphi =0,$$ $`(1.2)`$ where $`g^{\mu \nu }`$, $`h^\mu `$, and $`K`$ are suitable real functions of $`m`$ variables $`\{x^\mu \}`$, that we can think of as coordinates over an $`m`$-dimensional differentiable manifold $``$; without loss of generality, we can assume that $`g^{\mu \nu }=g^{\nu \mu }`$ so that these functions can be thought of as the components of the inverse of a Lorentzian metric $`g_{ab}`$ on $``$ given, in the coordinates $`\{x^\mu \}`$, by the inverse of the matrix $`g^{\mu \nu }`$; this metric, however, may not be the one normally used to determine intervals and causal relations (for a physical example in which it is not, see Ref 10). The metric $`g_{ab}`$ defines then a Riemannian connection $``$ that allows us to rewrite (1.2) equivalently as $$g^{ab}_a_b\varphi +H^a_a\varphi +K\varphi =0,$$ $`(1.3)`$ where $`H^a`$ is uniquely determined by $`h^\mu `$ and $`g^{\mu \nu }`$ (specifically, $`H^\mu =h^\mu _\nu g^{\mu \nu }^\mu \mathrm{ln}\sqrt{g}`$, where $`g`$ is the determinant of $`g_{\mu \nu }`$). Although physically one may wish to discuss the appearance of tails in waves emitted by very distant sources, the phenomenon is essentially a local one. Therefore, it will not represent a loss in generality to study wave propagation close to the hypersurfaces where data are given, and we can assume that $``$ is globally hyperbolic and normal. A more general manifold can always be covered by a collection of regions with these properties, and tails will develop iff they already do inside some of those smaller domains . In such a manifold, a very powerful tool for studying general properties of linear differential equations and their solutions is that of Green functions . A Green function $`G(x,x^{})`$ for (1.3) is defined by $$\left(g^{ab}_a_bH^a_a+K_aH^a\right)G(x,x^{})=\delta (x,x^{}),$$ $`(1.4)`$ and by appropriate boundary conditions; in (1.4), $`\delta (x,x^{})`$ is the delta function on $``$, such that for each test function $`f`$, $$\mathrm{d}^mx^{}\sqrt{g(x^{})}f(x^{})\delta (x,x^{})=f(x).$$ $`(1.5)`$ It should be noted that the differential operator acting on the first variable in $`G`$ is the adjoint of the one which acts on $`\varphi `$ in (1.3), and differs from it when $`H^a0`$ . If we choose to study the behaviour of the field at points to the future of a Cauchy hypersurface $`𝒮`$ on which data are assigned, it is convenient to work with the advanced Green function, which satisfies $`G(x,x^{})=0`$ when $`x^{}`$ is to the past of $`x`$. (Of course, the whole discussion can be carried out as well for the “time-reversed” situation.) In general, since the wave equation gives rise to causal propagation (; see also Ref 14, p 250), the advanced Green function can be written as $$G(x,x^{})=\stackrel{~}{\mathrm{\Sigma }}(x,x^{})+\stackrel{~}{\mathrm{\Delta }}(x,x^{}),$$ $`(1.6)`$ where $`\stackrel{~}{\mathrm{\Sigma }}`$ is a distribution term with support on light-like separated pairs $`(x,x^{})`$, and $`\stackrel{~}{\mathrm{\Delta }}`$ has support on timelike related ones, with $`x`$ to the past of $`x^{}`$ in both cases (in fact, this is a possible definition of causal propagation). Thus, $`\stackrel{~}{\mathrm{\Delta }}`$ is of the form $$\stackrel{~}{\mathrm{\Delta }}(x,x^{})=:\mathrm{\Delta }(x,x^{})\theta _+(\mathrm{\Gamma }(x,x^{})),$$ $`(1.7)`$ where $`\mathrm{\Gamma }(x,x^{})`$ is the square of the proper distance calculated along the unique geodesic connecting $`x`$ and $`x^{}`$, and $`\theta _+`$, with a common abuse of notation, represents an “advanced step function,” which is zero when $`x`$ lies in the future of $`x^{}`$ . We will allow the Cauchy hypersurface $`𝒮`$ to be piecewise smooth, like, e.g., the future null cone of a point, which is not differentiable at the vertex. The field $`\varphi `$ at each $`xD^+(𝒮)J^+(𝒮)`$ can then be expressed as $$\varphi (x)=_𝒮dS_a(x^{})\left[g^{ab}(x^{})G(x^{},x)\begin{array}{c}\mathrm{}\mathrm{-}\mathrm{}\hfill \\ \hfill \end{array}_b^{}\varphi (x^{})+H^a(x^{})G(x^{},x)\varphi (x^{})\right],$$ $`(1.8)`$ where $`\mathrm{d}S_a(x^{})`$ is the oriented volume element on the hypersurface $`𝒮`$ at $`x^{}`$; this expression is valid even if $`𝒮`$ is a (partially) null hypersurface, provided that $`\mathrm{d}S_a`$ is appropriately defined (see section 5). We start in section 2 with a general discussion of the two-dimensional case. In section 3 we give some definitions of properties of wave propagation related to the absence of tails, whose meanings and mutual relationships are investigated in sections 4–7. Section 8 contains some concluding remarks, as well as speculations about possible lines of future research on the topic. As we have already been doing, we will use latin indices $`a`$, $`b`$, … as abstract indices in spacetime , which just indicate the tensorial nature of an object without requiring a set of coordinates, while greek indices $`\mu `$, $`\nu `$, … will be used for equations valid in some chart. The notations $`D^\pm (𝒜)`$, $`J^\pm (𝒜)`$, $`I^\pm (𝒜)`$, where $`𝒜`$ is some subset of $``$, stand, respectively, for its future/past domain of dependence, causal future/past, and chronological future/past , all defined in terms of the causal relations induced by $`g_{ab}`$. Minkowskian coordinates, in which the coefficients of the metric have (at least at a point) the form $`\eta _{\mu \nu }=\eta ^{\mu \nu }=\mathrm{diag}(1,1,\mathrm{},1)`$, will be denoted by $`(t,x^i)`$, with $`i=1,\mathrm{},m1`$. 2.The two-dimensional wave equation We begin with a general discussion of the wave equation (1.3) in a two-dimensional spacetime. The motivation for devoting an entire section to such a specific subject is twofold. First, it will provide us with a concrete example with which to illustrate some general ideas. Second, and more important, much of the paper will turn out to deal exclusively with this case; it is thus convenient to establish a few preliminary points, to which we can refer later on. Any two-dimensional spacetime is conformally flat, and its metric can therefore be locally written as $`g_{ab}=\mathrm{\Omega }^2\eta _{ab}`$, where $`\mathrm{\Omega }`$ is a nonvanishing function and $`\eta _{ab}`$ the Minkowski metric. In a general $`m`$-dimensional conformally flat spacetime we have $$_aX^a=_aX^a+_a\mathrm{ln}\|\mathrm{\Omega }\|^mX^a=\|\mathrm{\Omega }\|^m_a\left(\|\mathrm{\Omega }\|^mX^a\right),$$ $`(2.1)`$ for an arbitrary vector field $`X^a`$. In the particular case $`m=2`$, and for $`X^a=g^{ab}_bf=\mathrm{\Omega }^2\eta ^{ab}_bf`$, equation (2.1) allows us to write $$g^{ab}_a_bf=\mathrm{\Omega }^2\eta ^{ab}_a_bf,$$ $`(2.2)`$ for any function $`f`$. Therefore, in $`1+1`$ dimensions, $$g^{ab}_a_bf+H^a_af+Kf=\mathrm{\Omega }^2\left(\eta ^{ab}_a_bf+\overline{H}^a_af+\overline{K}f\right),$$ $`(2.3)`$ where $`\overline{H}^a:=\mathrm{\Omega }^2H^a`$ and $`\overline{K}:=\mathrm{\Omega }^2K`$. Since $`\mathrm{\Omega }0`$ everywhere on $``$, we have that $`\varphi `$ satisfies the wave equation (1.3) in $`(,g)`$ iff it satisfies $$\eta ^{ab}_a_b\varphi +\overline{H}^a_a\varphi +\overline{K}\varphi =0$$ $`(2.4)`$ in the flat spacetime $`(,\eta )`$. We can therefore restrict ourselves to studying (2.4) without loss of generality. It is convenient to introduce locally on $``$ the null coordinates $$\begin{array}{ccc}& u:=\frac{1}{2}(tx),\hfill & (2.5)\hfill \\ & v:=\frac{1}{2}(t+x),\hfill & (2.6)\hfill \end{array}$$ in which the Minkowski metric and its inverse have only one independent nonvanishing component each, $`\eta _{uv}=2`$ and $`\eta ^{uv}=1/2`$, respectively, so that (2.4) becomes $$_{uv}^2\varphi +U_u\varphi +V_v\varphi +W\varphi =0,$$ $`(2.7)`$ with $`U:=\overline{H}^u`$, $`V:=\overline{H}^v`$, and $`W:=\overline{K}`$. In the rest of the paper, we shall often refer to this form of the two-dimensional wave equation. It is not difficult to check, using (2.1) again, that for an arbitrary function $`f`$ in two dimensions, $$g^{ab}_a_bf_a\left(H^af\right)+Kf=\mathrm{\Omega }^2\left[_{uv}^2f_u\left(Uf\right)_v\left(Vf\right)+Wf\right].$$ $`(2.8)`$ Since $`\delta (u,v;u^{},v^{})=\mathrm{\Omega }^2\delta (uu^{})\delta (vv^{})`$, it follows that $`G(u,v;u^{},v^{})`$ is a Green function for (1.3) iff it satisfies the equation $$_{uv}^2G_u\left(UG\right)_v\left(VG\right)+WG=\delta (uu^{})\delta (vv^{}).$$ $`(2.9)`$ By direct substitution, we can verify that $$G(u,v;u^{},v^{})=\mathrm{\Delta }(u,v;u^{},v^{})\theta (u^{}u)\theta (v^{}v),$$ $`(2.10)`$ where $`\theta `$ is the step function, satisfies (2.9) provided that the following conditions are fulfilled: $$_{uv}^2\mathrm{\Delta }_u\left(U\mathrm{\Delta }\right)_v\left(V\mathrm{\Delta }\right)+W\mathrm{\Delta }=0;$$ $`(2.11)`$ $$\mathrm{\Delta }(u,v^{};u^{},v^{})=\mathrm{exp}\left(_u^u^{}du^{\prime \prime }V(u^{\prime \prime },v^{})\right);$$ $`(2.12)`$ $$\mathrm{\Delta }(u^{},v;u^{},v^{})=\mathrm{exp}\left(_v^v^{}dv^{\prime \prime }U(u^{},v^{\prime \prime })\right).$$ $`(2.13)`$ Since (2.12) and (2.13) can be regarded as data in a characteristic initial value problem for (2.11), which has a unique solution, it follows that $`\mathrm{\Delta }(u,v;u^{},v^{})`$ is completely determined by (2.11)–(2.13); therefore, (2.10) is the general form of the advanced Green function for the two-dimensional wave equation. In the particular case $`U=V=K=0`$, we recover the well-known result $`\mathrm{\Delta }=1`$ . 3.Definitions Various properties are used to characterise the absence of wave tails. First of all, we need to distinguish between the fact that a particular solution of the wave equation may not have tails, and a possible intrinsically non-diffusive nature of the equation itself. We say that a set of Cauchy data with compact support $`𝒞`$ produces no tails iff the field $`\varphi `$ obtained evolving these data vanishes in $`K^+(𝒞)`$, the set of all points in the causal future of $`𝒞`$ which cannot be reached from $`𝒞`$ by a null geodesic . For such data and points $`xK^+(𝒞)`$, (1.8) becomes $$\varphi (x)=_𝒞dS_a(x^{})\left[g^{ab}(x^{})\mathrm{\Delta }(x^{},x)\begin{array}{c}\mathrm{}\mathrm{-}\mathrm{}\hfill \\ \hfill \end{array}_b^{}\varphi (x^{})+H^a(x^{})\mathrm{\Delta }(x^{},x)\varphi (x^{})\right]\theta (\mathrm{\Gamma }(x^{},x)),$$ $`(3.1)`$ since the $`\stackrel{~}{\mathrm{\Sigma }}`$ term and the derivative of the $`\theta `$-function in $`G(x^{},x)`$ do not contribute to the field at $`x`$ (we retain the factor $`\theta (\mathrm{\Gamma })`$ in order to cover also cases in which $`xK^+(𝒞)`$ is not in the causal future of some $`x^{}𝒞`$, as it may happen, for example, when $`𝒞`$ is not simply connected). Therefore, the data on $`𝒞`$ produce no tails iff the integral on the right hand side of (3.1) vanishes. A related notion, expressed directly in terms of the field rather than of the Cauchy data, generalises the concept of distortionless propagation, without necessarily referring to a differential equation . We say that a wave propagates without distortion if it can be written as $`\varphi (x)=R(u(x))`$, where the phase $`u`$ is a function on $``$, and the wave form $`R`$ is a function of one variable. Clearly, if $`R`$ has compact support $`[u_1,u_2]`$, the wave propagates permanently sandwiched between those two values of $`u`$. We have here no metric with which to verify whether this wave has tails in the sense defined above, but it is nevertheless obvious that “there are no tails with respect to $`u`$.” A slightly more general situation is that of relatively undistorted or simple progressing waves, those of the form $`\varphi (x)=f(x)R(u(x))`$, where the amplitude $`f`$ is a function on $``$ \[5,6,16–19\]. An example in Minkowski spacetime is the spherical wave $`\varphi (t,𝐱)=\mathrm{exp}\mathrm{i}(k\|𝐱\|\omega t)/\|𝐱\|`$, which progresses with speed $`\omega /k`$ in the radial direction, and is undistorted except for the decrease in the amplitude $`1/\|𝐱\|`$. Progressing waves become interesting for suitable choices of the phase; the most common ones in Minkowski spacetime have $`u(t,𝐱)=tF(𝐱)`$, for some wave front $`F`$, function of the spatial coordinates . A finite sum of simple progressing waves will also be called (relatively) undistorted, the simplest example being the general solution $`\varphi (t,x)=R(tx)+S(t+x)`$ of the two-dimensional equation (2.7) with $`U=V=W=0`$. A useful further generalisation is that of progressing waves of order $`N`$, those that can be written as $$\varphi (x)=\underset{i=0}{\overset{N}{}}f_i(x)R^{(i)}(u(x)),$$ $`(3.2)`$ where the superscript $`(i)`$ denotes the $`i^{\mathrm{th}}`$ derivative. For these waves it is still true that a compact support for $`R`$ implies a “sandwich propagation” for $`\varphi `$. Turning now to intrinsic properties of the wave equation, corresponding to the absence of tails in the general solution, the most frequently used one by far is Huygens’ principle: The improper use of the term “principle” to denote what is actually only a property is commonly adopted for historical reasons, with the probably unique exception of Hadamard, who refers to it as “Huygens’ minor premise” . In reality, Huygens made use of TF below rather than HP, and merely in an approximate version, by postulating that almost all the waves emitted by a pointlike source are concentrated on the wavefront (see Ref 1, e.g., pp. 18 and 22). It is interesting to notice that this hypothesis does not correspond to assuming that tails are absent, but only that they are small! HP: A wave equation is said to satisfy the Huygens principle (HP) iff the field at a point $`x`$ depends only on the data at the intersection of the past light cone through $`x`$ with the Cauchy hypersurface \[4–6\], in the sense that any two sets of data which coincide there must give the same field at $`x`$. It is obvious from (1.8) that HP is equivalent to the requirement that the advanced Green function $`G(x,x^{})`$ have support only on pairs of points such that $`x`$ lies on the past light cone of $`x^{}`$. The latter statement is often also referred to as “Huygens’ principle”. An obvious consequence of (2.10) is that HP is always violated in a two-dimensional spacetime. The concepts discussed at the beginning of this section give rise, however, to other definitions that can be found in the literature. Using the notion of data that produce no tails, we define: TF: A wave equation is said to be tail-free (TF) iff each set of Cauchy data with compact support produces no tails . It is almost trivial to see that TF and HP are equivalent; nevertheless, for the sake of clarity, we shall present the explicit proof in the next section. It is sometimes interesting to consider data assigned on a characteristic (i.e., null ) hypersurface, and to give therefore a modified version of TF. This property was originally formulated in terms of solutions of a wave equation , but we find it more meaningful to state it as applied to the equation itself : CPP: A wave equation is said to satisfy the characteristic propagation property (CPP) iff each set of null data of compact support, which vanish at any singular points of the support, produces no tails. Why restrict, in this definition, the admissible data to those that vanish at singular points? For a two-dimensional characteristic initial value problem, we can consider without loss of generality the initial data $`\varphi (u,0)=:\phi (u)`$ and $`\varphi (0,v)=:\psi (v)`$ to be assigned on the union of the two “half-axes” $`u0`$ and $`v0`$. Consider the simplest example of (2.7), with $`U=V=W=0`$. Its general solution is of the form $`\varphi (u,v)=R(u)+S(v)`$, which in terms of the data becomes $`\varphi (u,v)=\phi (u)+\psi (v)\varphi (0,0)`$. This means that, if $`\varphi (0,0)0`$, no data with compact supports $`[u_1,u_2]`$ and $`[v_1,v_2]`$ can lead to “sandwich propagation,” i.e., one for which the support of the full solution is contained in the union of the strips $`u_1uu_2`$ and $`v_1vv_2`$. There is, of course, nothing bad in this, since we know already that HP is violated in two dimensions. However, data $`\phi (u)`$ and $`\psi (v)`$ such that $`\phi (0)=\psi (0)=\varphi (0,0)=0`$ do lead to sandwich propagation. The point $`(0,0)`$ here is an example of singular point of $`𝒮`$, and the reason we excluded data which do not vanish at such points in the definition of CPP is that this allows us to give a reasonable no-tails characterisation that may hold even in cases in which HP does not. Actually, we shall see in section 5 that CPP, although in principle weaker than HP and TF, in practice always implies them, except in two spacetime dimensions. We wish to stress that the restriction of data we are considering is actually a very small one, concerning a subset of $`𝒞`$ of measure zero which contains points that are already pathological. Furthermore, such a restriction corresponds precisely to what is commonly done when assigning data on the asymptotic past (probably the only characteristic initial value problem with good physical motivations), where one allows $`\varphi 0`$ at past null infinity $`^{}`$ but requires $`\varphi =0`$ at past timelike infinity $`i^{}`$ . In addition, the notion of progressing waves motivates the following definition for two-dimensional equations: PW$`_N`$: A wave equation in two spacetime dimensions is said to satisfy the progressing wave propagation property of order $`N`$ (PW$`_N`$) iff its general solution can be written as a sum of two progressing waves, $$\varphi (x)=\underset{i=0}{\overset{N}{}}f_i(x)R^{(i)}(u(x))+\underset{i=0}{\overset{N}{}}g_i(x)S^{(i)}(v(x)),$$ $`(3.3)`$ where the amplitudes $`f_i`$ and $`g_i`$ are fixed functions on $``$ depending on the wave equation (and at least one of $`f_N`$ and $`g_N`$ is not identically zero), while the wave forms $`R`$ and $`S`$ are arbitrary sufficiently differentiable functions of one variable (this forces the coordinates $`u`$ and $`v`$ to be null) . In section 6 we shall show that, in two spacetime dimensions, CPP is equivalent to PW$`_0`$; as far as the PW$`_N`$, with $`N>0`$, are concerned, their resemblance to a no-tails property is only apparent. Furthermore, wave equations in more than two dimensions may well have progressing wave solutions for appropriate choices of the wave front (see, e.g., Ref 16), but it seems excessive to ask that their general solution be expressible as a finite sum of progressing waves. These issues, which diminish the appeal of PW$`_N`$ as regards the study of tails, will be discussed in the concluding section. Another property that can be satisfied by two-dimensional wave equations is their solvability by the method of Kundt and Newman . To apply this method, we start by writing (2.7) in the two equivalent normal forms $$\begin{array}{ccc}& _v(j_0_u\varphi _0)j_1\varphi _0=0,\hfill & (3.4)\hfill \\ \multicolumn{3}{c}{}\\ & _u(l_0_v\psi _0)l_1\psi _0=0,\hfill & (3.5)\hfill \end{array}$$ where $`\varphi _0`$ and $`\psi _0`$ are related to $`\varphi `$ by factor transformations $`\varphi _0=\varphi \mathrm{exp}\sigma `$ and $`\psi _0=\varphi \mathrm{exp}\tau `$, with $$\sigma (u,v):=^udu^{}V(u^{},v),$$ $`(3.6)`$ $$\tau (u,v):=^vdv^{}U(u,v^{}),$$ $`(3.7)`$ while the coefficients are related to those in (2.7) by $$j_0=l_0^1=\mathrm{exp}(\tau \sigma ),$$ $`(3.8)`$ $$j_1=(_vV+UVW)j_0,$$ $`(3.9)`$ and $$l_1=(_uU+UVW)l_0.$$ $`(3.10)`$ If we inductively define $`j_k`$ and $`\varphi _k`$ by $$\begin{array}{ccc}& j_{k+1}/j_k=j_k/j_{k1}_{uv}^2\mathrm{ln}\|j_k\|,\hfill & (3.11)\hfill \\ \multicolumn{3}{c}{}\\ & j_{k+1}\varphi _{k+1}=j_k_u\varphi _k,\hfill & (3.12)\hfill \end{array}$$ assuming of course $`j_k0`$ for all $`k\text{}`$, we obtain from (3.4) a countable set of equations $$_v(j_k_u\varphi _k)j_{k+1}\varphi _k=0,k\text{}.$$ $`(3.13)`$ Similarly, if we inductively define $`l_k`$ and $`\psi _k`$ by $$\begin{array}{ccc}& l_{k1}/l_k=l_k/l_{k+1}_{vu}^2\mathrm{ln}\|l_k\|,\hfill & (3.14)\hfill \\ \multicolumn{3}{c}{}\\ & l_{k1}\psi _{k1}=l_k_v\psi _k,\hfill & (3.15)\hfill \end{array}$$ assuming now that $`l_k0`$ for all $`k\text{}`$, we obtain from (3.5) a second countable set $$_u(l_k_v\psi _k)l_{k1}\psi _k=0,k\text{}.$$ $`(3.16)`$ We shall refer to equations (3.13) and (3.16) as being in the $`k^{\mathrm{th}}`$ $`v`$\- and $`u`$-normal form, respectively. It is not hard to check that for all $`k\text{}`$ the $`k^{\mathrm{th}}`$ $`v`$-normal form equation, corresponding to the coefficients $`j_k`$ and $`j_{k+1}`$, and the $`k^{\mathrm{th}}`$ $`u`$-normal form equation, corresponding to the coefficients $`l_k`$ and $`l_{k1}`$, are equivalent under the transformation $$\begin{array}{ccc}& j_kl_k=1,\hfill & (3.17)\hfill \\ \multicolumn{3}{c}{}\\ & \varphi _k=l_k\psi _k,\hfill & (3.18)\hfill \end{array}$$ for $`k\text{}`$. It is also easy to see that the equations within the set (3.13) (respectively, (3.16)) are equivalent in the sense that a solution of any one of them generates a solution of every other one of them through (3.11)–(3.13), (3.17), and (3.18) (respectively, (3.14)–(3.18)), and we thus obtain two equivalence classes of $`v`$\- and $`u`$-normal form equations labelled by the index $`k`$ ranging over both negative and positive integers. Given a wave equation in the $`0^{\mathrm{th}}`$ $`v`$-normal form (3.4), we say that its substitution sequence $`\{j_k\}`$ is double terminating in $`N`$ steps when $`j_{k_1+1}=0`$ and $`l_{k_21}=0`$ (but $`j_{k_1}0`$ and $`l_{k_2}0`$) for some $`k_10`$, $`k_20`$, and $`N=\mathrm{max}\{k_1,k_2\}`$. However, it was shown in Ref 20, using (3.11)–(3.18), that in this case the general solution of (3.4) is $`\varphi _0=\varphi _A+l_0\varphi _R`$, where $$\begin{array}{ccc}\hfill \varphi _A& :=\frac{1}{j_1}_v\left(\frac{j_1}{j_2}_v\left(\frac{j_2}{j_3}_v\left(\mathrm{}\frac{j_{k_11}}{j_{k_1}}_v\left(j_{k_1}S(v)\right)\mathrm{}\right)\right)\right),\hfill & (3.19)\hfill \\ \hfill \varphi _R& :=\frac{1}{l_1}_u\left(\frac{l_1}{l_2}_u\left(\frac{l_2}{l_3}_u\left(\mathrm{}\frac{l_{k_2+1}}{l_{k_2}}_u\left(l_{k_2}R(u)\right)\mathrm{}\right)\right)\right),\hfill & (3.20)\hfill \end{array}$$ with $`S(v)`$ and $`R(u)`$ arbitrary functions of one variable. It is obvious from (3.19) and (3.20) that such a $`\varphi `$ is a progressing wave of order $`N`$. This relationship motivates the following definition: KN$`_N`$: A wave equation in two spacetime dimensions is said to be solvable by the Kundt-Newman method in $`N`$ steps (KN$`_N`$) iff its substitution sequence is double terminating in $`N`$ steps. As we have just seen, all KN$`_N`$ wave equations are PW$`_N`$. We shall see in section 7 that the converse is also true, so that KN$`_N`$ and PW$`_N`$ are equivalent properties. 4.Equivalence between HP and TF We now begin studying the relationships between the various properties of wave equations we listed in the previous section. The first three are related in a simple way in any number of dimensions; as we will now show, TF and HP are always trivially equivalent, which justifies the fact that they are often identified, and they are almost always equivalent to CPP (see next section). These results generalise to arbitrary wave equations of the type (1.3) those of Ref 15. In order to show that TF $``$ HP, let us consider any point $`xD^+(𝒮)`$, and assign data with support $`𝒞I^{}(x)𝒮`$. Since TF holds by hypothesis, and $`xK^+(𝒞)`$, it follows that $`\varphi (x)=0`$ for each set of data prescribed on such a $`𝒞`$. By causality, $`\varphi (x)`$ cannot be influenced by data given outside $`I^{}(x)𝒮`$, hence it can depend only on their value on $`E^{}(x)𝒮`$, where $`E^{}(x)`$ stands for the set of points in $`J^{}(x)`$ that are null related to $`x`$. This is precisely the content of HP. Let us now prove that HP $``$ TF. First of all, let us notice that, by choosing data that vanish everywhere on $`𝒮`$, we get $`\varphi (x)=0`$ everywhere in $``$ by (1.3). If we now assign data on $`𝒮`$ with compact support $`𝒞`$, HP implies that the value of $`\varphi (x)`$, for any $`xK^+(𝒞)`$, must be independent of the data; in particular, it must have the same value than in the case of vanishing data, i.e., it must be equal to zero, by the remark above. This completes the proof that HP $``$ TF. It is interesting to notice that the structure of this proof allows one to extend it to more general wave equations than (1.3). Actually, in the case of the latter, a simpler proof can be given that makes use of the equivalence between HP and the property $`\stackrel{~}{\mathrm{\Delta }}(x^{},x)0`$. In fact, it follows immediately from (3.1) that, if HP is satisfied, then $`\varphi (x)=0`$ for all $`xK^+(𝒞)`$, i.e., the wave equation is TF. The converse is also true, as we can see by choosing $`𝒮`$ to be spacelike at a point $`\overline{x}`$, and data of support only at $`\overline{x}`$, such that, denoting by $`n^a`$ the unit vector normal to $`𝒮`$, $$\begin{array}{ccc}& \varphi \|_𝒮(x)0,\hfill & (4.1)\hfill \\ \multicolumn{3}{c}{}\\ & (n^a_a\varphi )\|_𝒮(x)=\delta _𝒮(x,\overline{x}),\hfill & (4.2)\hfill \end{array}$$ where $`\delta _𝒮`$ is the $`(m1)`$-dimensional delta function on $`𝒮`$. Then (3.1) reduces to $$\varphi (x)=\mathrm{\Delta }(\overline{x},x),$$ $`(4.3)`$ and if $`\varphi (x)=0`$ for all $`xK^+(𝒞)=I^+(\overline{x})`$, we must have $`\mathrm{\Delta }(\overline{x},x)=0`$ for all such pairs of points, i.e., $`\stackrel{~}{\mathrm{\Delta }}(x^{},x)0`$. However, this proof relies heavily on the existence of the integral representation (1.8), and appears much more difficult to generalise to the case of more complicated wave equations than the “geometrical” one given above. 5.Relationship between HP and CPP It is obvious from the definitions given in section 3 that TF, and thus HP, imply CPP. To check whether the converse holds we need to consider again the integral expression (3.1) for $`xK^+(𝒞)`$, where $`𝒞`$ is now part of a piecewise null hypersurface $`𝒮`$ locally defined by the condition $`w=\mathrm{const}`$, for some $`w:\text{}`$ with $`g^{ab}_aw_bw=0`$ at all points where $`𝒮`$ is differentiable (this condition may fail only in a subset of $`𝒮`$ of measure zero). One usually writes $`\mathrm{d}S_a=\mathrm{d}Sn_a`$, where $`\mathrm{d}S`$ and $`n_a`$ are often thought of as volume element and unit normal to the surface $`𝒮`$, respectively, which clearly does not make sense if $`𝒮`$ is null. However, the definitions of these objects can be extended to that case as well. Let us introduce local coordinates $`\{\xi ^i:i=1,\mathrm{},m1\}`$ on $`𝒮`$, and use $`w`$ as a coordinate transverse to $`𝒮`$, so that $`(w,\xi ^i)`$ are coordinates on $``$ adapted to $`𝒮`$. Then we can always write $`\mathrm{d}S_\mu =\mathrm{d}^{m1}\xi \sqrt{g}\delta _{\mu }^{}{}_{}{}^{w}`$, with the understanding that $`g`$ must be calculated in these coordinates, which in the null case can be split into $`\mathrm{d}S=\mathrm{d}^{m1}\xi \sqrt{g}`$, and $`n_\mu =\delta _{\mu }^{}{}_{}{}^{w}`$; the latter are the components of the form $`n_a=_aw`$. With these conventions, we can write the Gauss theorem for an arbitrary vector field $`Y^a`$ as $$_𝒩dV_aY^a=_𝒩dS_aY^a,$$ $`(5.1)`$ where $`𝒩`$ is a region in $``$, and $`\mathrm{d}V=\mathrm{d}^mx\sqrt{g}`$ is the spacetime volume element, even when (part) of $`𝒩`$ is null. This follows from the fact that the left hand side of (5.1) can be broken into a sum of integrals of the form $$_𝒟\mathrm{d}^{m1}\xi dw_\mu \left(\sqrt{g}Y^\mu \right),$$ $`(5.2)`$ where $`𝒟`$ is a domain in $`\text{}^m`$ one of whose “faces” (that we denote by $``$) is defined by $`w=0`$; but (5.2) can be transformed, using the Gauss theorem in $`\text{}^m`$, into a sum of terms corresponding to the various faces of $`𝒟`$, of which only $$_{}\mathrm{d}^{m1}\xi \sqrt{g}\delta _{\mu }^{}{}_{}{}^{w}Y^\mu $$ $`(5.3)`$ contributes, in the end, to the total expression. Therefore, (5.1) (hence (1.8) and (3.1)) is completely justified, independently of the causal character of $`𝒩`$. We can rewrite (3.1) in the coordinates $`(w,\xi ^i)`$ as $$\varphi (x)=_𝒞dS^{}\left\{_\mu ^{}(n^\mu \mathrm{\Delta }\varphi )\varphi \left[(_\mu ^{}n^\mu )\mathrm{\Delta }+n^\mu (2_\mu ^{}\mathrm{\Delta }H_\mu \mathrm{\Delta })\right]\right\}\theta (\mathrm{\Gamma })$$ $`(5.4)`$ (here and in the following integrals, all functions in the integrand depend on $`x^{}`$, and all two-point functions on the ordered pair $`(x^{},x)`$). Since we have $`n^\mu =g^{\mu w}`$, which implies $`n^w=g^{ww}=g^{\mu \nu }_\mu w_\nu w=0`$, the first term in the right hand side of (5.4) contains only derivatives performed along $`𝒮`$, and can be transformed as $$_𝒞dS^{}_\mu ^{}\left(n^\mu \mathrm{\Delta }\varphi \right)=_{}\mathrm{d}^{m1}\xi ^{}_i^{}\left(\sqrt{g}n^i\mathrm{\Delta }\varphi \right)=_𝒞d\sigma _a^{}n^a\mathrm{\Delta }\varphi ,$$ $`(5.5)`$ where $``$ stands again for the region of $`\text{}^m`$ defined by $`w=0`$, $`\mathrm{d}\sigma _a`$ is the oriented volume element on $`𝒞`$, and the Gauss theorem in $`\text{}^{m1}`$ has been applied in the last step (possible discontinuities of $`\varphi `$ do not prevent us from using the theorem, if they are correctly interpreted in the sense of distributions). We can therefore deduce that this term does not contribute to the expression, because $`\varphi (x^{})=0`$ outside $`𝒞`$ and we can think of the integration as being performed over a region larger than $`𝒞`$, on whose boundary $`\varphi (x^{})=0`$. Therefore, CPP is satisfied iff $$_𝒞dS^{}\varphi \left[(_\mu ^{}n^\mu )\mathrm{\Delta }+n^\mu (2_\mu ^{}\mathrm{\Delta }H_\mu \mathrm{\Delta })\right]\theta (\mathrm{\Gamma })=0$$ $`(5.6)`$ for all data on an arbitrary compact null $`𝒞`$. In particular, we can choose as data $$\varphi \|_𝒮(x)=\delta _𝒮(x,\overline{x}),$$ $`(5.7)`$ for any nonsingular point $`\overline{x}𝒞`$. Since $`\overline{x}`$ is arbitrary, this implies that either $`\mathrm{\Delta }(x^{},x)=0`$, and we get HP directly, or $$_an^a(x)+n^a(x)\left[2_a\mathrm{ln}\|\mathrm{\Delta }(x,x^{})\|H_a(x)\right]=0$$ $`(5.8)`$ for all pairs of points with $`x^{}I^+(x)`$, $`x𝒞`$ nonsingular, and all null vector fields $`n^a`$ which are gradients of a function. We will now study the consequences of (5.8), by writing it in a suitable coordinate system. In terms of $`w`$, (5.8) is of the form $$g^{ab}_a_bw+X^a_aw=0,$$ $`(5.9)`$ where $`X^a(x,x^{}):=2^a\mathrm{ln}\|\mathrm{\Delta }(x,x^{})\|H^a(x)`$. Choose any point $`x_0𝒞`$, and a system of Riemann normal coordinates based at $`x_0`$, for which $`g_{\mu \nu }(x)=\eta _{\mu \nu }+𝒪(2)`$, where we have introduced the notation $`𝒪(n)`$ for a quantity which is of order $`n`$ in the coordinate separation of $`x`$ from $`x_0`$; then (5.9) becomes $$\eta ^{\mu \nu }_\mu _\nu w+X^\mu _\mu w=𝒪(1).$$ $`(5.10)`$ Let us now consider the function $$w=t+\sqrt{x^ix_i}+c_{\mu \nu }x^\mu x^\nu +𝒪(3),$$ $`(5.11)`$ where the $`c_{\mu \nu }`$ are appropriate coefficients such that $`g^{\mu \nu }_\mu w_\nu w=0`$. Substituting into (5.10) we get that, for $`x^ix_i0`$, $$\frac{m2}{\sqrt{x^ix_i}}+2\eta ^{\mu \nu }c_{\mu \nu }+X^t(x,x^{})+\frac{\eta _{ij}X^i(x,x^{})x^j}{\sqrt{x_kx^k}}=𝒪(1).$$ $`(5.12)`$ The first term in the left hand side of (5.12) is $`𝒪(1)`$, whereas the others are all $`𝒪(0)`$; therefore (5.12) implies $`m=2`$ and $$X^t(x_0,x^{})+X^x(x_0,x^{})\underset{x0^\pm }{lim}\frac{x}{\|x\|}=2\eta ^{\mu \nu }c_{\mu \nu },$$ $`(5.13)`$ where the $`\pm `$ in the limit corresponds to the arbitrariness in the direction of approach of $`x^\mu `$ to the origin and, with some abuse of notation, we have denoted by $`x`$ the spatial coordinate of the two-dimensional Minkowski spacetime. This means that $$X^t(x_0,x^{})\pm X^x(x_0,x^{})=2\eta ^{\mu \nu }c_{\mu \nu }$$ $`(5.14)`$ must hold simultaneously for both signs, i.e., that $`X^x(x_0,x^{})=0`$. Since this is true in any system of local Minkowskian coordinates, and taking into account the arbitrariness of $`x_0`$, we deduce that $`X^a(x,x^{})=0`$ as a field. In conclusion, we have found that CPP implies either HP, or $`m=2`$ and $$H_a(x)=2_a\mathrm{ln}\|\mathrm{\Delta }(x,x^{})\|.$$ $`(5.15)`$ In other words, CPP is always equivalent, for $`m>2`$, to HP; in two spacetime dimensions $`\mathrm{\Delta }0`$ and HP is always violated, whereas CPP can be satisfied. Therefore, CPP can be regarded as a nontrivial generalisation of HP, although it becomes interesting by itself only for $`m=2`$. More explicitly, the definitions of CPP and HP differed in that for the former, the support of the data had to be null rather than achronal, and the data had to vanish at non-smooth points of $`𝒞`$. We now know that this is of no consequence other than in some two-dimensional cases, but it is perhaps surprising that such a small restriction in the data can have the effect of enlarging, in a nontrivial way, the class of equations to which the definition applies. We have derived (5.15) above as a necessary condition for the validity of CPP in two dimensions, but the following simple argument shows that it is also sufficient. Substituting (5.15) into (5.9) the latter reduces to $`g^{ab}_a_bw=0`$, and CPP holds if this equation is satisfied by any function $`w`$ that generates null hypersurfaces. For $`m=2`$ this is true, because $`w`$ is null iff it depends on only one of the coordinates $`u`$ and $`v`$ defined by (2.5) and (2.6), i.e., iff either $`w=w_1(u)`$ or $`w=w_2(v)`$. Using the identity (2.2) for $`f=w`$, we have $`g^{ab}_a_bw=0`$. Let us now see what restrictions the condition (5.15) places on the wave equation. First of all, since the Riemannian connection is torsion-free, we have $`_{[a}_{b]}\mathrm{ln}\|\mathrm{\Delta }\|=0`$, and (5.15) implies $$_{[a}H_{b]}=0;$$ $`(5.16)`$ this is an integrability condition that allows us to obtain $`\mathrm{\Delta }`$ by straightforward integration of (5.15). In fact, (5.16) implies that $`H_a=2_a\mathrm{\Lambda }`$, for some $`\mathrm{\Lambda }`$; substituting back into (5.15), we have $$\|\mathrm{\Delta }(x,x^{})\|=\frac{\mathrm{exp}\mathrm{\Lambda }(x)}{\mathrm{exp}\mathrm{\Lambda }(x^{})},$$ $`(5.17)`$ since $`\mathrm{\Delta }(x^{},x^{})=1`$ by (2.9) and (2.10). Furthermore, it is not difficult to verify, from the definition of the Green function, that $`\mathrm{\Delta }(x,x^{})`$ satisfies the “adjoint wave equation” in $`x`$ (see also (2.8) and (2.11)) $$g^{ab}_a_b\mathrm{\Delta }_a\left(H^a\mathrm{\Delta }\right)+K\mathrm{\Delta }=0.$$ $`(5.18)`$ Substituting (5.15) twice into (5.18), we get $$2_aH^a+g_{ab}H^aH^b4K=0,$$ $`(5.19)`$ hence $$g^{ab}_a_b\mathrm{\Lambda }+g^{ab}_a\mathrm{\Lambda }_b\mathrm{\Lambda }K=0.$$ $`(5.20)`$ Using (5.17), we can replace (5.20) by an equation for $`\mathrm{\Delta }`$ simpler than (5.18), $$g^{ab}_a_b\mathrm{\Delta }K\mathrm{\Delta }=0,$$ $`(5.21)`$ which can also be directly obtained by using (5.19) into (5.18). Equations (5.16) and (5.20) (or (5.19)) are also sufficient conditions for CPP, as one can easily see by considering a new field $`\psi :=\varphi \mathrm{exp}\mathrm{\Lambda }`$. From (1.3) and (5.20), we find that $`\psi `$ obeys the equation $`g^{ab}_a_b\psi =0`$, which is CPP because $`m=2`$. Since $`\psi `$ and $`\varphi `$ differ only by the nonvanishing factor $`\mathrm{exp}\mathrm{\Lambda }`$, it follows that $`\varphi `$ satisfies CPP as well. In the coordinates $`(u,v)`$ we have $`H^u=\mathrm{\Omega }^2U`$, $`H^v=\mathrm{\Omega }^2V`$, $`K=\mathrm{\Omega }^2W`$, so that (5.16) and (5.19) are equivalent to $$_uU=_vV$$ $`(5.22)`$ and $$_uU+UVW=0,$$ $`(5.23)`$ respectively. Similarly, we have $`U=_v\mathrm{\Lambda }`$ and $`V=_u\mathrm{\Lambda }`$; (5.20) and (5.21) become $$_{uv}^2\mathrm{\Lambda }+_u\mathrm{\Lambda }_v\mathrm{\Lambda }W=0$$ $`(5.24)`$ and $$_{uv}^2\mathrm{\Delta }W\mathrm{\Delta }=0.$$ $`(5.25)`$ These conditions for CPP represent strong constraints on the wave equation; for example, when $`U=V=0`$, (5.23) and (5.24) are satisfied only if $`W=0`$ . Notice that when (5.16), or equivalently (5.22), is satisfied, the functions $`\sigma `$ and $`\tau `$ defined in (3.6) and (3.7) both reduce to $`\mathrm{\Lambda }`$; upon substitution into (3.8)–(3.10), equations (5.22) and (5.23) then imply that $`j_0=l_0=1`$ and $`j_1=l_1=0`$. 6.Equivalence between CPP and PW$`_0`$ in two spacetime dimensions In this section we will show that, in two spacetime dimensions, CPP and PW$`_0`$ are equivalent properties; the proof is very simple if we rely on some results of the previous section. In fact, since the two-dimensional wave equation (2.7) satisfies CPP iff $`\varphi `$ can be written as $`\varphi =\mathrm{exp}(\mathrm{\Lambda })\psi `$, with $`_{uv}^2\psi =0`$, it follows that (2.7) is CPP iff its general solution has the form $$\varphi (u,v)=\mathrm{e}^{\mathrm{\Lambda }(u,v)}\left(r(u)+s(v)\right),$$ $`(6.1)`$ where $`r`$ and $`s`$ are arbitrary functions. By comparing (6.1) with the $`N=0`$ case of (3.3), $$\varphi (u,v)=f_0(u,v)R(u)+g_0(u,v)S(v),$$ $`(6.2)`$ we conclude immediately that CPP $``$ PW$`_0`$. In order to prove the converse, i.e., that PW$`{}_{0}{}^{}`$ CPP, substitute the expression (6.2) into (2.7); invoking the arbitrariness of $`R`$ and $`S`$, we get: $$\begin{array}{ccc}\hfill _{uv}^2f_0+U_uf_0& +V_vf_0+Wf_0=0;\hfill & (6.3)\hfill \\ \multicolumn{3}{c}{}\\ \hfill _{uv}^2g_0+U_ug_0& +V_vg_0+Wg_0=0;\hfill & (6.4)\hfill \\ \multicolumn{3}{c}{}\\ \hfill _vf_0& +Uf_0=0;\hfill & (6.5)\hfill \\ \multicolumn{3}{c}{}\\ \hfill _ug_0& +Vg_0=0.\hfill & (6.6)\hfill \end{array}$$ After some manipulations, these equations lead to (5.22) and (5.23), which are sufficient conditions for the validity of CPP, as we saw in section 5. We can therefore conclude that, in a two-dimensional spacetime, CPP and PW$`_0`$ are equivalent. An immediate corollary of this result is that not every expression of the type (6.2) is the general solution of a two-dimensional wave equation. This is evident by a comparison between (6.1) and (6.2), which shows that the amplitudes $`f_0`$ and $`g_0`$ must be related by $`f_0(u,v)/g_0(u,v)=\alpha (u)/\beta (v)`$, for some functions $`\alpha `$ and $`\beta `$. The same conclusion can be reached in a more explicit way by integration of (6.5) and (6.6), to obtain $$\begin{array}{ccc}\hfill f_0(u,v)& =f_0(u,0)\mathrm{exp}\left(_0^vdv^{}U(u,v^{})\right)=\alpha (u)\mathrm{exp}\left(\mathrm{\Lambda }(u,v)\right),\hfill & (6.7)\hfill \\ \hfill g_0(u,v)& =g_0(0,v)\mathrm{exp}\left(_0^udu^{}V(u^{},v)\right)=\beta (v)\mathrm{exp}\left(\mathrm{\Lambda }(u,v)\right),\hfill & (6.8)\hfill \end{array}$$ with $`\alpha (u):=f_0(u,0)\mathrm{exp}\mathrm{\Lambda }(u,0)`$ and $`\beta (v):=g_0(0,v)\mathrm{exp}\mathrm{\Lambda }(0,v)`$. By a generalisation of this argument, one could show that not every expression of the type (3.3), with arbitrary $`R`$ and $`S`$, is the general solution of a two-dimensional wave equation. It is possible to prove that CPP $``$ PW$`_0`$ without necessarily using the results of the previous section. This is obvious for what concerns the implication PW$`_0`$ $``$ CPP, since PW$`_0`$ leads us to the same equations (5.22) and (5.23) of section 5, and one can proceed to show that they are sufficient conditions for CPP exactly in the same way as we did there. The proof that CPP $``$ PW$`_0`$ is less compact, but we nevertheless present it because of its intrinsic interest. Consider the null data of compact support given by $`\phi (u)=\delta (uu_0)`$, with $`u_0>0`$, and $`\psi (v)0`$; then CPP implies that $`\varphi (u,v)`$ must be concentrated on the line $`u=u_0`$. This means that, for any fixed $`v`$, $`\varphi (u,v)`$ must be a distribution in $`u`$ that, applied to a test function $`F(u,v)`$, produces a number depending only on the behaviour of $`F`$ at the point $`(u_0,v)`$, i.e., on $`F^{(i)}(u_0,v)`$, with $`i0`$. Since distributions are linear functionals, the combination of the various $`F^{(i)}(u_0,v)`$ must be linear, so that we can write $$\varphi (u,v)=\underset{i=0}{\overset{+\mathrm{}}{}}f_i(u_0,v)\delta ^{(i)}(uu_0),$$ $`(6.9)`$ for some functions $`f_i`$. The sum on the right hand side of (6.9) actually consists only of a finite number of terms; this follows from a rigorous result of distribution theory , but can be heuristically justified as follows. Since the action of $`\varphi `$ as a functional must be defined on any test function $`F`$, the value of the series $$\underset{i=0}{\overset{+\mathrm{}}{}}f_i(u_0,v)F^{(i)}(u_0,v)$$ $`(6.10)`$ must be finite for each $`F`$. Suppose that the sum in (6.9) and (6.10) is infinite, that is, that there exist infinitely many $`f_i0`$. Then, taking $`F`$ such that $`F^{(i)}(u_0,v)=f_i(u_0,v)^1`$ when $`f_i0`$, and arbitrary otherwise, the expression (6.10) becomes infinite, which contradicts the hypothesis of regularity of $`\varphi `$. Hence, we must have $$\varphi (u,v)=\underset{i=0}{\overset{N_u}{}}f_i(u_0,v)\delta ^{(i)}(uu_0),$$ $`(6.11)`$ with $`N_u`$ finite. Similarly, if $`\phi (u)0`$ and $`\psi (v)=\delta (vv_0)`$, we have $$\varphi (u,v)=\underset{i=0}{\overset{N_v}{}}g_i(u,v_0)\delta ^{(i)}(vv_0),$$ $`(6.12)`$ for some finite $`N_v`$ and some functions $`g_i`$. From (6.11) and (6.12) it already follows that (2.7) is PW$`_N`$, with $`N=\mathrm{max}\{N_u,N_v\}`$, as can be seen just by using linearity and representing arbitrary data $`\phi `$ and $`\psi `$ as superpositions of $`\delta `$-like data, $$\begin{array}{ccc}& \phi (u)=du_0\phi (u_0)\delta (uu_0),\hfill & (6.13)\hfill \\ & \psi (v)=dv_0\psi (v_0)\delta (vv_0),\hfill & (6.14)\hfill \end{array}$$ but without making any other use of the differential equation for $`\varphi `$. We shall now show, however, that (2.7) allows us to set $`N=0`$, i.e., to conclude that CPP $``$ PW$`_0`$. For this purpose, let us first consider again data $`\phi (u)=\delta (uu_0)`$, $`\psi (v)0`$. Substituting the solution (6.11) into (2.7) we get $$\begin{array}{ccc}& \left[_vf_{N_u}(u_0,v)+U(u,v)f_{N_u}(u_0,v)\right]\delta ^{(N_u+1)}(uu_0)+\hfill & \\ & +\underset{i=1}{\overset{N_u}{}}[V(u,v)_vf_i(u_0,v)+W(u,v)f_i(u_0,v)+\hfill & \\ & +_vf_{i1}(u_0,v)+U(u,v)f_{i1}(u_0,v)]\delta ^{(i)}(uu_0)+\hfill & \\ & +\left[V(u_0,v)_vf_0(u_0,v)+W(u_0,v)f_0(u_0,v)\right]\delta (uu_0)=0.\hfill & (6.15)\hfill \end{array}$$ Smoothing (6.15) with a test function, and taking into account the arbitrariness of the latter, we obtain the following set of equations: $$\begin{array}{ccc}& _vf_{N_u}(u_0,v)+U(u_0,v)f_{N_u}(u_0,v)=0;\hfill & (6.16)\hfill \\ \multicolumn{3}{c}{}\\ & \underset{i=k}{\overset{N_u}{}}(1)^i\left(\genfrac{}{}{0pt}{}{i}{k}\right)[_u^{ik}V(u_0,v)_vf_i(u_0,v)+_u^{ik}W(u_0,v)f_i(u_0,v)+\hfill & \\ & +_u^{ik}U(u_0,v)f_{i1}(u_0,v)]+(1)^k_vf_{k1}(u_0,v)+\hfill & \\ & +(1)^{N_u+1}\left(\genfrac{}{}{0pt}{}{N_u+1}{k}\right)_u^{N_uk+1}U(u_0,v)f_{N_u}(u_0,v)=0;N_uk1;\hfill & (6.17)\hfill \\ \multicolumn{3}{c}{}\\ & \underset{i=1}{\overset{N_u}{}}(1)^i\left[_u^iV(u_0,v)_vf_i(u_0,v)+_u^iW(u_0,v)f_i(u_0,v)+_u^iU(u_0,v)f_{i1}(u_0,v)\right]+\hfill & \\ & +V(u_0,v)_vf_0(u_0,v)+W(u_0,v)f_0(u_0,v)=0+\hfill & \\ & +(1)^{N_u+1}_u^{N_u+1}U(u_0,v)f_{N_u}(u_0,v).\hfill & (6.18)\hfill \end{array}$$ From (6.16) and the fact that $`f_{N_u}(u_0,0)=0`$, which follows from using our data in (6.11), we obtain that $`f_{N_u}(u_0,v)=0`$ for all $`v`$. Using this in (6.17) and repeating the procedure, we obtain $`f_i(u_0,v)=0`$, $`i1`$, so that only $`f_0(u_0,v)`$ can be nonzero. An analogous argument, using data $`\phi (u)0`$, $`\psi (v)=\delta (vv_0)`$, shows that the only nonvanishing coefficient in (6.12) can be $`g_0(u,v_0)`$. Therefore, the general solution, obtained from (6.11) and (6.12) using arbitrary initial data as in (6.13) and (6.14), is $$\varphi (u,v)=f_0(u,v)\phi (u)+g_0(u,v)\psi (v),$$ $`(6.19)`$ from which the validity of PW$`_0`$ follows. Furthermore, it is not difficult to see that (6.17) for $`k=1`$ reduces to (6.5) in $`u=u_0`$, and that substituting this relation into (6.18) we obtain (5.23). By a completely symmetric procedure we can recover (6.6). 7.Equivalence between PW$`_N`$ and KN$`_N`$ As we saw in section 3, all KN$`_N`$ wave equations are PW$`_N`$; it has been conjectured that the converse is also true, namely that all PW$`_N`$ equations can in fact be solved exactly by the Kundt-Newman method in $`N`$ steps . To show that this is indeed the case, suppose we have a 2-dimensional wave equation whose general solution is the progressing wave (3.3). By suitable transformations, we can always write the wave equation in its $`0^{\mathrm{th}}`$ $`v`$\- or $`u`$-normal form, (3.4) or (3.5); let us concentrate on the $`v`$-normal form. Applying the differential operator in (3.4) to the progressing wave (3.3), and imposing that the coefficients of derivatives of the arbitrary functions $`R(u)`$ and $`S(v)`$ of different orders vanish separately in the resulting expression, we have: $$\begin{array}{ccc}& _v(j_0_uf_0)j_1f_0=0,\hfill & (7.1)\hfill \\ \multicolumn{3}{c}{}\\ & _v(j_0_uf_i)j_1f_i+_v(j_0f_{i1})=0,1iN,\hfill & (7.2)\hfill \\ \multicolumn{3}{c}{}\\ & _v(j_0f_N)=0,\hfill & (7.3)\hfill \end{array}$$ and $$\begin{array}{ccc}& _v(j_0_ug_0)j_1g_0=0.\hfill & (7.4)\hfill \\ \multicolumn{3}{c}{}\\ & _v(j_0_ug_i)j_1g_i+j_0_ug_{i1}=0,1iN,\hfill & (7.5)\hfill \\ \multicolumn{3}{c}{}\\ & _ug_N=0,\hfill & (7.6)\hfill \end{array}$$ Let us consider the group (7.4)–(7.6). Equations (7.4) and (7.5) can be rewritten more conveniently as $$_v\left(\frac{j_0}{j_1}_ug_0\right)+_v\mathrm{ln}\|j_1\|\frac{j_0}{j_1}_ug_0g_0=0$$ $`(7.7)`$ and $$_v\left(\frac{j_0}{j_1}_ug_i\right)+_v\mathrm{ln}\|j_1\|\frac{j_0}{j_1}_ug_ig_i+\frac{j_0}{j_1}_ug_{i1}=0,1iN,$$ $`(7.8)`$ respectively. By repeated differentiation of (7.8) with respect to $`u`$, one finds that, for any $`h1`$ and $`k=0,\mathrm{},N1`$, $$_v𝒟_u^hg_{Nk}+_v\mathrm{ln}\|j_h\|𝒟_u^hg_{Nk}𝒟_u^{h1}g_{Nk}+𝒟_u^hg_{Nk1}=0,$$ $`(7.9)`$ where the differential operator $`𝒟_u^k`$, containing $`k`$ derivatives $`_u`$, is defined as $$𝒟_u^kF:=\frac{j_{k1}}{j_k}_u\left(\frac{j_{k2}}{j_{k1}}_u\left(\mathrm{}\frac{j_1}{j_2}_u\left(\frac{j_0}{j_1}_uF\right)\mathrm{}\right)\right),$$ $`(7.10)`$ for any sufficiently differentiable function $`F`$. We shall prove (7.9) by induction over $`h`$. Assuming that it holds for some $`h`$, and differentiating it with respect to $`u`$, we have $$_{uv}^2𝒟_u^hg_{Nk}+_v\mathrm{ln}\|j_h\|_u𝒟_u^hg_{Nk}\frac{j_{h+1}}{j_h}𝒟_u^hg_{Nk}+_u𝒟_u^hg_{Nk1}=0,$$ $`(7.11)`$ where we have used (3.11). Rewriting the first term as $$_v\left(\frac{j_{h+1}}{j_h}\frac{j_h}{j_{h+1}}_u𝒟_u^hg_{Nk}\right)=_v\mathrm{ln}\left\|\frac{j_{h+1}}{j_h}\right\|_u𝒟_u^hg_{Nk}+\frac{j_{h+1}}{j_h}_v𝒟_u^{h+1}g_{Nk},$$ $`(7.12)`$ we see that (7.11) reduces to the form taken by (7.9) when $`hh+1`$. Moreover, for $`h=1`$ (7.9) becomes $$_v\left(\frac{j_0}{j_1}_ug_{Nk}\right)+_v\mathrm{ln}\|j_1\|\frac{j_0}{j_1}_ug_{Nk}g_{Nk}+\frac{j_0}{j_1}_ug_{Nk1}=0,$$ $`(7.13)`$ which is just (7.8) for $`i=Nk`$, hence certainly true. This completes the proof of (7.9). In the proof that the substitution sequence terminates, we need (7.9) only to derive the fact that, for $`0kN`$, $$g_N=𝒟_u^kg_{Nk}.$$ $`(7.14)`$ The proof of (7.14) is also by induction, this time over $`k`$. Assuming that the equation holds for some $`kN1`$, and writing (7.9) for the particular case $`h=k+1`$, we have $$_v\left(\frac{j_k}{j_{k+1}}_ug_N\right)+_v\mathrm{ln}\|j_{k+1}\|\frac{j_k}{j_{k+1}}_ug_Ng_N+𝒟_u^{k+1}g_{Nk1}=0.$$ $`(7.15)`$ Now, (7.6) allows us to obtain $$g_N=𝒟_u^{k+1}g_{Nk1},$$ $`(7.16)`$ i.e., (7.14) for $`kk+1`$. Since for $`k=0`$ (7.14) is trivially true (and even for $`k=1`$ it can be easily obtained by substituting (7.6) into (7.8) for $`i=N`$), its validity for $`0kN`$ is established. In particular, we shall be interested in the case $`k=N`$, for which $$g_N=𝒟_u^Ng_0.$$ $`(7.17)`$ The equation analogous to (7.9) with $`k=N`$ is found by repeated differentiation of (7.7) with respect to $`u`$; we have, for $`i1`$, $$_v𝒟_u^ig_0+_v\mathrm{ln}\|j_i\|𝒟_u^ig_0𝒟_u^{i1}g_0=0,$$ $`(7.18)`$ whose proof is perfectly analogous to that of (7.9). For $`i=N`$, (7.18) becomes, using (7.17), $$_vg_N+_v\mathrm{ln}\|j_N\|g_N𝒟_u^{N1}g_0=0.$$ $`(7.19)`$ The final step in the proof consists in taking a further derivative with respect to $`u`$ of (7.19), which gives, by (3.11), (7.6), and (7.17), $$\frac{j_{N+1}}{j_N}g_N=0,$$ $`(7.20)`$ i.e., $`j_{N+1}=0`$. We have thus shown that the substitution sequence of a PW$`_N`$ wave equation is upper terminating in $`N`$ steps. To prove that it is also lower terminating in $`N`$ steps, we could manipulate (7.1)–(7.3) to show that $`j_{N1}=\mathrm{}`$, i.e., $`l_{N1}=0`$ by (3.17). However, it is much easier to notice that, starting with the $`u`$-normal form of the wave equation, we can repeat the proof above in a completely symmetric way to get directly $`l_{N1}=0`$. Hence, PW$`{}_{N}{}^{}`$ KN$`_N`$ and, since we know already that the converse is also true, we can conclude that PW$`_N`$ and KN$`_N`$ are equivalent properties. 8.Conclusions and open questions The results presented in this paper clarify the relationships between several properties of wave equations related to the absence of tails in their solutions, some of which were obvious, while some others were never explicitly analysed in the literature, at least to our knowledge. We have seen that the Huygens principle (HP) and the tail-free property (TF) are satisfied by the same equations, that the characteristic propagation property (CPP) is more general only in that it is satisfied, in addition, by special two-dimensional equations, and that the progressing wave propagation property (PW$`_N`$) and the solvability by the Kundt-Newman method (KN<sub>N</sub>) for two-dimensional wave equations are equivalent. The two latter properties were not defined in more than two dimensions, and they are not given in geometrical terms, despite the motivation we gave for introducing PW$`_N`$ in section 3. However, we also showed that in two dimensions PW$`_0`$ is equivalent to CPP, and acquires thus a geometrical meaning. It would be therefore interesting to investigate possible geometrical aspects of PW$`_N`$ equations for higher $`N`$, and whether the definition of PW$`_N`$ can be meaningfully extended to higher dimensions; we will comment here on these issues. The result of section 6, that in $`1+1`$ dimensions CPP is equivalent to PW$`_0`$, has the obvious corollary that no PW$`_N`$ equation with $`N>0`$ can satisfy CPP. This consequence is at first surprising, because from the general form (3.3) of a progressing wave, one is tempted to think that, in two dimensions at least, all PW$`_N`$ wave equations satisfy the CPP. For, by choosing $`R(u)`$ and $`S(v)`$ to be of compact support, their derivatives will also be of compact support, and all of $`\varphi `$ will be made of pieces sandwiched between null coordinate lines, so it will have no tails. The problem with this argument is that, although solutions with no tails are indeed obtained when $`R`$ and $`S`$ are of compact support, the latter functions are not themselves the null data which are involved in the definition of CPP; rather, the data are $$\begin{array}{ccc}& \phi (u)=\underset{i=0}{\overset{N}{}}f_i(u,0)R^{(i)}(u)+\underset{i=0}{\overset{N}{}}g_i(u,0)S^{(i)}(0),\hfill & (8.1)\hfill \\ & \psi (v)=\underset{i=0}{\overset{N}{}}f_i(0,v)R^{(i)}(0)+\underset{i=0}{\overset{N}{}}g_i(0,v)S^{(i)}(v),\hfill & (8.2)\hfill \end{array}$$ and the $`R(u)`$ and $`S(v)`$ that correspond to generic $`\phi `$ and $`\psi `$ of compact support, will in general involve integrals of $`\phi `$ and $`\psi `$ if $`N>0`$, and will thus be of non-compact support. We can clarify this point by considering a typical example of PW$`_N`$ equations, obtained by separating the angular variables in the massless Klein-Gordon equation $`\text{ }\text{ }\text{ }\text{ }\text{ }\mathrm{\Phi }=0`$ in 4-dimensional Minkowski spacetime. Expanding $`\mathrm{\Phi }`$ in spherical harmonics as $`\mathrm{\Phi }(x)=_{lm}\chi _{lm}(t,r)Y_{lm}(\theta ,\phi )`$, and rescaling the radial components to $`\varphi _{lm}:=r\chi _{lm}`$, we have that the $`\varphi _{lm}`$ satisfy the $`l`$-dependent equation $$\left[_t^2_r^2+\frac{l(l+1)}{r^2}\right]\varphi _{lm}\left[_{uv}^2+\frac{l(l+1)}{(vu)^2}\right]\varphi _{lm}=0,$$ $`(8.3)`$ where the null coordinates $`u`$ and $`v`$ are defined as in (2.5) and (2.6), but now with $`x`$ replaced by $`r`$; the two-dimensional equation (8.3) is PW$`_l`$ and can be solved exactly . Since (5.23) is a necessary condition for CPP, and in the case of (8.3) we have $`U=V=0`$, it follows that CPP holds only if $`W=0`$; this is clearly false for $`l1`$. Let us see explicitly the failure of CPP in the simple case $`l=1`$. Equation (8.3) becomes $$\left[_{uv}^2+\frac{2}{(vu)^2}\right]\varphi (u,v)=0,$$ $`(8.4)`$ which is PW$`_1`$; its general solution is $$\varphi (u,v)=\frac{2R(u)}{vu}+R^{}(u)\frac{2S(v)}{vu}+S^{}(v).$$ $`(8.5)`$ We now look for the particular solution generated by null data with support at one point, e.g., $$\begin{array}{ccc}& \phi (u)=\delta (uu_0),\hfill & (8.6)\hfill \\ & \psi (v)=0.\hfill & (8.7)\hfill \end{array}$$ These data correspond to $`R(u)`$ and $`S(v)`$ satisfying the differential equations $$R^{}(u)\frac{2}{u}R(u)+\frac{2}{u}S(0)+S^{}(0)=\delta (uu_0)$$ $`(8.8)`$ and $$S^{}(v)\frac{2}{v}S(v)+\frac{2}{v}R(0)+R^{}(0)=0,$$ $`(8.9)`$ whose general solutions are, respectively, $$R(u)=\frac{u^2}{u_0^2}\theta (uu_0)+au^2+S^{}(0)u+S(0)$$ $`(8.10)`$ and $$S(v)=bv^2+R^{}(0)v+R(0),$$ $`(8.11)`$ with $`a`$ and $`b`$ arbitrary constants. Substituting (8.10) and (8.11) into (8.5), and using the relations $`R(0)=S(0)`$ and $`R^{}(0)=S^{}(0)`$ that follow, e.g., from (8.11), we have $$\varphi (u,v)=\delta (uu_0)+\frac{2}{u_0^2}\frac{uv}{vu}\left[\theta (uu_0)+c\right],$$ $`(8.12)`$ where $`c:=(ab)u_0^2`$ can be prescribed arbitrarily. It is clear from this expression that the only line $`v=`$ const on which $`\varphi `$ is of compact support is $`v=0`$; the choice $`c=0`$ gives a retarded wave, $`c=1`$ an advanced one, and any other value gives a solution which is non-zero almost everywhere on $``$! Turning now to the question of whether it is meaningful to define the class of wave equations whose general solution is a finite sum of progressing waves, in dimension $`m>2`$, the simplest candidate for such an equation would again be the example $`\text{ }\text{ }\text{ }\text{ }\text{ }\mathrm{\Phi }=0`$ in 4-dimensional Minkowski spacetime, which satisfies HP. Using the expansion in spherical harmonics given above, and (8.3), whose general solution is of the type (3.3) with $`N=l`$ and all the functions carrying labels $`l`$ and $`m`$, we can write $$\mathrm{\Phi }(x)=\underset{l=0}{\overset{+\mathrm{}}{}}\underset{i=0}{\overset{l}{}}\underset{m=l}{\overset{l}{}}\left(rf_{ilm}(t,r)Y_{lm}(\theta ,\phi )R_{lm}^{(i)}(tr)+rg_{ilm}(t,r)Y_{lm}(\theta ,\phi )S_{lm}^{(i)}(t+r)\right),$$ $`(8.13)`$ which contains derivatives up to arbitrarily high order of an infinite number of free functions $`R_{lm}`$ and $`S_{lm}`$, and is therefore not of the form (3.3). This shows that, at least in four dimensions, the definition of PW$`_N`$ as it stands is empty for the case of spherical wavefronts, and makes it plausible that it is ill-posed in general. It is still useful to define a “reduced” PW$`_N`$ property for an $`m`$-dimensional equation, in the sense that all two-dimensional equations, obtained by separating out appropriately chosen coordinates, may have progressing wave general solutions, as in the case of (8.3). A more general example, in four dimensions, would be obtained with equations of the type $$\left[f(y,z)K_{t,x}+g(t,x)H_{y,z}\right]\varphi (t,x,y,z)=0,$$ $`(8.14)`$ where $`f`$ and $`g`$ are nonvanishing functions and $`K_{t,x}`$, $`H_{y,z}`$ are suitable differential operators acting on the variables used as subscripts. With the separation of variables $`\varphi (t,x,y,z)=\psi (t,x)\phi (y,z)`$, we see that (8.14) satisfies the reduced progressing wave propagation property if, for each value of $`\alpha `$ admitted by the equation $$H_{y,z}\phi (y,z)+\alpha f(y,z)\phi (y,z)=0,$$ $`(8.15)`$ the reduced wave equation $$K_{t,x}\psi (t,x)\alpha g(t,x)\psi (t,x)=0$$ $`(8.16)`$ is PW$`_N`$ for some $`N(\alpha )`$. The high degree of dependence of the reduction procedure—hence of the properties of the reduced equation—from specific, noncovariant structure is evident. It is thus very difficult to make general statements about this property and its possible relationships with other ones we have discussed in this paper, independently of the explicit form of the wave equation and the choice of coordinates singled out for separation (i.e., of wave front). It would be interesting to know whether a meaningful, truly $`m`$-dimensional generalisation of the PW$`_N`$ property can be given, for another reason as well. The expressions for the advanced or retarded Green function of $`\text{ }\text{ }\text{ }\text{ }\text{ }\mathrm{\Phi }=0`$ in $`m`$-dimensional Minkowski space, which for even $`m4`$ are of the form $$G_m(t^{},𝐱^{};t,𝐱)=\underset{i=0}{\overset{N}{}}c_{mi}(\|𝐱^{}𝐱\|)\delta ^{(i)}(t^{}t\pm \|𝐱^{}𝐱\|),$$ $`(8.17)`$ with $`N=(m4)/2`$, remind one of definition (3.2), and suggest a possible connection between HP and progressing waves, this time of order $`N>0`$, in dimension $`m>2`$. So far we have not been able to find such a relationship, however. For now, PW$`_N`$ equations are left with no simple interpretation in terms of tails, and no true generalisation to higher dimensions, but their usefulness derives from the fact that they are exactly solvable by the method of Kundt and Newman , as we saw in section 7, and from their relationship to Toda lattices . Taking into account all of these remarks, there is little doubt that, as far as studies of tails are concerned, the definition that should be preferred is, for its generality, the one based on HP or, equivalently, on TF. With a notion of wave tails that is now completely clear and unambiguous, one can conduct further research in several directions. Here is a list of some problems that one may address: (i) Which are the conditions on $`g^{ab}`$, $`H^a`$, and $`K`$ that make (1.3) tail-free? Because of the geometrical interpretation of $`g_{ab}`$ as a metric in $``$, it is natural to restate this question by asking in which spacetimes (i.e., for which class of $`g^{ab}`$) is the wave equation tail-free. This is sometimes called the Hadamard problem, and has received some attention over the years . (ii) What happens for vector and tensor fields? These cases include the propagation of electromagnetic and gravitational waves and are more realistic—though somewhat more complicated—than the scalar field one. (iii) Which are the physical consequences and effects of wave tails? The situation that has been investigated most in detail is that of radiation from compact objects . The corresponding problem in a cosmological background has apparently been neglected, probably because of the belief that any observable effect should be extremely small. However, in this case curvature never drops off and scattering continues forever, so it might convey a relevant amount of radiation in the interior of the light cone . This brings immediately out a new problem, namely how much radiation goes into the tail, which leads us to the next question. (iv) How to quantify tails? A very natural and intuitive way would be to calculate a reflection coefficient; this is possible when backscattering is localised and there are regions in which the field is free (even asymptotically, or after a suitable coordinates transformation has been performed ). In a cosmological context, however, the property that makes the phenomenon potentially interesting—i.e., the fact that backscattering takes place always and everywhere—at the same time prevents us from defining purely ingoing and outgoing solutions of the wave equation, and hence from computing reflection and transmission coefficients . It is necessary to find alternative ways to quantify tails, perhaps based on the ratio between their energy content and the total energy of the field. Acknowledgements This work was supported by the Commission of the European Communities under the DG-XII contract no. CI1\*-0540-M(TT) and the DG-III contract no. 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warning/0002/astro-ph0002175.html	ar5iv	text	# Finite-Correlation-Time Effects in the Kinematic Dynamo Problem ## I Introduction The study of the statistics of magnetic fluctuations excited by a random Gaussian white-noise-like advecting velocity field, was pioneered by Kazantsev , and, in more recent times, has generated a considerable amount of research (see, e. g., Ref. and references therein, as well as Refs. ). While much attention has concentrated on resistive dynamo problems, most often for very large magnetic Prandtl numbers, it is well known that the fundamental Zeldovich’s, or “stretch–twist–fold,” mechanism of the magnetic-energy amplification (the so-called “fast dynamo”) is active regardless of the presence of the resistive (diffusive) regularization . If the initial seed magnetic field is concentrated on the scales of the same order as the characteristic scales of the advecting velocity, the stretching and folding of the magnetic-field lines by the random flow leads to an exponential growth of the magnetic fluctuations at scales that decrease exponentially fast, until the diffusive scales are reached . This scenario is common in astrophysical applications such as the turbulence in the interstellar medium or in the protogalaxy where the Prandtl number ranges from $`10^{14}`$ to $`10^{22}`$, giving rise to 7 to 11 decades of small (subviscous) scales available to the magnetic fluctuations . In fact, the initial diffusion-free regime may well be the only one practically important in such applications as far as the kinematic approximation is concerned, since the nonlinear saturation effects are likely to set in before the diffusion scales are reached . On the fundamental physical level, the diffusion-free regime, in which the magnetic-field lines are fully frozen into the flow, exhibits most clearly the underlying symmetry properties of the passive advection . With a few notable exceptions (such as Refs. ), the dominant approach in the existing literature on the turbulent kinematic dynamo problem has been to study the statistics of passive magnetic fields advected by a flow $`𝐮(t,𝐱)`$ whose two-time correlation function is approximated by a $`\delta `$ function, $`𝐮(t)𝐮(t^{})\delta (tt^{})`$. This white-noise property of the velocity greatly simplifies matters: the evolution equations for such statistical quantities as the correlation functions and probability density function of the magnetic field can be derived in closed form and yield themselves to exact solution. In this paper, we relax the white-noise assumption and explore the effects that arise when a finite-time correlated velocity field is introduced. This immediately raises the level of difficulty associated with solving the statistical problem. Within the theoretical framework adopted here, the difficulty can be described in the following terms. In the zero-correlation-time approximation, one essentially has to deal with only one closed differential equation that fully determines the desired statistics. Allowing for a finite velocity correlation time leads to an infinite number of interlinked integro-differential equations involving time-history integrals. These equations form an infinite open hierarchy that formally constitutes the exact description of the problem (these matters are explained in more detail in Sec. II A). Solving this hierarchy in its entirety without additional assumptions appears to be an impossible task. The most obvious way to make progress is clearly to try a perturbative approach, i.e., to consider the kinematic dynamo problem with an advecting field whose correlation time is short, but finite. If the correlation time $`\tau _\mathrm{c}`$ is assumed to be small, one can expect to be able to construct an expansion in the powers of $`\tau _\mathrm{c}`$ (in what follows, we will frequently refer to it as the $`\tau `$ expansion) and calculate corrections to the growth rates of the moments of the magnetic field. This is the program that we undertake here. We consider the one-point statistics of the passive magnetic field in the diffusion-free regime. In this context, the infinite hierarchy we have mentioned above interrelates the one-point probability density function (PDF) of the magnetic field and an infinite set of response functionals. These are averaged multiple functional derivatives of the magnetic field with respect to the velocity field and its gradients. We develop a functional expansion method that allows us to calculate successive terms in the $`\tau `$ expansion and derive in a closed form a Fokker–Planck equation for the one-point PDF of the magnetic field. We limit ourselves to advancing the expansion one order beyond the zero-correlation-time approximation. The result is a set of corrections to the growth rates of all moments of the magnetic field. These corrections are negative, so the growth rates are reduced. The expansion is carried out assuming that the velocity correlation time is small and keeping the time integral of the velocity correlation function fixed. The latter constraint ensures that the dynamo growth rate remains finite when the correlation time vanishes. An alternative way, which is sometimes deemed preferable on physical grounds (see, e.g, Ref. ), is to fix the total energy of the velocity field. Since a $`\delta `$-correlated velocity field must necessarily possess infinite energy, fixing the energy at a finite value leads to vanishing of the growth rates when $`\tau _\mathrm{c}=0`$. The relative ordering of the terms in the expansion is, however, the same, regardless of what is kept fixed, so the technical side of the expansion method is unaffected. Our expansion technique will be given detailed treatment in the body of this paper. Here, let us rather discuss the finite-correlation-time effects that can be distilled on the basis of our approach. As it turns out, a number of new interesting phenomena manifest themselves already at the level of the short-but-finite-correlation-time approximation. In the case of the $`\delta `$-correlated advecting flow, the one-point statistics of the passive magnetic field are universal in the sense that they only depend on one small-scale property of the velocity: the time integral of the one-point correlation tensor of its gradients, $`dt𝐮(t)𝐮(0)`$. The essential novelty in the case of finite correlation time is that this small-scale universality is lost on two accounts. First, the $`\tau `$ expansion exhibits a sensitive dependence on the specific shape of the time-correlation profile of the velocity field (in recent literature, this was first explicitly pointed out by Boldyrev ; see also Refs. ). Namely, multiple time integrals of products of velocity correlation functions enter the expressions for the expansion coefficients. Choosing different correlation profiles leads to order-one changes in the values of these coefficients. The root of this nonuniversality lies in the topology of the vertex-correction diagrams that contribute to the orders higher than the zeroth in the $`\tau `$ expansion (see Sec. II E). Second, the first-order terms of the $`\tau `$ expansion feature a part that arises from the fourth-order derivatives of the velocity correlation function, i.e., from the second derivatives of the velocity field. In the one-point statistical approach, this is the first manifestation of the more general tendency that introducing finite correlation times brings into play the large-scale structure of the velocity field. A related effect is the loss of Galilean invariance due to the fact that the expansion terms also depend on the actual energy of the velocity field, i.e., on the rms value of the sweeping velocity. Indeed, now that the trajectories of the fluid elements have a “memory” of themselves, which extends approximately one $`\tau _\mathrm{c}`$ back in time, we should naturally expect that there will appear an effective “correlation length” of the velocity (in what regards the one-point statistics of the fields it advects) approximately equal to $`u\tau _\mathrm{c}`$. Therefore, the one-point statistics of the passive fields now depend not only on the instantaneous velocity difference between two fluid particles that meet at a given time (i.e., the velocity gradient at a point), but also on the velocity that swept them into place and on the variation of the velocity gradient over the correlation length. This appearance of first-order corrections due to the second derivatives of the flow is a new effect, which indicates, in particular, that the customary approximation used in the Batchelor regime, where the advecting velocity is assumed to be locally linear , is only justified for the $`\delta `$-correlated-in-time advecting fields. Such are the main qualitative consequences of introducing a finite-time-correlated velocity field into the kinematic dynamo problem (or, in general, any passive-advection model). A few words are in order as to the quantitative impact of a finite correlation time on the dynamo action. As we have already mentioned, the effect of the first-order corrections is to reduce the growth rates of all moments of the magnetic field. Besides the nonuniversal dependence on the spatial and temporal structure of the velocity correlation function, the reduction depends in a universally calculable way on the usual set of parameters: the order of the moment, the dimension of space, and the degree of compressibility of the flow. The overall magnitude of this reductive effect is measured by the expansion parameter, which is of the order of $`\tau _\mathrm{c}\gamma `$, where $`\gamma `$ is the growth rate of the magnetic energy. It is not hard to demonstrate (see Sec. III D) that $`\tau _\mathrm{c}\gamma d(\tau _\mathrm{c}/\tau _{\mathrm{eddy}})^2`$, where $`\tau _{\mathrm{eddy}}`$ is the “eddy-turnover” time of the advecting turbulent velocity field and $`d`$ is the dimension of space. In a standard Kolmogorov-type turbulence setting, one would, of course, expect any such approximation to be valid at best marginally, since $`\tau _\mathrm{c}\tau _{\mathrm{eddy}}`$. Astrophysical plasmas offer more variety in this respect, as their driving forces (typically supernova explosions) can, in fact, decorrelate faster than the turbulent eddies turn over . In any event, the small-$`\tau _\mathrm{c}`$ expansion does not offer much more than qualitative, or, at best, semiquantitative, information about the way the dynamo action is modified by the finiteness of the correlation time. It is, of course, clear that introducing a finite correlation time cannot altogether suppress the fast-dynamo mechanism . On the other hand, our conclusion that some reduction of the growth rate should be expected, is corroborated by numerical evidence that suggests a reduction of about 40% to 50%. In fact, in Sec. IV, we offer a semiquantitative evaluation of the finite-$`\tau _\mathrm{c}`$ correction to the growth rate which yields a reduction of approximately 40% in the three-dimensional case and for $`\tau _\mathrm{c}\tau _{\mathrm{eddy}}`$. Of course, this is at best just an indication of the well-behaved character of our expansion, rather than a truly solid quantitative confirmation of it. The literature on the $`\tau `$ expansion and finite-correlation-time effects is not extensive. Kliatskin and Tatarskii were the first to propose the hierarchy of equations for the response functionals as a starting point for a method of successive approximations as applied to the description of waves propagating in a medium with random inhomogeneities. Vainshtein applied this method to the mean-field kinematic dynamo theory. The Kliatskin–Tatarskii method and its relation to our functional expansion method are discussed at the end of Sec. II C. Van Kampen and Terwiel developed the so-called cumulant expansion method; van Kampen’s review article also contains a good critical survey of other $`\tau `$-expansion schemes predating his work. His method was later applied in the kinematic-dynamo context by Knobloch and Chandran . Their treatment was Lagrangian and did not include any effects due to the explicit spatial dependence in the induction equation. Consequently, the nonuniversality of the $`\tau `$ expansion with respect to the spatial structure of the velocity correlator was not captured. The van Kampen method is discussed in detail in Sec. II D. Parallel to our development of the functional expansion method, Boldyrev proposed a $`\tau `$-expansion method that was based on the exact solution of the induction equation in the Lagrangian frame and offered a way to calculate the second moment of the magnetic field that elicited the nonuniversal character of the $`\tau `$ expansion with respect to both temporal and spatial properties of the velocity correlation tensor. Molchanov, Ruzmaikin, and Sokoloff considered the statistics of the kinematic dynamo in a renovating flow using the formalism of infinite products of random matrices. (See also Ref. for the treatment of the kinematic mean-field dynamo in a renovating flow.) A version of their approach was later advanced by Gruzinov, Cowley, and Sudan . Considerable progress was achieved in a nonperturbative way by Chertkov et al. , who studied the passive-scalar problem in two dimensions for arbitrary velocity correlation times. However, their method only works in the two-dimensional case. Thus, while we now seem to have a fairly good understanding of the structure of the $`\tau `$ expansion and such qualitative features as the loss of the small-scale universality, an adequate nonperturbative theory of the kinematic dynamo and passive advection in finite-time-correlated turbulent velocity fields remains an open problem. This paper is organized in the following way. In Sec. II, our functional expansion method is systematically developed on the example of the simplest available passive-advection problem: that of the Lagrangian passive vector in an incompressible flow. In this model, no explicit spatial dependence is present. In Sec. II A, Sec. II B, and Sec. II C, we present a functional formalism that allows one to systematically construct successive terms in the $`\tau `$-expanded Fokker–Planck equation. The dependence of the expansion coefficients on the specific functional form of the velocity time correlation profile emerges. The expansion is carried out up to the first order in $`\tau _\mathrm{c}`$. In Sec. II D, our method is compared with the van Kampen cumulant expansion method . We ascertain that results obtained via the van Kampen method are consistent with ours. Finally, in Sec. II E, we discuss the underlying structure of the $`\tau `$ expansion in diagrammatic terms. In Sec. III, the general arbitrarily compressible space-dependent dynamo problem is solved with the aid of the functional expansion. At this level, the nonuniversality with respect to the spatial structure of the velocity correlations, as well as the loss of Galilean invariance, become evident. In Sec. III A, we explain the emergence of an infinite hierarchy of equations for the characteristic function and various averaged response functionals of the magnetic field in the passive dynamo problem with finite-time-correlated advecting flow. The hierarchy is advanced up to the emergence of the second-order response functions. In Sec. III B, we construct the $`\tau `$ expansion up to first order in the correlation time, which leads to a closed equation for the characteristic function of the magnetic field. In Sec. III C, we derive the Fokker–Planck equation for the one-point PDF of the magnetic-field strength valid to first order in the correlation time. The distribution is lognormal. In Sec. III D, we calculate the rates of growth of all moments of the magnetic field with (negative) first-order corrections. Finally, in Sec. IV, we give a semiquantitative argument that relates the expansion parameter to the ratio of the correlation and eddy-turnover times of the velocity field. We also evaluate the finite-$`\tau _\mathrm{c}`$ reduction of the magnetic-energy growth rate in a model incompressible turbulence consisting of eddies of a fixed size. In Appendix A, we provide the basic relations that allow one to express the results we have obtained in the configuration space in terms of the spectral characteristics of the velocity field. Some of the more cumbersome technical details of the $`\tau `$ expansion are exiled to Appendix B. ## II The Gaussian Functional Expansion Formalism In this Section, we explain the Gaussian functional method for constructing the short-correlation-time expansion for passive advection problems. Working out such expansions for specific problems often involves a fair amount of algebra, which tends to obscure the otherwise transparent ideas behind them. In an attempt at the maximum possible clarity of exposition, we first consider a model that, while preserving most of the essential features of the passive-advection problems, offers much greater technical simplicity. Namely, let us consider the following stochastic equation in $`d`$ dimensions: $`_tB^i=\sigma _k^iB^k.`$ (1) All the fields involved explicitly depend on time only. The specific initial distribution of $`B^i`$ is not important for the derivation or the validity of the results below. Spatial isotropy is always assumed. The matrix field $`\sigma _k^i(t)`$ is Gaussian with zero mean and a given two-point correlation tensor: $`<\sigma _k^i(t)\sigma _l^j(t^{})>=T_{kl}^{ij}\kappa (tt^{}),`$ (2) $`T_{kl}^{ij}=\delta ^{ij}\delta _{kl}+a\left(\delta _k^i\delta _l^j+\delta _l^i\delta _k^j\right),a=1/(d+1).`$ (3) These equations can be interpreted to describe the evolution of a passive magnetic field in a Lagrangian frame, where the Lagrangian advecting velocity field is Gaussian and incompressible, and the tensor $`\sigma _k^i`$ is its gradient matrix. In the more general context of the theory of passive advection, equations (1) and (2) model the stochastic dynamics of a vector connecting two Lagrangian tracer particles in an ideal fluid. We assume that the temporal correlation function $`\kappa (tt^{})`$ of $`\sigma _k^i(t)`$ has a certain characteristic width $`\tau _\mathrm{c}`$, i.e., the field $`\sigma _k^i(t)`$ possesses a correlation time $`\tau _\mathrm{c}`$. Our task in this section is to construct an expansion of the statistics of $`B^i(t)`$ in powers of $`\tau _\mathrm{c}`$, which is assumed to be small. The limit $`\tau _\mathrm{c}0`$ ought to be taken in such a way that the time integral of the correlation function is kept constant: $`{\displaystyle _0^{\mathrm{}}}d\tau \kappa (\tau )={\displaystyle \frac{\overline{\kappa }}{2}}=\mathrm{const}\mathrm{and}\tau _\mathrm{c}\overline{\kappa }1.`$ (4) The white-noise limit of zero correlation time is realized by setting $`\kappa (\tau )=\overline{\kappa }\delta (\tau )`$. The first step in our averaging scheme is to define the characteristic function of the field $`B^i(t)`$, $`Z(t;\mu )=<\stackrel{~}{Z}(t;\mu )>=<\mathrm{exp}\left[i\mu _iB^i(t)\right]>.`$ (5) Here and in what follows, the overtildes designate unaveraged random functions. Upon differentiating $`\stackrel{~}{Z}(t;\mu )`$ with respect to time and making use of Eq. (1), we obtain a new stochastic equation: $`_t\stackrel{~}{Z}=\sigma _k^i\mu _i{\displaystyle \frac{}{\mu _k}}\stackrel{~}{Z}=\widehat{\mathrm{\Lambda }}_i^k\sigma _k^i\stackrel{~}{Z},`$ (6) where the auxiliary operator $`\widehat{\mathrm{\Lambda }}_i^k`$ has been introduced for the sake of notational compactness. Our objective now is to learn how to obtain a closed equation for the averaged characteristic function $`Z(t;\mu )`$, i.e., how to average Eq. (6) when $`\sigma _k^i`$ has a nonzero correlation time. The inverse Fourier transform (with respect to $`\mu _i`$) of the resulting equation will be the Fokker–Planck equation for the PDF $`P(t;𝐁)`$ of the passive field $`B^i`$. ### A The Hierarchy of Response Functions We start the construction of the functional expansion by developing an exact formalism that describes the one-point statistics of the field $`B^i(t)`$. Let us average both sides of Eq. (6) and “split” the mixed average that arises on the right-hand side with the aid of the well-known Furutsu–Novikov (or “Gaussian-integration”) formula : $`_tZ(t)=\widehat{\mathrm{\Lambda }}_i^k<\sigma _k^i(t)\stackrel{~}{Z}(t)>=\widehat{\mathrm{\Lambda }}_i^kT_{k\alpha }^{i\beta }{\displaystyle _0^t}dt^{}\kappa (tt^{})G_\beta ^\alpha (t\|t^{}),`$ (7) where we have suppressed the $`\mu `$’s in the arguments, used the formula (2) for the second-order correlation tensor of $`\sigma _k^i(t)`$, and introduced the averaged first-order response function: $`G_\beta ^\alpha (t\|t^{})=<\stackrel{~}{G}_\beta ^\alpha (t\|t^{})>=<{\displaystyle \frac{\delta \stackrel{~}{Z}(t)}{\delta \sigma _\alpha ^\beta (t^{})}}>.`$ (8) This function is subject to the causality constraint: $`G_\beta ^\alpha (t\|t^{})=0`$ if $`t^{}>t`$ \[hence the upper integration limit in Eq. (7)\]. Integrating Eq. (6) from $`0`$ to $`t`$, taking the functional derivative $`\delta /\delta \sigma _\alpha ^\beta (t^{})`$ of both sides, averaging, setting $`t^{}=t`$, and taking causality into account, we get $`G_\beta ^\alpha (t\|t)=\widehat{\mathrm{\Lambda }}_\beta ^\alpha Z(t).`$ (9) We have thus obtained the equal-time form of $`G_\beta ^\alpha (t\|t^{})`$. In order to determine the response function at $`t^{}<t`$, we simply take the functional derivative $`\delta /\delta \sigma _\alpha ^\beta (t^{})`$ of both sides of Eq. (6) and find that each element of the unaveraged tensor $`\stackrel{~}{G}_\beta ^\alpha (t\|t^{})`$ satisfies an equation identical in form to Eq. (6). Upon averaging this, we obtain an evolution equation for $`G_\beta ^\alpha (t\|t^{})`$ subject to the initial condition (9) at $`t=t^{}`$: $`_tG_\beta ^\alpha (t\|t^{})=\widehat{\mathrm{\Lambda }}_j^l<\sigma _l^j(t)\stackrel{~}{G}_\beta ^\alpha (t\|t^{})>.`$ (10) This equation can now be handled in the same fashion as Eq. (6), the average on the right-hand side split via the Furutsu–Novikov formula in terms of the correlation tensor of $`\sigma _l^j(t)`$ and the appropriately defined second-order averaged response function $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t\|t_1,t_2)`$. At equal times, the latter can be expressed in terms of $`G_\beta ^\alpha (t\|t^{})`$ just as $`G_\beta ^\alpha (t\|t)`$ was expressed in terms of $`Z(t)`$. At different times, we obtain the evolution equation for $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t\|t_1,t_2)`$ by taking the functional derivative of the equation for $`\stackrel{~}{G}_\beta ^\alpha (t\|t^{})`$ and averaging. An infinite linked hierarchy can be constructed by further iterating this procedure and introducing response functions of ascending orders. Let us give the general form of this hierarchy. Define the $`n`$th-order averaged response function: $`G_{\beta _1\mathrm{}\beta _n}^{\alpha _1\mathrm{}\alpha _n}(t\|t_1,\mathrm{},t_n)=<{\displaystyle \frac{\delta \stackrel{~}{Z}(t)}{\delta \sigma _{\alpha _1}^{\beta _1}(t_1)\mathrm{}\delta \sigma _{\alpha _n}^{\beta _n}(t_n)}}>.`$ (11) This function has two essential properties: (i) it is causal: $`G_{\beta _1\mathrm{}\beta _n}^{\alpha _1\mathrm{}\alpha _n}(t\|t_1\mathrm{}t_n)=0`$ if any $`t_i>t`$; (ii) it remains invariant under all simultaneous permutations of the times $`t_1,\mathrm{},t_n`$ and indices $`\alpha _1,\mathrm{},\alpha _n`$, $`\beta _1,\mathrm{},\beta _n`$, which correspond to changes of the order of functional differentiation in the definition (11). The $`n`$th-order response function satisfies the following recursive relations: if $`t_1,\mathrm{},t_n<t`$, $`_tG_{\beta _1\mathrm{}\beta _n}^{\alpha _1\mathrm{}\alpha _n}(t\|t_1\mathrm{}t_n)=\widehat{\mathrm{\Lambda }}_i^kT_{k\alpha _{n+1}}^{i\beta _{n+1}}{\displaystyle _0^t}dt_{n+1}\kappa (tt_{n+1})G_{\beta _1\mathrm{}\beta _{n+1}}^{\alpha _1\mathrm{}\alpha _{n+1}}(t\|t_1\mathrm{}t_{n+1});`$ (12) if, say, $`t_n=t`$ and $`t_1,\mathrm{},t_{n1}t_n`$, $`G_{\beta _1\mathrm{}\beta _{n1}\beta _n}^{\alpha _1\mathrm{}\alpha _{n1}\alpha _n}(t\|t_1\mathrm{}t_{n1},t)=\widehat{\mathrm{\Lambda }}_{\beta _n}^{\alpha _n}G_{\beta _1\mathrm{}\beta _{n1}}^{\alpha _1\mathrm{}\alpha _{n1}}(t\|t_1\mathrm{}t_{n1}).`$ (13) The hierarchy is “forward” at different times and “backward” at equal times. The characteristic function $`Z(t)`$ is formally treated as the zeroth-order response function. ### B The White-Noise Approximation The white-noise approximation is obtained by setting $`\kappa (tt^{})=\overline{\kappa }\delta (tt^{})`$. We are then left with just Eq. (7), where the time history integral reduces to $`\frac{1}{2}\overline{\kappa }G_\beta ^\alpha (t\|t)`$, which is substituted from Eq. (9). This produces a closed evolution equation for $`Z(t)`$. The Fourier transform of it is the Fokker–Planck equation for the PDF of $`B^i`$ at time $`t`$ in the $`\delta `$-correlated regime: $`_tP(t)={\displaystyle \frac{\overline{\kappa }}{2}}T_{k\alpha }^{i\beta }\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_\beta ^\alpha P(t)={\displaystyle \frac{\overline{\kappa }}{2}}\widehat{L}P(t).`$ (14) In order to be not overly burdened by notation, we typically use the same symbol for denoting an operator in the Fourier space of the $`\mu `$’s and its analog in the configuration space of the $`B`$’s. This should lead to no confusion, as the context will always be clear. Thus, $`\widehat{\mathrm{\Lambda }}_i^k=\mu _i{\displaystyle \frac{}{\mu _k}}={\displaystyle \frac{}{B^i}}B^k=\left(\delta _i^k+B^k{\displaystyle \frac{}{B^i}}\right).`$ (15) Due to isotropy, the probability density function $`P`$ depends on the absolute value $`B=\|𝐁\|`$ only. The operator $`\widehat{L}`$ in Eq. (14) can therefore be written in the following isotropic form: $`\widehat{L}=T_{k\alpha }^{i\beta }\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_\beta ^\alpha ={\displaystyle \frac{d1}{d+1}}\left[B^2{\displaystyle \frac{^2}{B^2}}+(d+1)B{\displaystyle \frac{}{B}}\right],`$ (16) which turns Eq. (14) into the familiar Fokker–Planck equation for the one-point PDF of the magnetic field in the kinematic $`\delta `$-correlated dynamo problem taken for the incompressible flow . The resulting distribution is lognormal and the moments of $`B`$ satisfy $`_tB^n={\displaystyle \frac{1}{2}}{\displaystyle \frac{d1}{d+1}}n\left(n+d\right)\overline{\kappa }B^n.`$ (17) This is the expected outcome, because, as was shown in Ref. , the Lagrangian and Eulerian statistics are the same for the $`\delta `$-correlated incompressible flow. Thus, the solution in the $`\delta `$-correlated limit is quite elementary. Things become much more complicated once the white-noise assumption is relaxed and a nonzero, however small, velocity correlation time is introduced. ### C The Recursive Expansion In order to construct an expansion in small correlation time, it is convenient to combine the equations (12) and (13) into one recursive integral relation that expresses the $`n`$th-order response function in terms of its immediate precursor and its immediate successor: $`G_{\beta _1\mathrm{}\beta _n}^{\alpha _1\mathrm{}\alpha _n}(t\|t_1\mathrm{}t_n)=\widehat{\mathrm{\Lambda }}_{\beta _n}^{\alpha _n}G_{\beta _1\mathrm{}\beta _{n1}}^{\alpha _1\mathrm{}\alpha _{n1}}(t_n\|t_1\mathrm{}t_{n1})`$ (18) $`+\widehat{\mathrm{\Lambda }}_i^kT_{k\alpha _{n+1}}^{i\beta _{n+1}}{\displaystyle _{t_n}^t}dt^{}{\displaystyle _0^t^{}}dt_{n+1}\kappa (t^{}t_{n+1})G_{\beta _1\mathrm{}\beta _{n+1}}^{\alpha _1\mathrm{}\alpha _{n+1}}(t^{}\|t_1\mathrm{}t_{n+1}).`$ (19) The above relation is exact and valid for $`t_1,\mathrm{},t_{n1}t_nt`$. Due to the permutation symmetry of the response functions, this does not limit the generality. The desired expansion is constructed by repeated application of the formula (19). Let us substitute the formula (19) with $`n=1`$ for the first-order response function into the right-hand side of Eq. (7): $`_tZ(t)=\widehat{L}{\displaystyle _0^t}dt_1\kappa (tt_1)Z(t_1)`$ (20) $`+T_{k\alpha _1}^{i\beta _1}T_{n\alpha _2}^{m\beta _2}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_m^n{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^t^{}}dt_2\kappa (tt_1)\kappa (t^{}t_2)G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{}\|t_1,t_2),`$ (21) where the operator $`\widehat{L}`$ is defined in (16). We now use the formula (19) to express the second-order response function on the right-hand side of the above equation: for $`t_2>t_1`$, we have $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{}\|t_1,t_2)=\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}G_{\beta _1}^{\alpha _1}(t_2\|t_1)`$ (22) $`+\widehat{\mathrm{\Lambda }}_p^qT_{q\alpha _3}^{p\beta _3}{\displaystyle _{t_2}^t^{}}dt^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}dt_3\kappa (t^{\prime \prime }t_3)G_{\beta _1\beta _2\beta _3}^{\alpha _1\alpha _2\alpha _3}(t^{\prime \prime }\|t_1,t_2,t_3),`$ (23) while for $`t_2<t_1`$ we flip the variables, $`t_1t_2`$, to make sure that the first-order response function on the right-hand side do not vanish: $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{}\|t_1,t_2)=\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}G_{\beta _2}^{\alpha _2}(t_1\|t_2)`$ (24) $`+\widehat{\mathrm{\Lambda }}_p^qT_{q\alpha _3}^{p\beta _3}{\displaystyle _{t_1}^t^{}}dt^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}dt_3\kappa (t^{\prime \prime }t_3)G_{\beta _2\beta _1\beta _3}^{\alpha _2\alpha _1\alpha _3}(t^{\prime \prime }\|t_2,t_1,t_3).`$ (25) The recursion relation (19) is now applied to the first-order response functions in the formulas (23) and (25): $`G_{\beta _1}^{\alpha _1}(t_2\|t_1)=\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1)+\widehat{\mathrm{\Lambda }}_p^qT_{q\alpha _3}^{p\beta _3}{\displaystyle _{t_1}^{t_2}}dt^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}dt_3\kappa (t^{\prime \prime }t_3)G_{\beta _1\beta _3}^{\alpha _1\alpha _3}(t^{\prime \prime }\|t_1,t_3),`$ (26) $`G_{\beta _2}^{\alpha _2}(t_1\|t_2)=\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2)+\widehat{\mathrm{\Lambda }}_p^qT_{q\alpha _3}^{p\beta _3}{\displaystyle _{t_2}^{t_1}}dt^{\prime \prime }{\displaystyle _0^{t^{\prime \prime }}}dt_3\kappa (t^{\prime \prime }t_3)G_{\beta _2\beta _3}^{\alpha _2\alpha _3}(t^{\prime \prime }\|t_2,t_3).`$ (27) All this must be substituted into Eq. (21): $`_tZ(t)`$ $`=`$ $`\widehat{L}{\displaystyle _0^t}dt_1\kappa (tt_1)Z(t_1)`$ (28) $`+`$ $`T_{k\alpha _1}^{i\beta _1}T_{n\alpha _2}^{m\beta _2}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_m^n\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^{t_1}}dt_2\kappa (tt_1)\kappa (t^{}t_2)Z(t_2)`$ (29) $`+`$ $`T_{k\alpha _1}^{i\beta _1}T_{n\alpha _2}^{m\beta _2}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_m^n\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _{t_1}^t^{}}dt_2\kappa (tt_1)\kappa (t^{}t_2)Z(t_1)`$ (30) $`+`$ $`R(t),`$ (31) where the remainder $`R(t)`$ contains the assembled terms that involve quintuple time integrals. So far, all the manipulations we have carried out have been exact. It is now not hard to perceive the emerging contours of the small-$`\tau _\mathrm{c}`$ expansion. Since the time-correlation function $`\kappa (tt^{})`$ is a profile of width $`\tau _\mathrm{c}`$, the area under which is constant and equal to $`\overline{\kappa }`$, the triple time integrals in Eq. (31) are of the order of $`\tau _\mathrm{c}\overline{\kappa }^2`$, while the quintuple time integrals absorbed into $`R(t)`$ are of the order of $`\tau _\mathrm{c}^2\overline{\kappa }^3`$. Further application of the recursion formula (19) to the second- and third-order response functions in the equations (23)–(27) leads to the appearance of more multiple time integrals of the time-correlation function $`\kappa (tt^{})`$. These integrals are of orders $`\tau _\mathrm{c}^2\overline{\kappa }^3`$, $`\tau _\mathrm{c}^3\overline{\kappa }^4`$, etc. We would only like to keep terms up to first order in the correlation time. The remainder term $`R(t)`$ in Eq. (31) can therefore be dropped. Remark on the physics of the $`\tau `$ expansion. The above argument is based on the stipulation made at the beginning of this Section that the $`\tau `$ expansion must be carried out keeping the integral of the velocity time correlation function constant \[formula (4)\]. This requirement is natural because it leads to finite dynamo growth rates in the limit of zero correlation time \[see Sec. II B, Eq. (17)\]. However, it is also acceptable to institute an alternative, arguably more physical, requirement that the total energy of the velocity field (i.e., the rms velocity) remain constant (as, e.g, in the numerics of Ref. ). Quantitatively, this means that $`\kappa (0)`$, rather than $`_0^{\mathrm{}}d\tau \kappa (\tau )`$, is kept fixed. Under this constraint, the terms that we have previously estimated to be of orders $`\overline{\kappa }`$, $`\tau _\mathrm{c}\overline{\kappa }^2`$, $`\tau _\mathrm{c}^2\overline{\kappa }^3`$, etc., and hence, $`\overline{\kappa }`$ being constant, to represent the zeroth, first, second, etc. orders of the $`\tau `$ expansion, should now be reevaluated as follows. Since $`\overline{\kappa }\tau _\mathrm{c}\kappa (0)`$, these terms are of orders $`\tau _\mathrm{c}\kappa (0)`$, $`\tau _\mathrm{c}^3\kappa (0)^2`$, $`\tau _\mathrm{c}^5\kappa (0)^3`$, etc., and therefore constitute the first, third, fifth, etc. orders of the expansion. The shortcoming of this approach is that the dynamo growth rates vanish when $`\tau _\mathrm{c}=0`$, so formally there is no nontrivial zero-correlation-time limit. With $`\overline{\kappa }=\mathrm{const}`$, this problem was avoided because the energy was formally infinite when $`\tau _\mathrm{c}=0`$ (a $`\delta `$-correlated velocity field cannot have a finite energy). In any event, we see that, since the difference between keeping $`\overline{\kappa }`$ and $`\kappa (0)`$ constant does not affect the relative magnitudes of the terms in the expansion, our expansion scheme remains valid in both cases. Let us therefore proceed with our construction. The dependence of the right-hand side of Eq. (31) on the “past” values of $`Z`$ (i.e., on its values at times preceding $`t`$) can also be resolved in the framework of the small-$`\tau _\mathrm{c}`$ expansion. Formally integrating Eq. (31), we get, at times $`t_1<t`$, $`Z(t_1)=Z(t)\widehat{L}{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^t^{}}dt_2\kappa (t^{}t_2)Z(t_2)+\mathrm{}`$ (32) Upon substituting this onto the right-hand side of Eq. (31) and again discarding all the terms of orders higher than the first in $`\tau _\mathrm{c}`$, we get $`_tZ(t)`$ $`=`$ $`\widehat{L}{\displaystyle _0^t}dt_1\kappa (tt_1)Z(t)`$ (33) $``$ $`\widehat{L}^2{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^t^{}}dt_2\kappa (tt_1)\kappa (t^{}t_2)Z(t)`$ (34) $`+`$ $`(\widehat{L}^2\widehat{L}_1){\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^{t_1}}dt_2\kappa (tt_1)\kappa (t^{}t_2)Z(t)`$ (35) $`+`$ $`(\widehat{L}^2\widehat{L}_2){\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _{t_1}^t^{}}dt_2\kappa (tt_1)\kappa (t^{}t_2)Z(t).`$ (36) We have introduced the following two operators: $`\widehat{L}_1`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_i^kT_{k\alpha _1}^{i\beta _1}[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\widehat{\mathrm{\Lambda }}_m^n]T_{n\alpha _2}^{m\beta _2}\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}={\displaystyle \frac{d^2}{d+1}}\widehat{L},`$ (37) $`\widehat{L}_2`$ $`=`$ $`\widehat{\mathrm{\Lambda }}_i^kT_{k\alpha _1}^{i\beta _1}[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\widehat{L}]=2d\widehat{L},`$ (38) where the square brackets denote commutators. We see that the terms in Eq. (36) that contain $`\widehat{L}^2`$ cancel out, and only those terms remain that are due to the non-self-commuting nature of the operator $`\widehat{\mathrm{\Lambda }}_i^k`$. Finally, we inverse-Fourier transform Eq. (36) into the $`𝐁`$ space and take the long-time limit, $`t\tau _\mathrm{c}`$. The following Fokker–Planck equation with constant coefficients results: $`_tP={\displaystyle \frac{\overline{\kappa }}{2}}\left[1\tau _\mathrm{c}\overline{\kappa }d\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{d+1}}K_1+K_2\right)\right]\widehat{L}P,`$ (39) where the coefficients, $`K_1`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt_2{\displaystyle _0^{t_1}}dt_3\kappa (tt_1)\kappa (t_2t_3),`$ (40) $`K_2`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt_2{\displaystyle _{t_1}^{t_2}}dt_3\kappa (tt_1)\kappa (t_2t_3),`$ (41) are constants that depend on the particular shape of the time-correlation function $`\kappa (tt^{})`$. Thus, the $`\tau `$ expansion is nonuniversal in the sense that the specific choice of the functional form of the small-time regularization directly affects the values of the expansion coefficients (cf. Ref. ). As an example, let us give the values of the coefficients $`K_1`$ and $`K_2`$ for two popular choices of $`\kappa (tt^{})`$: $`\kappa (tt^{})={\displaystyle \frac{\overline{\kappa }}{2\tau _\mathrm{c}}}\mathrm{exp}\left[\|tt^{}\|/\tau _\mathrm{c}\right]`$ $``$ $`K_1=K_2=0.5`$ (42) $`\kappa (tt^{})={\displaystyle \frac{\overline{\kappa }}{\sqrt{\pi }\tau _\mathrm{c}}}\mathrm{exp}\left[(tt^{})^2/\tau _\mathrm{c}^2\right]`$ $``$ $`K_10.33,K_20.23.`$ (43) In Sec. III, we will apply the method we have presented above to the more realistic general compressible kinematic dynamo problem in the Eulerian frame. Remark on the Kliatskin–Tatarskii method. The Gaussian hierarchy given by the equations (12) and (13) and based on repeated application of the Furutsu–Novikov formula was proposed by Kliatskin and Tatarskii as a basis for constructing successive-approximation solutions of the problem of light propagation in a medium with randomly distributed inhomogeneities. Their method in its original form was carried over to the mean-field dynamo theory with finite-time-correlated velocity field by Vainshtein . The method we have outlined in this section, while also based on the response-function hierarchy (12)–(13), differs substantially from that developed and applied by these authors. Their successive-approximation scheme consisted essentially in writing out the first $`n`$ equations in the hierarchy (12)–(13) and then truncating it at the $`n`$th step by replacing $`\kappa (tt_{n+1})`$ by $`\overline{\kappa }\delta (tt_{n+1})`$ in the equation for the $`n`$th-order response function. This gave a closed system of equations that could be solved. Carried out in the first order, such a procedure would correspond to setting $`\kappa (t^{}t_2)=\overline{\kappa }\delta (t^{}t_2)`$ in Eq. (21) and consequently $`\kappa (t_2t_3)=\overline{\kappa }\delta (t_2t_3)`$ in the expressions (40) and (41) for the coefficients $`K_1`$ and $`K_2`$. Such a substitution leads to $`K_1=0`$ and $`K_2=(2/\tau _\mathrm{c}\overline{\kappa })_0^{\mathrm{}}d\tau \tau \kappa (\tau )`$, which is incorrect. The reason for this discrepancy is that, in the time integrals involving multiple products of the correlation functions $`\kappa (tt_1)`$, $`\kappa (t_2t_3)`$, etc., the latter cannot be approximated by $`\delta `$ functions plus first-order corrections even in the small-$`\tau _\mathrm{c}`$ limit. ### D Comparison with the Van Kampen Cumulant Expansion Method The evolution equation $`(\text{6})`$ for the “unaveraged characteristic function” $`\stackrel{~}{Z}(t;\mu )`$ is a stochastic linear differential equation whose form agrees exactly with that of the general such equation considered by van Kampen and simultaneously by Terwiel : $`_t\stackrel{~}{Z}(t)=\widehat{A}(t)\stackrel{~}{Z}(t),`$ (44) where, in our case, $`\widehat{A}(t)=\widehat{\mathrm{\Lambda }}_i^k\sigma _k^i(t)`$. In his work, van Kampen developed a formalism that allowed one to construct successive terms in the short-correlation-time expansion of $`Z=\stackrel{~}{Z}`$ in terms of the cumulants of the operator $`\widehat{A}`$. Terwiel’s projection-operator method was shown by its author to be equivalent to that of van Kampen. Let us see what happens if the passive-advection problem given by Eq. (1), or, equivalently, by Eq. (6) is subjected to van Kampen’s expansion algorithm. The latter proceeds as follows. Start by writing the formal solution of Eq. (6) in terms of the time-ordered exponential: $`\stackrel{~}{Z}(t)=\mathrm{exp}{\displaystyle _0^t}dt^{}\widehat{A}(t^{})\stackrel{~}{Z}(0)`$ (45) $`=\left[1+{\displaystyle _0^t}dt_1\widehat{A}(t_1)+{\displaystyle _0^t}dt_1{\displaystyle _0^{t_1}}dt_2\widehat{A}(t_1)\widehat{A}(t_2)+\mathrm{}\right]\stackrel{~}{Z}(0).`$ (46) This solution is averaged assuming that the initial distribution of $`\stackrel{~}{Z}`$ is independent of the statistics of $`\widehat{A}`$: $`Z(t)=[1+{\displaystyle _0^t}\mathrm{d}t_1{\displaystyle _0^{t_1}}\mathrm{d}t_2<\widehat{A}(t_1)\widehat{A}(t_2)>.`$ (47) $`+.{\displaystyle _0^t}\mathrm{d}t_1{\displaystyle _0^{t_1}}\mathrm{d}t_2{\displaystyle _0^{t_2}}\mathrm{d}t_3{\displaystyle _0^{t_3}}\mathrm{d}t_4<\widehat{A}(t_1)\widehat{A}(t_2)\widehat{A}(t_3)\widehat{A}(t_4)>+\mathrm{}]Z(0).`$ (48) Here all the odd-order averages have vanished (recall that $`\widehat{A}=\widehat{\mathrm{\Lambda }}_i^k\sigma _k^i`$). The closed equation for $`Z(t)`$ is now obtained as follows. First, the formal solution (48) is differentiated with respect to time: $`_tZ(t)=[{\displaystyle _0^t}\mathrm{d}t_1<\widehat{A}(t)\widehat{A}(t_1)>.`$ (49) $`.+{\displaystyle _0^t}\mathrm{d}t_1{\displaystyle _0^{t_1}}\mathrm{d}t_2{\displaystyle _0^{t_2}}\mathrm{d}t_3<\widehat{A}(t)\widehat{A}(t_1)\widehat{A}(t_2)\widehat{A}(t_3)>+\mathrm{}]Z(0).`$ (50) Second, $`Z(0)`$ is expressed in terms of $`Z(t)`$ by formally inverting the operator series on the right-hand side of Eq. (48), whereupon $`Z(0)`$ is substituted into Eq. (50). Keeping only the terms that contain up to three time integrations, as we did in the previous section, we get $`_tZ(t)`$ $`=`$ $`[{\displaystyle _0^t}\mathrm{d}t_1<\widehat{A}(t)\widehat{A}(t_1)>.`$ (51) $`+`$ $`{\displaystyle _0^t}dt_1{\displaystyle _0^{t_1}}dt_2{\displaystyle _0^{t_2}}dt_3<\widehat{A}(t)\widehat{A}(t_1)\widehat{A}(t_2)\widehat{A}(t_3)>`$ (52) $``$ $`.{\displaystyle _0^t}\mathrm{d}t_1{\displaystyle _0^t}\mathrm{d}t_2{\displaystyle _0^{t_2}}\mathrm{d}t_3<\widehat{A}(t)\widehat{A}(t_1)><\widehat{A}(t_2)\widehat{A}(t_3)>+\mathrm{}]Z(t).`$ (53) The quadruple average in the above expression splits into three products of second-order averages in the usual Gaussian way. Since $`<\widehat{A}(t)\widehat{A}(t_1)>=\kappa (tt_1)T_{kl}^{ij}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_j^l=\kappa (tt_1)\widehat{L},`$ (54) we have $`<\widehat{A}(t)\widehat{A}(t_1)\widehat{A}(t_2)\widehat{A}(t_3)>`$ $`=`$ $`\kappa (tt_1)\kappa (t_2t_3)T_{kl}^{ij}T_{pq}^{mn}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_j^l\widehat{\mathrm{\Lambda }}_m^p\widehat{\mathrm{\Lambda }}_n^q`$ (55) $`+`$ $`\kappa (tt_2)\kappa (t_1t_3)T_{kp}^{im}T_{lq}^{jn}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_j^l\widehat{\mathrm{\Lambda }}_m^p\widehat{\mathrm{\Lambda }}_n^q`$ (56) $`+`$ $`\kappa (tt_3)\kappa (t_1t_2)T_{kq}^{in}T_{lp}^{jm}\widehat{\mathrm{\Lambda }}_i^k\widehat{\mathrm{\Lambda }}_j^l\widehat{\mathrm{\Lambda }}_m^p\widehat{\mathrm{\Lambda }}_n^q`$ (57) $`=`$ $`\kappa (tt_1)\kappa (t_2t_3)\widehat{L}^2`$ (58) $`+`$ $`\kappa (tt_2)\kappa (t_1t_3)\left(\widehat{L}^2\widehat{L}_1\right)`$ (59) $`+`$ $`\kappa (tt_3)\kappa (t_1t_2)\left(\widehat{L}^2\widehat{L}_2\right),`$ (60) $`<\widehat{A}(t)\widehat{A}(t_1)><\widehat{A}(t_2)\widehat{A}(t_3)>`$ $`=`$ $`\kappa (tt_1)\kappa (t_2t_3)\widehat{L}^2.`$ (61) The operators $`\widehat{L}`$, $`\widehat{L}_1`$, and $`\widehat{L}_2`$ are the same as those in the previous section \[see definitions (16), (37), and (38)\]. The averages (54), (60), and (61) are now substituted into Eq. (53). The triple time integrals can be argued to represent (all of the) first-order terms in the small-$`\tau _\mathrm{c}`$ expansion in the same way as it was done in Sec. II C. In the limit $`t\tau _\mathrm{c}`$, the coefficients in Eq. (53) do not depend on time. The resulting Fokker–Planck equation for the PDF $`P(t;B)`$ \[which is the inverse Fourier transform of $`Z(t;\mu )`$\] obtained by the van Kampen method and analogous to Eq. (39) is then $`_tP={\displaystyle \frac{\overline{\kappa }}{2}}\left[1\tau _\mathrm{c}\overline{\kappa }d\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{d+1}}C_1+C_2\right){\displaystyle \frac{1}{2}}\tau _\mathrm{c}\overline{\kappa }\left(C_0C_1C_2\right)\widehat{L}\right]\widehat{L}P,`$ (62) where the coefficients $`C_0`$, $`C_1`$, and $`C_2`$, that depend on the shape function $`\kappa (tt^{})`$, are as follows: $`C_0`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt_2{\displaystyle _0^{t_2}}dt_3\kappa (tt_1)\kappa (t_2t_3),`$ (63) $`C_1`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _0^{t_1}}dt_2{\displaystyle _0^{t_2}}dt_3\kappa (tt_2)\kappa (t_1t_3),`$ (64) $`C_2`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _0^{t_1}}dt_2{\displaystyle _0^{t_2}}dt_3\kappa (tt_3)\kappa (t_1t_2).`$ (65) By comparing the definition of $`C_0`$ with those of the coefficients $`K_1`$ and $`K_2`$ in Sec. II C \[see formulas (40) and (41)\], we immediately establish that $`C_0=K_1+K_2`$. Furthermore, it is also not hard to ascertain that $`C_1=K_1`$ and $`C_2=K_2`$. Therefore, the last term in Eq. (62) vanishes, and the first-order Fokker–Planck equations (39) and (62) are identical. Thus, the results obtained via the van Kampen method are consistent with ours. Unlike the van Kampen method, however, our method does not involve any nontrivial operator algebra and is therefore better suited for a wide variety of applications. In particular, the stochastic equations containing spatial derivatives (such as the convective derivatives present in all Eulerian passive-advection problems) can be handled without much additional difficulty (this will be done in detail for the full kinematic dynamo problem in Sec. III). ### E Discussion: The Vertex Corrections While the particular methods one employs to obtain the successive terms in the $`\tau `$ expansion may vary and depend on one’s taste and the specific demands of the stochastic problem at hand, the underlying structure of the $`\tau `$ expansion remains the same and is rooted in the common properties of all turbulence closure problems (see, e.g., Ref. ). As we have stated in general terms in the introduction to this paper, and as was clear from our construction of the response-function formalism in Sec. II A–Sec. II C or of van Kampen’s explicit series solution (48) in Sec. II D, averaged solutions of stochastic equations such as Eq. (6) can be represented in terms of infinite sums of multiple time-history integrals containing products of time-correlation functions $`\kappa (t_it_j)`$ in the integrands. This summation can be visualized in terms of Feynman-style diagrams. The $`n`$-point diagrams represent the terms containing $`n`$ time-history integrations. As an example, Fig. 1 lists the three possible four-point diagrams. It was noted by Kazantsev (see also Ref. ) that the white-noise approximation corresponds to the partial summation of all ladder-type diagrams such as the four-point one shown in Fig. 1(a). The distinctive property of these diagrams is that the pairs of points $`t_i`$$`t_j`$ at which the time-correlation functions in the integrands of the time-history integrals are taken, can be fused without interfering with each other. No essential information is therefore lost when the time-correlation functions $`\kappa (t_it_j)`$ are approximated by $`\delta `$ functions. However, in all orders of the $`\tau `$ expansion but the zeroth, diagrams with more tangled topology appear: e.g., in the first order, these are the diagrams 1(b) and 1(c)\]. Such diagrams are often referred to as the vertex corrections. Fusing points in these diagrams leads to the loss of terms that cannot be neglected . This is the context in which the emerging nonuniversality with respect to the shape of the time-correlation profile should be viewed. In this paper, we restrict our consideration to the first-order terms in the $`\tau `$ expansion. The relevant diagrams are the four-point ones shown in Fig. 1. The diagrams 1(b) and 1(c) give rise to the coefficients $`C_1`$ and $`C_2`$, respectively \[see Eq. (62) and formulas (64) and (65)\]. Upon changing variables $`t_1t_2`$ in the diagram 1(b) and $`t_1t_2`$, $`t_2t_3`$, $`t_3t_1`$ in the diagram 1(c), we see that these diagrams equally well correspond to the coefficients $`K_1`$ and $`K_2`$ \[Eq. (39) and formulas (40) and (41)\]. ## III The Functional Expansion for the Kinematic Dynamo in a Finite-Time-Correlated Velocity Field In this Section, we use the functional expansion method developed in Sec. II to construct the $`\tau `$ expansion for the general diffusion-free kinematic dynamo problem in the Eulerian frame with an arbitrarily compressible velocity field. Through the convective derivative, an explicit spatial dependence is now present in the problem. This leads to the appearance of the new effect advertised in the Introduction: while the zeroth-order terms in the expansion only depend on the one-point correlation properties of the velocity gradients, the first-order terms also depend on the energy of the advecting velocity field and on the one-point correlation function of its second derivatives. The former represents the loss of Galilean invariance, the latter the loss of the small-scale universality and the advent of the sensitive dependence of the statistics on the large-scale structure of the velocity correlations. In this Section, all statistics are Eulerian. For the questions regarding the transformation of PDF’s of passive fields from the Eulerian to the Lagrangian frame, we address the reader to Ref. , as well as to Ref. , where the $`\tau `$ expansion is treated as a problem in stochastic calculus and Lagrangian statistics are discussed. ### A The Gaussian Hierarchy The magnetic field passively advected by the velocity field $`u^i(t,𝐱)`$ evolves according to the Hertz induction equation (formally in $`d`$ dimensions): $`_tB^i=u^kB_{,k}^i+u_{,k}^iB^ku_{,k}^kB^i,`$ (66) where $`u_{,k}^i=u^i/x^k`$, $`B_{,k}^i=B^i/x^k`$, and the Einstein summation convention is used throughout. Let the advecting velocity field $`u^i(t,𝐱)`$ be a homogeneous and isotropic Gaussian random field whose statistics are defined by its second-order correlation tensor: $`<u^i(t,𝐱)u^j(t^{},𝐱^{})>=\kappa ^{ij}(tt^{},𝐱𝐱^{}),`$ (67) where, as a function of the time separation $`tt^{}`$, the correlator $`\kappa ^{ij}`$ is assumed to have a finite width $`\tau _\mathrm{c}`$, which we will call the velocity correlation time. As we will only study the one-point statistics of the magnetic field, all relevant information about the velocity correlation properties is contained in the Taylor expansion of $`\kappa ^{ij}`$ around the origin: $`\kappa ^{ij}(\tau ,𝐲)=\kappa _0(\tau )\delta ^{ij}`$ $``$ $`{\displaystyle \frac{1}{2}}\kappa _2(\tau )\left[y^2\delta ^{ij}+2ay^iy^j\right]`$ (68) $`+`$ $`{\displaystyle \frac{1}{4}}\kappa _4(\tau )y^2\left[y^2\delta ^{ij}+2by^iy^j\right]+\mathrm{},`$ (69) as $`y0`$. Here $`a`$ and $`b`$ are the compressibility parameters. Between the purely incompressible and the purely irrotational cases, they vary in the intervals $`{\displaystyle \frac{1}{d+1}}a1,{\displaystyle \frac{2}{d+3}}b2.`$ (70) We should like to mention here that the choice of the coefficients of the small-scale expansion (69) of the velocity correlation tensor is, strictly speaking, not entirely unconstrained. As $`\kappa ^{ij}(\tau ,𝐲)`$ is a correlation function, it must be an inverse Fourier transform of a proper correlation function in the Fourier space . In Appendix A, we give the expressions for the coefficients of the expansion (69) in terms of the spectral characteristics of the velocity field. We further note that, while $`a`$ and $`b`$ can certainly be functions of $`\tau `$, we will not overly shrink the limits of physical generality by assuming that they are either constant or slowly-varying functions of time, i.e. that they do not change appreciably over one correlation time. In order to determine the one-point statistics of the magnetic field, we follow the standard procedure and introduce the characteristic function of $`B^i(t,𝐱)`$ at an arbitrary fixed spatial point $`𝐱`$: $`Z(t;\mu )=<\stackrel{~}{Z}(t,𝐱;\mu )>=<\mathrm{exp}\left[i\mu _iB^i(t,𝐱)\right]>.`$ (71) As usual, the angle brackets denote ensemble averages and the overtildes mark unaveraged quantities. The function $`Z`$ is the Fourier transform of the PDF of the vector elements $`B^i`$. Due to spatial homogeneity, $`Z`$ does not depend on the point $`𝐱`$, where $`B^i(t,𝐱)`$ is taken. Upon differentiating $`\stackrel{~}{Z}`$ with respect to time and using Eq. (66), we get $`_t\stackrel{~}{Z}=u^k\stackrel{~}{Z}_{,k}+u_{,k}^i\mu _i{\displaystyle \frac{}{\mu _k}}\stackrel{~}{Z}u_{,k}^k\mu _l{\displaystyle \frac{}{\mu _l}}\stackrel{~}{Z}=u^k\stackrel{~}{Z}_{,k}+\widehat{\mathrm{\Lambda }}_i^ku_{,k}^i\stackrel{~}{Z},`$ (72) where, for the sake of future convenience, we introduce the operator $`\widehat{\mathrm{\Lambda }}_i^k=\mu _i{\displaystyle \frac{}{\mu _k}}\delta _i^k\mu _l{\displaystyle \frac{}{\mu _l}},`$ (73) which will turn up repeatedly in this calculation. In order to obtain an evolution equation for the characteristic function $`Z(t;\mu )`$ of the random magnetic field, we average both sides of Eq. (72). Since, due to the homogeneity of the problem, $`u^k\stackrel{~}{Z}_{,k}=u_{,k}^k\stackrel{~}{Z}`$, we may write the equation for $`Z(t;\mu )`$ in the following form: $`_tZ(t)=\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)<u_{,k}^i(t,𝐱)\stackrel{~}{Z}(t,𝐱)>`$ (74) $`=\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right){\displaystyle _0^t}dt_1{\displaystyle \mathrm{d}^dx_1\kappa _{,k\alpha _1}^{i\beta _1}(tt_1,𝐱𝐱_1)G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)},`$ (75) where the mixed average on the right-hand side has been “split” with the aid of the Furutsu–Novikov (“Gaussian-integration”) formula , and the $`\mu `$ dependence in the arguments has been suppressed for the sake of notational compactness. We have introduced the first-order averaged response function of the following species: $`G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)=<\stackrel{~}{G}_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)>={\displaystyle \frac{\delta \stackrel{~}{Z}(t,𝐱)}{\delta u_{,\alpha _1}^{\beta _1}(t_1,𝐱_1)}}.`$ (76) As a response function, $`G_{\beta _1}^{\alpha _1}`$ satisfies the causality constraint: $`G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)=0`$ for $`t_1>t`$. The same-time form of $`G_{\beta _1}^{\alpha _1}`$ can be obtained in terms of the characteristic function $`Z(t)`$: integrating Eq. (72) from $`0`$ to $`t_1`$, taking the functional derivative $`\delta /\delta u_{,\alpha _1}^{\beta _1}(t^{},𝐱_1)`$, averaging, setting $`t^{}=t_1`$, and taking causality into account, we get $`G_{\beta _1}^{\alpha _1}(t_1,𝐱\|t_1,𝐱_1)=\delta (𝐱𝐱_1)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1).`$ (77) In order to find $`G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)`$ at $`t>t_1`$, we take the functional derivative $`\delta /\delta u_{,\alpha _1}^{\beta _1}(t^{},𝐱^{})`$ of both sides of Eq. (72) and establish that each element of the unaveraged tensor $`\stackrel{~}{G}_{\beta _1}^{\alpha _1}`$ satisfies an equation identical in form to Eq. (72): $`_t\stackrel{~}{G}_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)=u^m(t,𝐱)\stackrel{~}{G}_{\beta _1,m}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)+\widehat{\mathrm{\Lambda }}_m^nu_{,n}^m(t,𝐱)\stackrel{~}{G}_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1).`$ (78) Subscripts such as “<sub>,m</sub>” in the above equation mean, in accordance with the usual notation, the partial differentiation with respect to $`x^m`$, viz., $`/x^m`$. We must now average Eq. (78) in its turn, to obtain an evolution equation for $`G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)`$ at $`t>t_1`$. Using the initial condition (77), let us write this evolution equation in the integral form valid for all $`tt_1`$: $`G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)=\delta (𝐱𝐱_1)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1)`$ (79) $`{\displaystyle _{t_1}^t}\mathrm{d}t^{}{\displaystyle _0^t^{}}\mathrm{d}t_2{\displaystyle }\mathrm{d}^dx_2[\kappa ^{m\beta _2}(t^{}t_2,𝐱𝐱_2)G_{\beta _1\beta _2,m}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2).`$ (80) $`+.\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,𝐱𝐱_2)G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)],`$ (81) where the mixed averages have again been “split” by the Furutsu–Novikov formula, at the price of introducing two new second-order response functions: $`G_{\beta _1\beta _2}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ $`=`$ $`{\displaystyle \frac{\delta ^2\stackrel{~}{Z}(t^{},𝐱)}{\delta u_{,\alpha _1}^{\beta _1}(t_1,𝐱_1)\delta u^{\beta _2}(t_2,𝐱_2)}},`$ (82) $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ $`=`$ $`{\displaystyle \frac{\delta ^2\stackrel{~}{Z}(t^{},𝐱)}{\delta u_{,\alpha _1}^{\beta _1}(t_1,𝐱_1)\delta u_{,\alpha _2}^{\beta _2}(t_2,𝐱_2)}}.`$ (83) In the same way that the equal-time first-order response function was expressed in terms of $`Z(t)`$ \[Eq. (77)\], the second-order response functions at $`t^{}=t_1`$ or $`t^{}=t_2`$ can be expressed in terms of $`G_{\beta _2}^{\alpha _2}(t_1,𝐱\|t_2,𝐱_2)`$ or $`G_{\beta _1}^{\alpha _1}(t_2,𝐱\|t_1,𝐱_1)`$, respectively. Because of causality, the former representation would be valid provided $`t_1t_2`$, the latter in the opposite case $`t_2t_1`$. At other times, $`t_1,t_2t^{}`$, the functions $`G_{\beta _1\beta _2}^{\alpha _1}`$ and $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}`$ satisfy integral equations analogous to Eq. (81), where third-order response functions make their appearance. An infinite open hierarchy can thus be obtained by further iterating this procedure and introducing response functions of ascending orders. This hierarchy constitutes the exact description of the statistics of the kinematic dynamo problem with arbitrary velocity correlation time. ### B The $`\tau `$ Expansion The expansion in small correlation time must be carried out in such a way that the time integral of the velocity correlator $`\kappa ^{ij}(\tau ,𝐲)`$ remains constant. Since $`\kappa ^{ij}(\tau ,𝐲)`$ has a finite (small) width $`\tau _\mathrm{c}`$, we can conclude that the double time integral on the right-hand side of Eq. (81) must be of first order in the correlation time $`\tau _\mathrm{c}`$. As we are only interested in constructing the $`\tau `$ expansion up to first order, it is now sufficient to calculate the second-order response functions $`G_{\beta _1\beta _2}^{\alpha _1}`$ and $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}`$ with zeroth-order precision. We have already mentioned that recursive relations completely analogous to the relation (81) can be derived for the second-order response functions. The latter are thereby expressed as their equal-time values plus double time integrals of the same sort as that which appeared on the right-hand side of Eq. (81). These time integrals are first order in the correlation time and can therefore be neglected. The equal-time values of the second-order response functions are obtained by formally integrating Eq. (72), taking functional derivatives of it, averaging, and using causality. The second-order response functions are thus expressed to zeroth order in terms of the first-order ones. These latter can by the same token be replaced by their equal-time values, which only contain the characteristic function $`Z(t)`$. The resulting expressions, valid to zeroth order, must be substituted into the first-order term (the double time integral) in Eq. (81). All these manipulations, which require a fair amount of algebra, are relegated to Appendix B. Here we simply give the resulting expression for the first-order response function, valid to first order in $`\tau _\mathrm{c}`$: $`G_{\beta _1}^{\alpha _1}(t,𝐱\|t_1,𝐱_1)=\delta (𝐱𝐱_1)[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1).`$ (84) $`{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^{t_1}}dt_2\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\left(\delta _m^n+L_m^n\right)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2)`$ (85) $`{\displaystyle _{t_1}^t}dt^{}{\displaystyle _{t_1}^t^{}}dt_2\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\left(\delta _m^n+L_m^n\right)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1)`$ (86) $`.+{\displaystyle _{t_1}^t}\mathrm{d}t^{}{\displaystyle _0^{t_1}}\mathrm{d}t_2\kappa _{,\beta _1\alpha _2}^{\alpha _1\beta _2}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2)]`$ (87) $`+{\displaystyle \frac{^2\delta (𝐱𝐱_1)}{x^mx^n}}{\displaystyle _{t_1}^t}dt^{}{\displaystyle _{t_1}^t^{}}dt_2\kappa ^{mn}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1)+𝒪(\tau _\mathrm{c}^2).`$ (88) This expression must now be substituted into the time-history integral on the right-hand side of Eq. (75). This gives a closed integro-differential equation for the characteristic function $`Z(t)`$. However, the dependence on the past values of $`Z`$ is spurious and can be resolved to first order in $`\tau _\mathrm{c}`$. Indeed, we can formally integrate Eq. (75) from $`t_1`$ to $`t`$ and, using the zeroth-order value of the first-order response function \[the first term in the formula (88)\], get $`Z(t_1)=Z(t)+\left(\delta _m^n+\widehat{\mathrm{\Lambda }}_m^n\right){\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^t^{}}dt_2\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2)+𝒪(\tau _\mathrm{c}^2).`$ (89) The double time integral in this equation is of first order in $`\tau _\mathrm{c}`$, as usual. Upon assembling the equations (75), (88), and (89), we finally arrive at the following closed partial differential equation for $`Z(t)`$: $`_tZ(t)={\displaystyle _0^t}dt_1\kappa _{,k\alpha _1}^{i\beta _1}(tt_1,0)\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t)`$ (90) $`{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^{t_1}}dt_2\kappa _{,k\alpha _1}^{i\beta _1}(tt_1,0)`$ (91) $`\times \kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\widehat{\mathrm{\Lambda }}_m^n]\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t)`$ (92) $`{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _{t_1}^t^{}}dt_2\kappa _{,k\alpha _1}^{i\beta _1}(tt_1,0)`$ (93) $`\times \kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\left(\delta _m^n+\widehat{\mathrm{\Lambda }}_m^n\right)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}]Z(t)`$ (94) $`{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _0^{t_1}}dt_2\kappa _{,k\alpha _1}^{i\beta _1}(tt_1,0)\kappa _{,\beta _1\alpha _2}^{\alpha _1\beta _2}(t^{}t_2,0)\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t)`$ (95) $`{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt^{}{\displaystyle _{t_1}^t^{}}dt_2\kappa _{,k\alpha _1mn}^{i\beta _1}(tt_1,0)\kappa ^{mn}(t^{}t_2,0)\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t).`$ (96) The square brackets denote commutators. Note that, besides the second derivatives of the velocity correlation tensor, the first-order terms contain the fourth ones, as well as the undifferentiated tensor itself \[in the last term in Eq. (96)\]. The latter implies the loss of Galilean invariance, the former the loss of small-scale universality in the sense that the large-scale structure of the velocity correlator starts to play a role. This effect could not have been captured if the velocity field had been assumed to be purely a combination of the instantaneous velocity at a given point and a linear shear. ### C The Fokker–Planck Equation In order to obtain the Fokker–Planck equation for the PDF of the magnetic field, we must inverse-Fourier transform Eq. (96) back to $`𝐁`$ dependence. The inverse Fourier transform of $`Z(t;\mu )`$ is the one-point PDF $`P(t;𝐁)`$. We will continue using the symbol $`\widehat{\mathrm{\Lambda }}_i^k`$ to denote the counterpart of the operator $`\widehat{\mathrm{\Lambda }}_i^k`$ in the $`𝐁`$ space: $`\widehat{\mathrm{\Lambda }}_i^k=(d1)\delta _i^kB^k{\displaystyle \frac{}{B^i}}+\delta _i^kB^l{\displaystyle \frac{}{B^l}}.`$ (97) Due to the isotropy of the problem, the PDF $`P(t;𝐁)`$ will in fact be a scalar function of the field strength $`B`$ only. Thus, all the operators that appear on the right-hand side of the $`𝐁`$-space counterpart of Eq. (96) must, after they are convolved with the velocity correlation tensors, be expressible in terms of $`B`$. Let us use the Taylor expansion (69) of the velocity correlator to calculate the tensor convolutions in Eq. (96). We have $`\kappa ^{ij}(\tau ,0)`$ $`=`$ $`\kappa _0(\tau )\delta ^{ij},`$ (98) $`\kappa _{,kl}^{ij}(\tau ,0)`$ $`=`$ $`\kappa _2(\tau )\left[\delta ^{ij}\delta _{kl}+a\left(\delta _k^i\delta _l^j+\delta _l^i\delta _k^j\right)\right]=\kappa _2(\tau )T_{kl}^{ij},`$ (99) $`\kappa _{,klmm}^{ij}(\tau ,0)`$ $`=`$ $`\kappa _4(\tau )\left[2(d+2+b)\delta ^{ij}\delta _{kl}+(d+4)b\left(\delta _k^i\delta _l^j+\delta _l^i\delta _k^j\right)\right]`$ (100) $`=`$ $`\kappa _4(\tau )U_{kl}^{ij}.`$ (101) A number of second-order differential operators (with respect to $`B`$) arise in Eq. (96). In the zeroth-order term, we have $`\widehat{L}=T_{k\alpha _1}^{i\beta _1}\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}={\displaystyle \frac{d1}{d+1}}\left(B{\displaystyle \frac{}{B}}+d\right)\left(\left(1+\beta \right)B{\displaystyle \frac{}{B}}+(d+1)\beta \right);`$ (102) two operators appearing in the first-order terms result from the non-self-commuting nature of the operator $`\widehat{\mathrm{\Lambda }}_i^k`$ \[see the second and the third terms in Eq. (96)\] : $`\widehat{L}_1`$ $`=`$ $`T_{k\alpha _1}^{i\beta _1}T_{n\alpha _2}^{m\beta _2}\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\widehat{\mathrm{\Lambda }}_m^n]\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}`$ (103) $`=`$ $`{\displaystyle \frac{d^2(d1)}{(d+1)^2}}\left(B{\displaystyle \frac{}{B}}+d\right)\left(1+{\displaystyle \frac{\beta }{d^2}}\right)^2B{\displaystyle \frac{}{B}},`$ (104) $`\widehat{L}_2`$ $`=`$ $`T_{k\alpha _1}^{i\beta _1}T_{n\alpha _2}^{m\beta _2}\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\left(\delta _m^n+\widehat{\mathrm{\Lambda }}_m^n\right)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}]=T_{k\alpha _1}^{i\beta _1}\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)[\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1},\widehat{L}]`$ (105) $`=`$ $`{\displaystyle \frac{2d(d1)}{d+1}}\left(B{\displaystyle \frac{}{B}}+d\right)\left(1+{\displaystyle \frac{\beta }{d^2}}\right)B{\displaystyle \frac{}{B}};`$ (106) and, finally, there are two other operators due to the presence of the convective term (i.e., explicit spatial dependence) in the induction equation \[see the fourth and the fifth terms in Eq. (96)\]: $`\widehat{M}_1`$ $`=`$ $`T_{k\alpha _1}^{i\beta _1}T_{\beta _1\alpha _2}^{\alpha _1\beta _2}\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}={\displaystyle \frac{d(d1)}{(d+1)^2}}\left(B{\displaystyle \frac{}{B}}+d\right)`$ (108) $`\times \left[\left(1+{\displaystyle \frac{2\beta }{d^2}}+{\displaystyle \frac{d(d+1)1}{d^3}}\beta ^2\right)B{\displaystyle \frac{}{B}}+{\displaystyle \frac{(d+1)^2}{d^2}}\beta \right],`$ $`\widehat{M}_2`$ $`=`$ $`U_{k\alpha _1}^{i\beta _1}\left(\delta _i^k+\widehat{\mathrm{\Lambda }}_i^k\right)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}={\displaystyle \frac{2(d1)(d+4)}{d+3}}\left(B{\displaystyle \frac{}{B}}+d\right)`$ (110) $`\times \left[\left(1+{\displaystyle \frac{d^2+4d+2}{2d(d+4)}}\zeta \right)B{\displaystyle \frac{}{B}}+{\displaystyle \frac{(d+2)(d+3)}{2(d+4)}}\zeta \right].`$ In all of the above, $`\beta =d\left[1+(d+1)a\right]`$ and $`\zeta =d\left[2+(d+3)b\right]`$ are compressibility parameters that vanish in the case of incompressible flow. In this latter case, the operators defined above simplify considerably: $`\widehat{L}_1={\displaystyle \frac{d^2}{d+1}}\widehat{L},`$ $`\widehat{L}_2=2d\widehat{L},`$ (111) $`\widehat{M}_1={\displaystyle \frac{d}{d+1}}\widehat{L},`$ $`\widehat{M}_2={\displaystyle \frac{2(d+1)(d+4)}{d+3}}\widehat{L}.`$ (112) If we take the long-time limit, i.e., $`t\tau _\mathrm{c}`$, the coefficients in Eq. (96) do not depend on time $`t`$. We can now use the inverse Fourier transform of Eq. (96) taken in this limit and the isotropic operators listed above to assemble the Fokker–Planck equation for the PDF of the magnetic field. This equation contains the desired corrections that are of first order in the velocity correlation time $`\tau _\mathrm{c}`$ and represent the first available manifestation of the finite-correlation-time effects. We have $`_tP={\displaystyle \frac{\overline{\kappa }_2}{2}}\left(\widehat{L}{\displaystyle \frac{1}{2}}\tau _\mathrm{c}\overline{\kappa }_2\left[K_1\left(\widehat{L}_1+\widehat{M}_1\right)+K_2\widehat{L}_2+\stackrel{~}{K}_2\widehat{M}_2\right]\right)P,`$ (113) where the overall dimensional factor is $`\overline{\kappa }_2=2{\displaystyle _0^{\mathrm{}}}d\tau \kappa _2(\tau ),`$ (114) and the coefficients , $`K_1`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }_2^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt_2{\displaystyle _0^{t_1}}dt_3\kappa _2(tt_1)\kappa _2(t_2t_3),`$ (115) $`K_2`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }_2^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt_2{\displaystyle _{t_1}^{t_2}}dt_3\kappa _2(tt_1)\kappa _2(t_2t_3),`$ (116) $`\stackrel{~}{K}_2`$ $`=`$ $`{\displaystyle \frac{4}{\tau _\mathrm{c}\overline{\kappa }_2^2}}\underset{t\mathrm{}}{lim}{\displaystyle _0^t}dt_1{\displaystyle _{t_1}^t}dt_2{\displaystyle _{t_1}^{t_2}}dt_3\kappa _4(tt_1)\kappa _0(t_2t_3),`$ (117) are constants that depend on the particular shapes of the time-correlation functions $`\kappa _0(\tau )`$, $`\kappa _2(\tau )`$, and $`\kappa _4(\tau )`$. Such sensitive dependence is a new feature and represents a loss of universality with respect to the specific time-correlation profiles (cf. Ref. ). As we have pointed out in Sec. III B, the universality with respect to the functional form of the velocity correlator in space is also lost (this effect is incorporated into the coefficient $`\stackrel{~}{K}_2`$). Let us also list the much more compact form that the Fokker–Planck equation (113) assumes in the case of an incompressible velocity field: $`_tP={\displaystyle \frac{\overline{\kappa }_2}{2}}\left[1\tau _\mathrm{c}\overline{\kappa }_2d\left({\displaystyle \frac{1}{2}}K_1+K_2+{\displaystyle \frac{(d+1)(d+4)}{d(d+3)}}\stackrel{~}{K}_2\right)\right]\widehat{L}P.`$ (118) Here it is especially manifest that the true expansion parameter in the problem is $`\tau _\mathrm{c}\overline{\kappa }_2d`$. This is a general statement that holds regardless of the degree of compressibility, as can be readily verified by counting powers of $`d`$ in the general expressions for the operators $`\widehat{L}`$, $`\widehat{L}_1`$, $`\widehat{L}_2`$, $`\widehat{M}_1`$, and $`\widehat{M}_2`$ \[formulas (102)–(110)\]. It is evident that the distribution resulting from Eq. (113) is lognormal, which is a well-known fact in the kinematic-dynamo and passive-advection theory. Since we are interested in the quantitative description of the fast-dynamo effect, we will now proceed to calculate the growth rates of the moments of the magnetic field. ### D The Dynamo Growth Rates The evolution of all moments of $`B`$ can be determined from Eq. (113). The $`n`$th moment is calculated according to $`B^n={\displaystyle \frac{2\pi ^{d/2}}{\mathrm{\Gamma }(d/2)}}{\displaystyle _0^{\mathrm{}}}dBB^{n+d1}P(t;B).`$ (119) Upon multiplying both sides of Eq. (113) by $`B^{d+n1}`$ and integrating over $`B`$, we find that $`B^n`$ satisfies: $`_tB^n`$ $`=`$ $`\gamma (n)B^n`$ (120) $`=`$ $`{\displaystyle \frac{\overline{\kappa }_2}{2}}\left\{\mathrm{\Gamma }(n)\tau _\mathrm{c}\overline{\kappa }_2d\left[K_1\mathrm{\Gamma }_1(n)+K_2\mathrm{\Gamma }_2(n)+\stackrel{~}{K}_2\stackrel{~}{\mathrm{\Gamma }}_2(n)\right]\right\}B^n,`$ (121) where the nondimensionalized zeroth-order growth rates are (cf. Ref. ) $`\mathrm{\Gamma }(n)={\displaystyle \frac{d1}{d+1}}n\left[n+d+(n1)\beta \right],`$ (122) and the universal parts of the (negative) first-order corrections arising from the second- and fourth-order terms in the velocity correlator (69) are $`\mathrm{\Gamma }_1(n)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d1}{d+1}}n\left[\left(n+d\right)\left(1+{\displaystyle \frac{2\beta }{d^2}}\right)+\left(n1\right){\displaystyle \frac{\beta ^2}{d^2}}\right],`$ (123) $`\mathrm{\Gamma }_2(n)`$ $`=`$ $`{\displaystyle \frac{d1}{d+1}}n\left(n+d\right)\left(1+{\displaystyle \frac{\beta }{d^2}}\right),`$ (124) $`\stackrel{~}{\mathrm{\Gamma }}_2(n)`$ $`=`$ $`{\displaystyle \frac{(d1)(d+4)}{d(d+3)}}n\left[n+d+{\displaystyle \frac{n}{2}}\left({\displaystyle \frac{d^2+4d+2}{d(d+4)}}1\right)\zeta \right].`$ (125) We observe that, for $`n=0`$, $`\mathrm{\Gamma }=\mathrm{\Gamma }_1=\mathrm{\Gamma }_2=\stackrel{~}{\mathrm{\Gamma }}_2=0`$. This simply means that both zeroth- and first-order terms in the $`\tau `$ expansion preserve the normalization of the PDF, i.e., our expansion is conservative, as it should be. In the incompressible flow, the total growth rate of the $`n`$th moment can be written in a more compact form: $`\gamma (n)={\displaystyle \frac{\overline{\kappa }_2}{2}}{\displaystyle \frac{d1}{d+1}}n(n+d)\left[1\tau _\mathrm{c}\overline{\kappa }_2d\left({\displaystyle \frac{1}{2}}K_1+K_2+{\displaystyle \frac{(d+1)(d+4)}{d(d+3)}}\stackrel{~}{K}_2\right)\right].`$ (126) We see that the corrections to the growth rates of the magnetic-field moments are negative, so the growth rates are reduced. The amount of reduction depends on a variety of factors including the dimension of space, the order of the moment, the degree of compressibility, the functional form of the velocity correlator in time and space, and, of course, the velocity correlation time. Let us note that our general results derived for an arbitrarily compressible velocity field reveal no qualitatively essential effect of compressibility on the behavior of the first-order finite-correlation-time corrections to the dynamo growth rates in the diffusion-free regime. Compressibility of the flow simply leads to additive (and positive) corrections to the incompressible values of $`\mathrm{\Gamma }(n)`$, $`\mathrm{\Gamma }_1(n)`$, $`\mathrm{\Gamma }_2(n)`$, and $`\stackrel{~}{\mathrm{\Gamma }}_2(n)`$. Quantitatively, these corrections may affect the exact conditions for the break-down of the first-order approximation. For more discussion of the compressibility effects in the kinematic dynamo (with a $`\delta `$-correlated velocity field), we address the reader to Refs. . We remind the reader that here we have studied magnetic fluctuations in the diffusion-free regime and therefore dropped the term in the induction equation that is responsible for the resistive regularization. Such an approach is justified for plasmas with very large magnetic Prandtl numbers (e.g., the ISM or the prototogalaxy) and applies to the initial stage of the small-scale dynamo that lasts for a time of order $`t\mathrm{log}\mathrm{Pr}`$ that elapses before the magnetic fluctuations reach resistive scales . After that, or if the Prandtl number is of order unity or small (as is, e.g., the case for the Sun), resistive effects must be taken into account. In this case, the calculation of the moments of the magnetic field via the Fokker–Planck equation for its PDF as presented in this Section does not apply because of the closure problem associated with the diffusion term \[the equations for $`\stackrel{~}{Z}(t,𝐱;\mu )`$ and $`Z(t;\mu )`$ do not close\]. However, the general $`\tau `$-expansion method proposed in this paper can, in principle, be applied to multipoint correlators of the magnetic field, for which treating the diffusive case presents no conceptual difficulty. One-point moments can then be obtained by fusing the points at which the multipoint correlators are taken (cf. Refs. ). Although it is the diffusive case that is studied in most numerical simulations, where $`\mathrm{Pr}`$ rarely exceeds $`100`$, it is not necessarily the most relevant one in the context of the (proto)galactic dynamo, for which $`\mathrm{Pr}10^{14}÷10^{22}`$. Indeed, as we already pointed out in the Introduction, the initial (proto)galactic seed field may well be strong enough for the kinematic approximation to break down while the dynamo is still in the diffusion-free stage . If this is the case, the effect of magnetic diffusion must be studied in conjunction with nonlinear saturation of the magnetic fluctuations . ## IV A Physical Example: The One-Eddy Model In real astrophysical environments, such as the interstellar medium and the protogalactic plasmas, the magnetic fields are acted upon by a Kolmogorov-like turbulence with a fully developed inertial range about three decades wide ($`\mathrm{Re}10^4`$). While the velocities of the turbulent eddies excited by the Kolmogorov cascade decrease with the scale of the eddy, the velocity gradients increase (see, e.g., Ref. ). Therefore, the dominant role in the process of amplification of the small-scale magnetic fluctuations is played by the smallest eddies. With this circumstance in mind, one often considers, for modeling purposes, a synthetic incompressible turbulent velocity field consisting of eddies all of which have the same fixed size but random isotropic orientation (for detailed discussions of the galactic and protogalactic dynamo, we refer the reader to Refs. ). In this Section, we will present a brief discussion of the implications of the $`\tau `$-expansion theory developed in Sec. III for such a model problem, which will henceforth be referred to as the one-eddy model. The velocity field in the one-eddy model is specified as follows: $`u^i(t,𝐱)`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}e^{i𝐤𝐱}u^i(t,𝐤)},`$ (127) where the Fourier modes $`u^i(t,𝐤)`$ are random variables that satisfy $`<u^i(t,𝐤)u^j(t^{},𝐤^{})>`$ $`=`$ $`(2\pi )^d\delta (𝐤+𝐤^{})\left(\delta ^{ij}{\displaystyle \frac{k_ik_j}{k^2}}\right)\delta (kk_0)\kappa (tt^{}).`$ (128) In this case, $`\kappa _2(\tau )\kappa (\tau )`$, and, upon using the relations listed in Appendix A, we get $`\kappa _0(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{k_0^2}}{\displaystyle \frac{(d1)(d+2)}{d+1}}\kappa _2(\tau ),`$ (129) $`\kappa _4(\tau )`$ $`=`$ $`k_0^2{\displaystyle \frac{d+3}{2(d+4)(d+1)}}\kappa _2(\tau ).`$ (130) Let us specify a plausible velocity time-correlation profile: $`\kappa _2(\tau )={\displaystyle \frac{\overline{\kappa }_2}{2\tau _\mathrm{c}}}\mathrm{exp}\left({\displaystyle \frac{\|\tau \|}{\tau _\mathrm{c}}}\right).`$ (131) For this correlation function, which corresponds, for example, to the well-known Ornstein–Uhlenbeck random process (see, e.g., Ref. ), the coefficients of the $`\tau `$ expansion (126) are $`K_1=K_2=1/2`$. The relations (129) and (130) provide the value of $`\stackrel{~}{K}_2`$: $`\stackrel{~}{K}_2={\displaystyle \frac{(d1)(d+2)(d+3)}{2(d+1)^2(d+4)}}K_2.`$ (132) Let us define the “eddy-turnover” time $`\tau _{\mathrm{eddy}}(k_0u)^1`$ of such a velocity field according to the following relation: $`{\displaystyle \frac{1}{\tau _{\mathrm{eddy}}^2}}=k_0^2{\displaystyle \frac{1}{\tau _\mathrm{c}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}d\tau \kappa ^{ii}(\tau ,𝐲=0)=k_0^2d{\displaystyle \frac{\overline{\kappa }_0}{\tau _\mathrm{c}}}={\displaystyle \frac{d(d1)(d+2)}{d+1}}{\displaystyle \frac{\overline{\kappa }_2}{\tau _\mathrm{c}}},`$ (133) where we have used Eq. (129) to express $`\overline{\kappa }_0`$ in terms of $`\overline{\kappa }_2`$. Note that the same expression is obtained if $`\tau _{\mathrm{eddy}}(𝐮:𝐮)^{1/2}`$ is formally defined in terms of the velocity gradients (without recourse to the one-eddy model): $`{\displaystyle \frac{1}{\tau _{\mathrm{eddy}}^2}}={\displaystyle \frac{1}{\tau _\mathrm{c}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}d\tau \|\kappa _{,jj}^{ii}(\tau ,𝐲=0)\|={\displaystyle \frac{d(d1)(d+2)}{d+1}}{\displaystyle \frac{\overline{\kappa }_2}{\tau _\mathrm{c}}}.`$ (134) We recall that the zeroth-order growth rate $`\gamma _0`$ of the magnetic-fluctuation energy $`B^2`$ is \[see formula (126)\] $`\gamma _0={\displaystyle \frac{(d1)(d+2)}{d+1}}\overline{\kappa }_2.`$ (135) Formulas (133) and (135) then imply $`\left({\displaystyle \frac{\tau _\mathrm{c}}{\tau _{\mathrm{eddy}}}}\right)^2=\tau _\mathrm{c}\gamma _0d={\displaystyle \frac{(d1)(d+2)}{d+1}}\tau _\mathrm{c}\overline{\kappa }_2d.`$ (136) We have established a correspondence between the small parameter that has arisen in our expansion of the dynamo growth rates and the “physical” small parameter, which is the ratio of the correlation and eddy-turnover times. Of course, the above expression hinges on the definitions (133) or (134) of $`\tau _{\mathrm{eddy}}`$. A simple physical argument can be made in favor of these definitions and the resulting formula (136). Namely, let us observe that when $`\tau _\mathrm{c}\tau _{\mathrm{eddy}}`$ the eddy only stretches the magnetic field line in one of the $`d`$ available directions during one turnover time, whence $`\tau _\mathrm{c}\gamma 1/d`$. The same estimate follows from the formula (136). Let us now evaluate the first-order correction to the growth rate of the magnetic energy. In the one-eddy model, one gets, upon using formulas (126) and (132) and taking $`K_1=K_2=1/2`$ for the Ornstein–Uhlenbeck time-correlation profile (131), $`\gamma =\gamma (2)=\gamma _0\left(1C_d\tau _\mathrm{c}\gamma _0d\right),C_d={\displaystyle \frac{2d(d+1)1}{2d(d1)(d+2)}}.`$ (137) We note that in three dimensions, $`C_d=23/6040\%`$. When $`\tau _\mathrm{c}\tau _{\mathrm{eddy}}`$, we have $`\tau _\mathrm{c}\gamma _0d1`$, and the resulting growth-rate reduction of $`40\%`$ is in a good qualitative agreement with the available numerical results . Of course, as we have already stressed in the Introduction, our $`\tau `$ expansion is not designed for the case of $`\tau _\mathrm{c}\tau _{\mathrm{eddy}}`$, so the fact that it gives a fairly reasonable prediction should not be considered as an adequate quantitative corroboration of our theory. At best, one might conclude that the first-order expansion is well behaved for not-too-small values of the expansion parameter. Let us emphasize, however, that such a well-behaved expression has resulted from a number of essentially arbitrary (albeit physically reasonable) specifications of the parameters involved in the $`\tau `$ expansion. One of the most physically important points that we have tried to make in this work is, in fact, that the inclusion of finite-correlation-time effects leads to nonuniversal statistics, so the quantitative predictions of the theory can and will change appreciably if such factors as the shapes of the time-correlation profiles are changed. Namely, one would obtain expressions of the form (137) with different values of the coefficient $`C_d`$. For sufficiently large values of $`\tau _\mathrm{c}\gamma _0d`$, the validity of the expansion (137) will break down, and the expression in the brackets may even become negative. However, the following heuristic argument can be envisioned in this context. Let us recall that the finite-correlation-time effect was due to the presence of time-history integrals such as those that appear in the equations (75), (81), and (88). The first-order corrections in the Fokker–Planck equation (113) arose from systematically approximating the time evolution of the statistical quantities \[response functions and characteristic function $`Z(t;\mu )`$\] that entered these time-history integrals. The corrected (“true”) value of $`\gamma `$ represents, in a rough way, the rate at which these quantities change. It would appear then that a better estimate of $`\gamma `$ would be obtained if $`\gamma _0`$ in the first-order term in the brackets in Eq. (137) were replaced with the corrected value $`\gamma `$. With this caveat, we would find that $`\gamma ={\displaystyle \frac{\gamma _0}{1+C_d\tau _\mathrm{c}\gamma _0d}}.`$ (138) To first order, this formula is equally accurate as Eq. (137). However, it better represents the fact that, as $`\tau _\mathrm{c}\gamma _0d`$ increases, the corrected value of of the growth rate should be expected to saturate . Of course, such considerations cannot substitute for an adequate nonperturbative theory of the passive advection and kinematic dynamo in finite-time-correlated flows, which remains an open problem. ## Acknowledgments The authors would like thank S. A. Boldyrev for extensive and very fruitful discussions of the physics and the formalism of the finite-time-correlated kinematic dynamo problem. Both the substance and the style of the presentation have benefited from suggestions made by J. A. Krommes who read an earlier manuscript of this work. We are also grateful to G. Falkovich, V. Lebedev, S. Cowley, and the anonymous referee for several useful comments. This work was supported by the U. S. Department of Energy under Contract No. DE-AC02-76-CHO-3073. ## A Small-Scale-Expansion Coefficients of the Velocity Correlation Tensor in Terms of Velocity Spectra In this Appendix, we list the basic formulas that relate the coefficients of the small-scale expansion (69) of the velocity correlation tensor to the spectral characteristics of the velocity field. These relations allow one to apply the results on the small-$`\tau _\mathrm{c}`$ expansion obtained in Sec. III to velocity fields that are specified in the Fourier, rather than configuration, space. They also provide a set of consistency constraints that must be respected when the specific functional forms of $`\kappa _0(\tau )`$, $`\kappa _2(\tau )`$, and $`\kappa _4(\tau )`$ are chosen. Let the advecting velocity field be given as a sum of spatial Fourier modes, $`u^i(t,𝐱)`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}e^{i𝐤𝐱}u^i(t,𝐤)}.`$ (A1) and let the Fourier coefficients $`u^i(t,𝐤)`$ be random variables that satisfy $`<u^i(t,𝐤)u^j(t^{},𝐤^{})>`$ $`=`$ $`(2\pi )^d\delta (𝐤+𝐤^{})\left[\kappa (k,tt^{})\delta ^{ij}+\stackrel{~}{\kappa }(k,tt^{}){\displaystyle \frac{k_ik_j}{k^2}}\right].`$ (A2) For the incompressible flows, $`\stackrel{~}{\kappa }(k,\tau )=\kappa (k,\tau )`$; for the irrotational ones, $`\kappa (k,\tau )=0`$. The coefficients of the expansion (69) can then be expressed as follows: $`\kappa _0(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{d}}{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}\left[d\kappa (k,\tau )+\stackrel{~}{\kappa }(k,\tau )\right]},`$ (A3) $`\kappa _2(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}k^2\left[(d+2)\kappa (k,\tau )+\stackrel{~}{\kappa }(k,\tau )\right]},`$ (A4) $`\kappa _4(\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2d(d+2)(d+4)}}{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}k^4\left[(d+4)\kappa (k,\tau )+\stackrel{~}{\kappa }(k,\tau )\right]},`$ (A5) $`a(\tau )`$ $`=`$ $`\kappa _2(\tau )^1{\displaystyle \frac{1}{d(d+2)}}{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}k^2\stackrel{~}{\kappa }(k,\tau )},`$ (A6) $`b(\tau )`$ $`=`$ $`\kappa _4(\tau )^1{\displaystyle \frac{1}{d(d+2)(d+4)}}{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}k^4\stackrel{~}{\kappa }(k,\tau )},`$ (A7) where the $`𝐤`$-space integrations of radial functions can, of course, be written more explicitly as $`{\displaystyle \frac{\mathrm{d}^dk}{(2\pi )^d}}={\displaystyle \frac{S_d}{(2\pi )^d}}{\displaystyle _0^{\mathrm{}}}dkk^{d1},S_d={\displaystyle \frac{2\pi ^{d/2}}{\mathrm{\Gamma }(d/2)}}.`$ (A8) The derivation of the above relations is straightforward and based on the expressions for the correlation functions of isotropic fields in configuration space in terms of their spectra. For the 3-D case, these expressions can be found in Ref. . A detailed derivation of the formulas (A3)-(A7) for the $`d`$-dimensional case is also given in Appendix A of Ref. . ## B Second-Order Response Functions In this Appendix, we provide the zeroth-order expressions for the second-order response functions that we used in Sec. III B. They are all derived in the same fashion: Eq. (72) is formally integrated, functional derivatives of it are taken with respect to the velocity field $`u^i`$ or its gradients $`u_{,k}^i`$ at the appropriate moments, the result is averaged, and the causality property of the response functions is used. Here we simply list the results. When $`t_2t_1`$, we have, to zeroth order in the correlation time $`\tau _\mathrm{c}`$, $`G_{\beta _1\beta _2}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)=G_{\beta _1\beta _2}^{\alpha _1}(t_2,𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ (B1) $`=\delta (𝐱𝐱_2)G_{\beta _1,\beta _2}^{\alpha _1}(t_2,𝐱\|t_1,𝐱_1)+\left[{\displaystyle \frac{}{x^n}}\delta (𝐱𝐱_2)\right]\widehat{\mathrm{\Lambda }}_{\beta _2}^nG_{\beta _1}^{\alpha _1}(t_2,𝐱\|t_1,𝐱_1),`$ (B2) $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)=G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t_2,𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ (B3) $`=\mathrm{\Delta }^{\alpha _2}(𝐱𝐱_2)G_{\beta _1,\beta _2}^{\alpha _1}(t_2,𝐱\|t_1,𝐱_1)+\delta (𝐱𝐱_2)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}G_{\beta _1}^{\alpha _1}(t_2,𝐱\|t_1,𝐱_1),`$ (B4) where we have introduced the following notation: by definition, $`{\displaystyle \frac{\delta u^m(t,𝐱)}{\delta u_{,\alpha _2}^{\beta _2}(t_2,𝐱_2)}}=\delta _{\beta _2}^m\delta (tt_2)\mathrm{\Delta }^{\alpha _2}(𝐱𝐱_2).`$ (B5) The function $`\mathrm{\Delta }^{\alpha _2}(𝐱𝐱_2)`$ is nonrandom and has the following property, which will be all that we need to know about it: $`{\displaystyle \frac{}{x^n}}\mathrm{\Delta }^{\alpha _2}(𝐱𝐱_2)=\delta _n^{\alpha _2}\delta (𝐱𝐱_2).`$ (B6) When $`t_1>t_2`$, the expressions (B2) and (B4) vanish by causality, so we have to flip the order of functional differentiation: $`G_{\beta _1\beta _2}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)=G_{\beta _2\beta _1}^{\alpha _1}(t_1,𝐱\|t_2,𝐱_2;t_1,𝐱_1)`$ (B7) $`=\mathrm{\Delta }^{\alpha _1}(𝐱𝐱_1)G_{\beta _2,\beta _1}(t_1,𝐱\|t_2,𝐱_2)+\delta (𝐱𝐱_1)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}G_{\beta _2}(t_1,𝐱\|t_2,𝐱_2),`$ (B8) $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)=G_{\beta _2\beta _1}^{\alpha _2\alpha _1}(t_1,𝐱\|t_2,𝐱_2;t_1,𝐱_1)`$ (B9) $`=\mathrm{\Delta }^{\alpha _1}(𝐱𝐱_1)G_{\beta _2,\beta _1}^{\alpha _2}(t_1,𝐱\|t_2,𝐱_2)+\delta (𝐱𝐱_1)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}G_{\beta _2}^{\alpha _2}(t_1,𝐱\|t_2,𝐱_2).`$ (B10) In Eq. (B8), the following obvious notation was used: $`G_{\beta _2\beta _1}^{\alpha _1}(t_1,𝐱\|t_2,𝐱_2;t_1,𝐱_1)={\displaystyle \frac{\delta ^2\stackrel{~}{Z}(t_1,𝐱)}{\delta u^{\beta _2}(t_2,𝐱_2)\delta u_{,\alpha _1}^{\beta _1}(t_1,𝐱_1)}},`$ (B11) and a new first-order response function appeared: $`G_{\beta _2}(t_1,𝐱\|t_2,𝐱_2)={\displaystyle \frac{\delta \stackrel{~}{Z}(t_1,𝐱)}{\delta u^{\beta _2}(t_2,𝐱_2)}}.`$ (B12) The equal-time form of this function is $`G_{\beta _2}(t_2,𝐱\|t_2,𝐱_2)=\left[{\displaystyle \frac{}{x^n}}\delta (𝐱𝐱_2)\right]\widehat{\mathrm{\Lambda }}_{\beta _2}^nZ(t_2).`$ (B13) The first-order response functions that appear in the formulas (B2), (B4), (B8), and (B10) can be written as their equal-time values (77) and (B13) plus first-order terms. To zeroth order, we have therefore: for $`t_2t_1`$, $`G_{\beta _1\beta _2}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ $`=`$ $`\left[{\displaystyle \frac{}{x^{\beta _2}}}\delta (𝐱𝐱_1)\right]\delta (𝐱𝐱_2)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1)`$ (B14) $`+`$ $`\delta (𝐱𝐱_1)\left[{\displaystyle \frac{}{x^n}}\delta (𝐱𝐱_2)\right]\widehat{\mathrm{\Lambda }}_{\beta _2}^n\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1),`$ (B15) $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ $`=`$ $`\left[{\displaystyle \frac{}{x^{\beta _2}}}\delta (𝐱𝐱_1)\right]\mathrm{\Delta }^{\alpha _2}(𝐱𝐱_2)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1)`$ (B16) $`+`$ $`\delta (𝐱𝐱_1)\delta (𝐱𝐱_2)\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1);`$ (B17) for $`t_1t_2`$, $`G_{\beta _1\beta _2}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ $`=`$ $`\mathrm{\Delta }^{\alpha _1}(𝐱𝐱_1)\left[{\displaystyle \frac{^2}{x^{\beta _1}x^n}}\delta (𝐱𝐱_2)\right]\widehat{\mathrm{\Lambda }}_{\beta _2}^nZ(t_2)`$ (B18) $`+`$ $`\delta (𝐱𝐱_1)\left[{\displaystyle \frac{}{x^n}}\delta (𝐱𝐱_2)\right]\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}\widehat{\mathrm{\Lambda }}_{\beta _2}^nZ(t_2),`$ (B19) $`G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)`$ $`=`$ $`\mathrm{\Delta }^{\alpha _1}(𝐱𝐱_1)\left[{\displaystyle \frac{}{x^{\beta _1}}}\delta (𝐱𝐱_2)\right]\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2)`$ (B20) $`+`$ $`\delta (𝐱𝐱_1)\delta (𝐱𝐱_2)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2).`$ (B21) These expressions must be substituted into Eq. (81). The volume integrals with respect to $`𝐱_2`$ can be done, taking into account the extremely useful fact that all odd spatial derivatives of the velocity correlator $`\kappa ^{ij}(\tau ,𝐲)`$ vanish at the origin (at $`𝐲=0`$). The results are: for $`t_2t_1`$, $`{\displaystyle \mathrm{d}^dx_2\kappa ^{m\beta _2}(t^{}t_2,𝐱𝐱_2)G_{\beta _1\beta _2,m}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)}`$ (B22) $`=\left[{\displaystyle \frac{^2\delta (𝐱𝐱_1)}{x^{\beta _2}x^m}}\kappa ^{m\beta _2}(t^{}t_2,0)+\delta (𝐱𝐱_1)\kappa _{,mn}^{m\beta _2}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_{\beta _2}^n\right]\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1),`$ (B23) $`{\displaystyle \mathrm{d}^dx_2\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,𝐱𝐱_2)G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)}`$ (B24) $`=\delta (𝐱𝐱_1)\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_m^n\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}Z(t_1);`$ (B25) for $`t_1t_2`$, $`{\displaystyle \mathrm{d}^dx_2\kappa ^{m\beta _2}(t^{}t_2,𝐱𝐱_2)G_{\beta _1\beta _2,m}^{\alpha _1}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)}`$ (B26) $`=\delta (𝐱𝐱_1)\left[\kappa _{,\beta _1n}^{\alpha _1\beta _2}(t^{}t_2,0)+\kappa _{,mn}^{m\beta _2}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}\right]\widehat{\mathrm{\Lambda }}_{\beta _2}^nZ(t_2),`$ (B27) $`{\displaystyle \mathrm{d}^dx_2\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,𝐱𝐱_2)G_{\beta _1\beta _2}^{\alpha _1\alpha _2}(t^{},𝐱\|t_1,𝐱_1;t_2,𝐱_2)}`$ (B28) $`=\delta (𝐱𝐱_1)\kappa _{,n\alpha _2}^{m\beta _2}(t^{}t_2,0)\widehat{\mathrm{\Lambda }}_m^n\widehat{\mathrm{\Lambda }}_{\beta _1}^{\alpha _1}\widehat{\mathrm{\Lambda }}_{\beta _2}^{\alpha _2}Z(t_2).`$ (B29) With the aid of these expressions and Eq. (81), one obtains Eq. (88) of Sec. III B.
warning/0002/hep-th0002041.html	ar5iv	text	# References
warning/0002/hep-ph0002271.html	ar5iv	text	# Scaling and Duality in Semi-exclusive Processes ## I Introduction Scaling is a well established phenomenon in deep inelastic scattering (DIS). The cross section with specific kinematic factors removed gives structure functions that depend on only the scaling variable $`x_B`$, up to calculable logarithmic corrections. In addition, an inclusive-exclusive connection—“Bloom-Gilman scaling” —is observed in these totally inclusive (at least on the hadronic side) reactions. Duality in this situation means that resonance bumps observed in the structure functions at low momentum transfers $`Q^2`$ average out to the smooth structure function measured at higher momentum transfers but the same $`x_B`$. Usually, but not always, duality is realized in such a way that as the resonance peak moves in $`x_B`$ with changing $`Q^2`$, the ratio of the peak height to the height of the scaling curve evolved from higher $`Q^2`$ is constant. Both scaling and scaling violation have played a crucial role in understanding the constituents of elementary particles and in establishing QCD as as the accepted theory of the strong interactions. Duality is in detail less well understood . It seems, however, to show that the fundamental single quark QCD process is still decisive in setting the scale of the reaction in the resonance region, and that the crucial role of the final state interactions in forming the resonance becomes moot when averaged over, say, the resonance width. This last observation, if reliably understood, could allow one to use duality to study the structure functions in the interesting and still experimentally uncertain $`x_B1`$ region. For a fixed available energy, $`x_B1`$ means getting into the resonance region and if one were sure of the connection of the resonance region average to the scaling curve, one could determine the scaling curve significantly closer to the kinematic upper endpoint. Departing from DIS, we want to continue test our ability to understand and apply QCD to describe hadronic processes. A set of processes that can be a new testing ground for both scaling and duality phenomena are semi-exclusive reactions typified by $$\gamma +p\pi +X,$$ (1) where the photon may be real or virtual. These processes are the topic of this paper. We shall study suitable kinematic variables for the general case and, when we get more detailed, give special attention to photoproduction with large photon to pion momentum transfer, $`t`$. A first requirement is to find a scaling region. This problem has been studied in the high $`Q^2`$–low $`(t/Q^2)`$ limit, focusing on the totally exclusive reaction but with extension to the semi-exclusive case . These authors found that scaling functions would exist, provided the photon and pion currents directly and successively interacted with the same quark while the rest acted as spectators. We here, concentrating on photoproduction at high $`\|t\|`$, show that perturbative QCD (pQCD) predicts there is indeed a scaling region. We shall below show the kinematic factors that connect the cross section to the expected scaling function. We shall also see that the scaling region does require kinematics where photopion production is dominated by direct interactions of both the photon and the pion , such as seen in Fig. 1. In particular, one must avoid regions where the pion comes from soft processes or comes as part of a jet from a fragmenting parton. In earlier work we were able to show that regions of direct pion production exist, and therefore there are regions where we can find a scaling function. When scaling is established at high $`\|t\|`$, one can study duality. One can ask whether the scaling curve from high $`Q^2`$ or $`t`$ a decent average over the resonance bumps seen at the same $`x`$ but lower $`Q^2`$ or $`t`$? Duality in this sense appears to be true for all the resonances seen in DIS. Further, one can ask if the bump to continuum ratio is constant as $`Q^2`$ or $`t`$ changes? This constancy is seen in DIS for most resonances, but not for the $`\mathrm{\Delta }`$(1232). While studying the kinematics and working in photoproduction context, we will see that it is possible in a single experiment with good kinematic coverage to probe a given $`x`$ region over a wide range of $`m_X`$ from the resonance region to well into the continuum region. The paper proceeds as follows. Section II will discuss the kinematics and scaling variable for the semiexclusive process. Section III will show how a scaling function emerges for semi-exclusive hard pion photoproduction, and also show the existence of a region where direct pion production dominates, specifically for a situation of 30 GeV incoming photons. Section IV will show that pQCD expectations for the resonance peak/scaling curve ratio at changing $`\|t\|`$ are similar to what one sees in DIS. Section V will offer some conclusions. ## II kinematic variables For the process $`\gamma +p\pi +X`$, define the Mandelstam variables by $$s=(p+q)^2,t=(qk)^2,u=(pk)^2.$$ (2) Define $`x`$ in general by, $$x=\frac{t}{s+u2m_N^2q^2m_\pi ^2},$$ (3) and note that all quantities defining $`x`$ are experimentally measurable . One can show $`0x1`$, and $`x=1`$ corresponds to the case that $`X`$ is a nucleon. Also generally, the hadronic mass recoiling against the pion is given by $$m_X^2=m_N^2t\left(\frac{1}{x}1\right).$$ (4) Specializing to the case where direct pion production, Fig. 1, is the underlying process, in the limit of high $`t`$ or $`s`$ and high recoil mass $`m_X`$, one can show that this $`x`$ is the fraction of the target’s momentum carried by the struck quark. The proof involves defining Mandelstam variables for the subprocess $`\gamma +q\pi +q^{}`$. We anticipate the result by letting the momentum of the struck quark be called $`xp`$, and get $`\widehat{s}`$ $`=`$ $`(xp+q)^2=x(sq^2)+q^2,`$ (5) $`\widehat{t}`$ $`=`$ $`(qk)^2=t,`$ (6) $`\widehat{u}`$ $`=`$ $`(xpk)^2=xu,`$ (7) where we have neglected masses but not $`q^2`$. For the direct pion production subprocess, $$\widehat{s}+\widehat{t}+\widehat{u}=q^2,$$ (8) and substituting Eqn. (5) leads to the identification of $`x`$ the momentum fraction to $`x`$ the experimental observable, when individual particle masses can be neglected. Thus $`x`$ is a precise analog of the observable $`x_B=Q^2/2m_N\nu `$ in deep inelastic scattering. Our formulas should and do connect to the well known ones for deep inelastic kinematics in the limit $`k0`$. In this limit, $`um_N^2`$ and $`tq^2Q^2`$ and $`m_X^2`$ $`=`$ $`m_N^2+Q^2\left({\displaystyle \frac{1}{x}}1\right),`$ (9) $`x`$ $`=`$ $`{\displaystyle \frac{t}{sm_N^2q^2}}={\displaystyle \frac{Q^2}{2pq}},`$ (10) without approximation. Still regarding deep inelastic scattering, Bloom and Gilman found that near threshold scaling worked better if one used a revised variable defined as $`1/x_B^{}=1/x_B+m_N^2/Q^2`$. By analogy to Bloom and Gilman’s proposal we could define a modified scaling variable with $`t`$ replacing $`Q^2`$: $$\frac{1}{x^{}}=\frac{1}{x}+\frac{m_N^2}{t},$$ (11) whence $$m_X^2=t\frac{1x^{}}{x^{}}.$$ (12) One should keep this possibility in mind here also. Another situation, related to the one we are pursuing, is semi-exclusive deep inelastic scattering with parallel kinematics. This means high $`Q^2`$ and an observed meson with three-momentum parallel to the incoming photon, in the lab. In this case, there is a variable $`z`$ defined by $$z\frac{pk}{pq}.$$ (13) and obtain $`m_X^2`$ $`=`$ $`m_N^2+Q^2(1z)\left({\displaystyle \frac{1}{x}}1\right),`$ (14) $`x`$ $`=`$ $`{\displaystyle \frac{Q^2}{2pq}},`$ (15) with the neglect of terms of $`𝒪(m_N^2x^2/Q^2)`$. ## III scaling and kinematic regions Now we shall focus on hard pion photoproduction, where $`q^2=0`$, $`k_T`$ is large, and $`\|t\|`$ is large. We will be mainly interested in direct pion production with $`m_X`$ large, and in the transition to the exclusive reactions $`\gamma +N\pi X`$. Other processes do contribute. In particular there are soft processes, and processes where the pion is produced as part of the fragmentation of a quark or gluon into a jet. These processes can be evaded if one can go to sufficient transverse momentum. We will comment on them briefly before proceeding. Soft processes are frequently approximated using vector meson dominance of the photon interaction, illustrated in Fig. 2. They are important at low transverse momenta, although the boundary between “low” and “high” is higher than one might expect, namely around 2 GeV. We have considered these processes in a fashion suitable for the present context in ; one can also find a representation of them in PYTHIA . Moderate transverse momenta hard pions can be produced by a fragmentation of a parton. The process is perturbatively calculable and could be a way to learn about polarized of unpolarized gluon distributions of the target ; one example is illustrated in Fig. 3. Neither the fragmentation nor the soft process is useful for the present duality study. The reason is that the experimental $`x`$ variable for them does not have a unique connection to the quark momentum fraction, and we will not be able to prove a scaling relation for them. Direct pion production, however, does have the nice connection between an experimentally observable $`x`$ and the struck quark momentum fraction, and it is calculable in pQCD. It is a higher twist process. Factors of the decay constant enter the amplitude, representing the quark-antiquark wave function of the pion at the origin, and must be dimensionally compensated by an extra power of $`s`$ in the cross section. Nonetheless, it can dominate over fragmentation at high $`k_T`$ because it always gives all the transverse momentum in the pion direction to the one pion. For the direct process, we can operationally define a scaling function $`F(x,t)`$ by $`E_\pi {\displaystyle \frac{d\sigma }{d^3k}}`$ $`=`$ $`{\displaystyle \frac{(sm_N^2)x^2}{\pi t}}{\displaystyle \frac{d\sigma }{dxdt}}`$ (16) $`=`$ $`{\displaystyle \frac{(sm_N^2)x^2}{\pi t}}{\displaystyle \frac{d\widehat{\sigma }}{dt}}(\gamma q\pi q^{})F(x,t).`$ (17) I.e., the scaling function is related to the cross section by some kinematic factors, which are partly explicit above and partly given in terms of the cross section for the subprocess $`{\displaystyle \frac{d\widehat{\sigma }(\gamma qMq^{})}{dt}}`$ $`=`$ $`{\displaystyle \frac{128g_F^2\pi ^2\alpha \alpha _s^2}{27(t)\widehat{s}^2}}I_M^2\left({\displaystyle \frac{e_q}{\widehat{s}}}+{\displaystyle \frac{e_q^{}}{\widehat{u}}}\right)^2`$ (18) $`\times `$ $`\left[\widehat{s}^2+\widehat{u}^2+\lambda h\left(\widehat{s}^2\widehat{u}^2\right)\right],`$ (19) where we should substitute quark charges relevant for pion being produced, for example $`e_q=e_u`$ and $`e_q^{}=e_d`$ for the $`\pi ^+`$. The flavor factor $`g_F`$ is unity for the $`\pi ^\pm `$ and $`1/\sqrt{2}`$ for the $`\pi ^0`$. $`I_M`$ is for the present purpose a constant factor, but if the perturbative calculation is valid, it will be given in terms of the distribution amplitude of the meson as $`𝑑\xi _1\varphi _M(\xi ,\mu ^2)/\xi _1`$. For the asymptotic distribution amplitude, $`I_\pi =\sqrt{3}f_\pi /2`$ with $`f_\pi 93`$ MeV. We have included polarization dependence for future use: $`\lambda `$ is the helicity of the photon and $`h`$ is twice the helicity of the target quark. Of course, duality can be tested with polarization as well as without. The reason to believe that the above expression, Eqn. (16), produces a scaling formula is that the perturbative formula valid for short distance pion production (the process of Fig. 1) is, $`E_\pi {\displaystyle \frac{d\sigma }{d^3k}}`$ $`=`$ $`{\displaystyle \frac{(sm_N^2)x^2}{\pi t}}`$ (20) $`\times `$ $`{\displaystyle \underset{q}{}}{\displaystyle \frac{d\widehat{\sigma }}{dt}}(\gamma q\pi q^{})G_{q/T}(x,\mu ^2)`$ (21) Thus where perturbation theory works, there is a scaling function $`F(x,t)`$ is mainly dependent on $`x`$. We can relate it to the quark distributions (with weak dependence on the scale $`\mu ^2`$, which we may set to $`t`$), as in DIS. We expect the formulas will be mainly applied in the high $`x`$ region, where valence quarks dominate. Hence the comment on the choice of the quark charges just above. Let us comment on the fact that the presence of a hard gluon exchange (see Fig. 1) indicates that one needs sufficiently high energies to apply the pQCD formalism. However, since only one pion distribution amplitude is involved for the direct process, if the photon attaches to the produced $`q\overline{q}`$ pair of Fig. 1 (the worse case), the average virtuality of the gluon in question corresponds to the one determining the pion electromagnetic form factor at $`Q^220(35)`$ GeV<sup>2</sup> scale, for the asymptotic (Chernyak-Zhitnitsky) pion distribution amplitude assuming a CEBAF energy of 12 GeV, pion emission angle of 22, and $`m_X=2`$ GeV (see Ref. for details). Therefore one may hope to observe a single-gluon exchange, which is a higher twist effect, in inclusive photoproduction of pions even at CEBAF energies generally considered not high enough to reach the perturbative QCD domain. Indications of direct pion production off a quark were obtained in $`\pi N`$ scattering (see Ref. for references and discussion). We may now ask if this scaling function dual, in a Bloom-Gilman like sense, to the bumpier curve one will get in the resonance region? The resonance region, of course, is what we have at the very highest transverse momentum, where there is very little energy left over to put into recoil mass. Formally, the duality relation then may be written as an integral of the differential cross section $`d\sigma /dxdt(\gamma N\pi X)`$ and in the region of the direct process dominance reads $`{\displaystyle _{(1\frac{m_X^2m_N^2}{t})^1}^1}`$ $`dx`$ $`{\displaystyle \underset{q}{}}G_{q/N}(x){\displaystyle \frac{d\sigma }{dt}}(\gamma q\pi q^{})`$ (22) $``$ $`{\displaystyle \underset{R}{}}`$ $`{\displaystyle \frac{d\sigma }{dt}}(\gamma +N\pi +R).`$ (23) Summation in the right hand side of Eqn. (22) is done over all resonances $`R`$ with masses $`m_Rm_X`$, with the nucleon final state included. If the parton distribution function of the nucleon is $`G_{q/N}(1x)^3`$ at $`x1`$ and the subprocess $`\gamma q\pi q^{}`$ cross section is determined by the one-gluon exchange mechanism of Fig. 1, then as will be shown in the next section duality as in Eqn. (22) requires that the resonance excitation cross section $`d\sigma /dt(\gamma +N\pi +R)1/s^7`$ at fixed $`t/s`$—the result known from the constituent counting rules . The duality relation above could also be written using the modified scaling variable $`x^{}`$ from Eqn. (11). One should ask if the proper regions exist. There needs to be a region where direct pion production dominates, where one can measure the scaling curve and see how it tails off into the resonance region. Such a region does exist. ¿From earlier studies we have the machinery to calculate the direct pion and fragmentation process, and estimate the VMD processes, and have shown that the direct process, even though it is higher twist, does take over at some point if we have enough initial photon energy. A useful presentation of our calculated results is shown in Fig. 4. The figure attempts to show that we can follow a given $`x`$ region from the resonance region until well into the scaling region, and do so in a single experiment. The axes of Fig. 4 give the outgoing pion transverse and longitudinal momenta, in the target rest frame. Some labeled straight lines give the pion angle relative to the incoming beam. The three solid elliptical curves each correspond to a fixed value of recoil mass $`m_X`$. The outermost curve has $`m_X=m_N`$ and thus corresponds to the quasielastic process $`\gamma N\pi N`$, and also marks the kinematic limit of pion momenta. The next curve has $`m_X=2`$ GeV, and the innermost solid curve has $`m_X=3.5`$ GeV. Thus the region between the two outermost curves is the resonance region, and the region within the middle solid curve is the continuum region. The segment above the grey band is the region where direct pion production dominates. For us, this is the “good region.” (As a side note, the grey band is straighter than we might have guessed, especially since it is made up of two parts. The central part comes from the fragmentation process growing larger. Both ends come from VDM, as modeled in an earlier note, which we think was conservative in estimating the size of the VDM contributions.) Finally comes the important dotted elliptical curve, which has a constant $`x`$, specifically $`x=0.7`$ in this case. We see that we can thinkably measure the putative scaling function in the resonance region at small pion angles, and then by moving to larger angle, follow its behavior at the same $`x`$ but larger $`m_X`$ (and larger $`\|t\|`$) well out of the resonance region, before running into a region where fragmentation or soft processes dominate. (As another aside, lines of constant $`\|t\|0`$ on this plot would be parabolas opening to the right, and passing through the small line segment between the origin and the lower (negative) limit of $`k_L`$; $`\|t\|=0`$ occurs along the positive $`k_L`$ axis.) ## IV resonance bumps vs. the scaling curve There is always a resonance region. In plots of $`F(x,t)`$ vs. $`x`$, the bumpy resonance region slides to the right with increasing $`\|t\|`$. In the corresponding DIS case the bumps slide neatly down the curve, with the resonance/smooth curve ratio observed to stay the same, for most resonances. Within pQCD, this is expected theoretically as a consequence of the known behaviors of the scaling curve as $`x1`$ and the predicted falloff of the resonance transition form factors at high $`Q^2`$. We can show that the resonance/continuum constancy is consistent with pQCD in the semi-exclusive case also. We need to find the behavior of $`F_{res}(x,t)={\displaystyle \frac{d\sigma }{dxdt}}(\gamma N\pi R)/{\displaystyle \frac{d\sigma }{dt}}(\gamma q\pi q^{})`$ (24) at (say) the resonance peak for large $`\|t\|`$ (and $`x1`$). The denominator in this limit is $$\frac{d\sigma }{dt}=g(t/s)\|t\|^3(1x)^3$$ (25) where $`g(t/s)`$ is a known function (see Eqn. (18)) which does not go to zero for $`t/s`$ finite, and we have used $`1/t(1x)`$ in the stated limit. The numerator for a finite width resonance can be approximated by (for $`x1`$), $`\left({\displaystyle \frac{d\sigma }{dxdt}}\right)_{res}`$ $``$ $`{\displaystyle \frac{\|t\|}{2m_X}}\left({\displaystyle \frac{d\sigma }{dm_Xdt}}\right)_{res}`$ (26) $``$ $`{\displaystyle \frac{\|t\|}{2m_R}}\left({\displaystyle \frac{d\sigma }{dt}}\right)_{res}{\displaystyle \frac{\mathrm{\Gamma }/2\pi }{(m_Xm_R)^2+\mathrm{\Gamma }^2/4}}`$ (27) where $`\mathrm{\Gamma }`$ is the width of the resonance and we have used a simple lorentzian form to give the resonance shape. The pQCD scaling rules tell us that $$\left(\frac{d\sigma }{dt}\right)_{res}=f(t/s)\|t\|^7$$ (28) where $`f(t/s)`$ is not known but in general it should not go to zero for finite $`t/s`$. Thus, $`\left({\displaystyle \frac{d\sigma }{dxdt}}\right)_{respeak}{\displaystyle \frac{1}{\pi m_R\mathrm{\Gamma }}}f(t/s)\|t\|^6(1x)^6`$ (29) Thus, $$F_{respeak}(x,t)(1x)^3.$$ (30) This is how the height of a resonance peak fall with $`x`$ as $`x1`$. It is also precisely the pQCD expectation for the scaling curve. Hence the resonance/continuum ratio is in general constant, at least at high $`\|t\|`$, as it is for DIS. In DIS, the Delta(1232) is an exception, as it falls markedly with $`Q^2`$ ; $`Q^2`$ in lepton scattering is the analog of $`t`$ in hard meson photoproduction. It will be interesting to see if the Delta(1232) disappears with increasing $`\|t\|`$ and if the, say, S<sub>11</sub>(1535) stays up at high $`\|t\|`$. Recall that in pQCD, the disappearing Delta in electron scattering is explained as an accident having to do with the specifics of the Delta and nucleon wave functions . We should not expect this to be necessarily replicated in pion photoproduction since the integrals over the distribution amplitudes will involve different weightings. ## V Conclusions and Discussion Semi-exclusive processes give an opportunity to extend the studies of scaling and duality, which in deep inelastic scattering have been fruitful in verifying our understanding of QCD and in pushing our effort to deepen that understanding. It appears that scaling in the sense that the cross section is directly related to a scaling function that depends, up to logarithmic corrections, on just one variable. The scaling variable for semi-exclusive processes, given in the text, is related to the momentum fraction of the struck quark, just like the scaling variable in deep inelastic scattering. However, scaling, at least as we have been able to present it in this paper, works in semi-exclusive process only when the pion is produced directly off the same quark that absorbs the incoming photon. We have been able to show, theoretically, that such a scaling region does exist. One should bear in mind that there are soft kinematic regions where one does not know where the pion comes from, and fragmentation regions where the pion is produced at some remove from the fundamental process that initiates the reaction. We do not know of a scaling function for these regions, and it is not trivial that one can avoid them, but one can. One should also bear in mind that a certain amount of initial energy is needed to be able to produce a scaling region. For incoming photon lab energy 16 GeV or below and our present estimates of the vector meson dominance contributions, it does not appear that there is a region where VMD is not the biggest process for photoproduction, at least if one does not make any additional cuts. However, a there are possibilities for reducing the necessary incoming energy. One follows from noting that a directly produced pion is also a pion produced in kinematic isolation, not as part of a jet, and one can consider an “isolation cut,” a requirement that there be no other particles collinear with the pion. Another possibility is to have the photon off shell, since then the vector meson propagator is significantly reduced, reducing the VMD contributions without there being an equal reduction for other contributions. We are hopeful that using electroproduction and isolation cuts can make the incoming energy requirement low enough to fit an upgraded CEBAF range, but are deferring detailed elaboration. The existence of a scaling region also allows one to consider the inclusive-exclusive connection with the resonance region. Will the resonance bumps average out to the smooth scaling curve measured at higher $`t`$ and evolved to lower $`t`$ ? Will the resonance peak to scaling curve ratio be independent of $`t`$? In deep inelastic physics, it does appear that the final state interactions which produce the resonance are irrelevant to the overall rate of resonance region production, if one does a suitable average. And we have shown that for the semi-exclusive case, as in the deep inelastic case, the resonance to continuum ratio should be constant, barring special circumstances. A special circumstance in the deep inelastic case occurs for the $`\mathrm{\Delta }(1232)`$, which disappears into the scaling curve with increasing $`Q^2`$. One would like to know if similar phenomena occur in other situations. Testing scaling and duality in inclusive photoproduction of mesons requires coverage of the large-$`x`$ region, where the cross sections are rapidly falling as $`x`$ approaches it upper limit. Therefore such a uniquely designed high-luminosity machine as CEBAF, with a bit more energy, could do an excellent job in these duality studies. We thank Nathan Isgur, Wally Melnitchouk, Chris Armstrong, Rolf Ent, and Cynthia Keppel for useful discussions about duality and give the latter three a second thanks for showing us their experimental duality results from CEBAF. AA thanks the US Department of Energy for support under contract DE-AC05-84ER40150; CEC and CW thank the NSF for support under grant PHY-9900657.
warning/0002/astro-ph0002048.html	ar5iv	text	# ” ## 1 We cannot measure magnetic fields of the majority of stars, which are distant far from us. Therefore, the existence of magnetic fields for the majority of stars can be considered only hypothetically. However, the magnetic field of the Sun is known over than a hundred years, and in the last decades the astronomers managed to measure magnetic fields for a number of stars (so-called $`A_p`$-stars) and some pulsars . It is interesting to construct a model describing the generation of magnetic fields by stars and to compare it with the data of the astronomers. The mechanism examined below is based on the gravity-induced electric polarization of matter. It is capable to explain also the generation of magnetic fields by planets , however, in the case of stars, this mechanism works in the purest manner. ## 2 The action of gravity on metals has often been a topic of discussion before -. The basic result of these researches is reduced to the statement that gravity induces inside a metal an electric field with an intensity $$\stackrel{}{E}\frac{m_i\stackrel{}{g}}{e},$$ (1) where $`m_i`$ is the mass of an ion, $`\stackrel{}{g}`$ is gravity acceleration, $`e`$ is the electron charge. This field is so small that it is not possible to measure it experimentally. It is a direct consequence of the presence of an ion lattice in a metal. This lattice is deformed by gravity and then the electron gas adapts its density to this deformation. The resulting field becomes very small. Under superhigh pressure, all substances transform into ultradense matter usually named nuclear-electron plasma . It occurs when external pressure enhances the density of matter several times . Such values of pressure exist inside celestial bodies. In nuclear-electron plasma the electrons form the degenerated Fermi gas. At the same time, the positively charged ions form inside plasma a dense packing lattice ,. As usually accepted, this lattice may be replaced by a lattice of spherical cells of the same volume. The radius $`r_s`$ of such a spherical cell in plasma of the mass density $`\gamma `$ is given by $$\frac{4\pi }{3}r_s^3=\left(\frac{\gamma }{m_i}\right)^1=\frac{Z}{n},$$ (2) where Z is the charge of the nucleus, $`m_i=Am_p`$ is the mass of the nucleus, A is the atomic number of the nucleus, $`m_p`$ is the mass of a proton, and n is the electron number density $$n=\frac{3Z}{4\pi r_s^3}.$$ (3) The equilibrium condition in matter is described by the constancy of its electrochemical potential . In plasma, the direct interaction between nuclei is absent, therefore the equilibrium in a nuclear subsystem of plasma (at $`T=0`$) looks like $$\mu _i=m_i\psi +Ze\phi =const.$$ (4) Here $`\phi `$ is the potential of an electric field and $`\psi `$ is the potential of a gravitational field. The direct action of gravitation on electrons can be neglected. Therefore, the equilibrium condition in the electron gas is $$\mu _e=\frac{p_F^2}{2m_e}(e\delta q)\phi =const,$$ (5) where $`m_e`$ is the mass of an electron and $`p_F`$ is the Fermi momentum. By introducing the charge $`\delta q`$, we take into account that the charge of the electron cloud inside a cell can differ from $`e`$. A small number of electrons can stay on the surface of a plasma body where the electric potential is absent. It results that the charge in a cell, subjected to the action of the electric potential, is effectively decreased on a small value $`\delta q`$. The electric polarization in plasma is a result of changing in density of both nuclear and electron gas subsystems. The electrostatic potential of the arising field is determined by the Gauss’ law $$^2\phi =\frac{1}{r^2}\frac{d}{dr}\left[r^2\frac{d}{dr}\phi \right]=4\pi \left[Ze\delta (r)en\right],$$ (6) where the position of nuclei is described by the function $`\delta (r)`$. According to the Thomas - Fermi method, $`n`$ is approximated by $$n=\frac{8\pi }{3h^3}p_F^3.$$ (7) With this substitution, Eq.(6) is converted into a nonlinear differential equation for $`\phi `$, which for $`r>0`$ is given by $$\frac{1}{r^2}\frac{d}{dr}\left(r^2\frac{d}{dr}\phi (r)\right)=4\pi \left[\frac{8\pi }{3h^3}\right]\left[2m_e(\mu _e+(e\delta q)\phi )\right]^{3/2}.$$ (8) It can be simplified by introducing the following variables : $$\mu _e+(e\delta q)\phi =Ze^2\frac{u}{r}$$ (9) and $`r=ax`$, where $`a=\{\frac{9\pi ^2}{128Z}\}^{1/3}a_0`$ with $`a_0=\frac{\mathrm{}^2}{m_ee^2}=`$ Bohr radius. With the account of Eq.(4) $$Ze^2\frac{u}{r}=const\frac{m_i\psi }{Z}\delta q\phi .$$ (10) Then Eq.(8) gives $$\frac{d^2u}{dx^2}=\frac{u^{3/2}}{x^{1/2}}.$$ (11) In terms of u and x, the electron density within a cell is given by $$n_{TF}=\frac{8\pi }{3h^3}p_F^3=\frac{32Z^2}{9\pi ^3a_0^3}\left(\frac{u}{x}\right)^{3/2}.$$ (12) Under the influence of gravity the charge of the electron gas in a cell becomes equal to $$Q_e=4\pi e_0^{r_s}n(r)r^2𝑑r=\frac{8\pi e}{3h^3}\left[2m_e\frac{Ze^2}{a}\right]^{3/2}4\pi a^3_0^{x_s}x^2𝑑x\left[\frac{u}{x}\right]^{3/2}.$$ (13) Using Eq.(11), we obtain $$Q_e=Ze_0^{x_s}x𝑑x\frac{d^2u}{dx^2}=Ze_0^{x_s}𝑑x\frac{d}{dx}\left[x\frac{du}{dx}u\right]=Ze\left[x_s\frac{du}{dx}\|_{x_s}u(x_s)+u(0)\right].$$ (14) At $`r0`$ the electric potential is due to the nucleus alone $`\phi (r)\frac{Ze}{r}`$. It means that $`u(0)1`$ and each cell of plasma obtains a small charge $$\delta q=Ze\left[x_s\frac{du}{dx}\|_{x_s}u(x_s)\right]=Zex_{s}^{}{}_{}{}^{2}\left[\frac{d}{dx}\left(\frac{u}{x}\right)\right]_{x_s}.$$ (15) For a cell placed in a point $`R`$ inside a star $$\delta q=Zer_s^2\left[\frac{d}{dR}\left(\frac{u}{r}\right)\right]\left[\frac{dR}{dr_s}\right].$$ (16) Considering that gravity acceleration $`\stackrel{}{g}=\frac{d\psi }{dR}`$ and the electric field intensity $`\stackrel{}{E}=\frac{d\phi }{dR}`$ $$\frac{dr_s}{dR}=\frac{r_s^2}{e}\left[\frac{\frac{m_i}{Z}\stackrel{}{g}+\delta q\stackrel{}{E}}{\delta q}\right].$$ (17) This equation has the following solution $$\frac{dr_s}{dR}=0$$ (18) and $$\frac{m_i}{Z}\stackrel{}{g}+\delta q\stackrel{}{E}=0.$$ (19) In plasma, the equilibrium value of the electric field on nuclei according to Eq.(4) is determined by Eq.(1) as well as in a metal. But there is one more additional effect in plasma. Simultaneously with the supporting of nuclei in equilibrium, each cell obtains an extremely small positive electric charge. As $`div\stackrel{}{g}=4\pi Gnm_i`$ and $`div\stackrel{}{E}=4\pi n\delta q`$, the gravity-induced electric charge in a cell $$\delta q=\sqrt{G}\frac{m_i}{Z}10^{18}e,$$ (20) where $`G`$ is the gravity constant. However, because the sizes of bodies may be very large, the electric field intensity may be very large as well $$\stackrel{}{E}=\frac{\stackrel{}{g}}{\sqrt{G}}.$$ (21) In accordance with Eqs.(18,19), the action of gravity on matter is compensated by the electric force and the gradient of pressure is absent. Thus, a celestial body is electrically neutral as a whole, because the positive volume charge is concentrated inside the charged core and the negative electric charge exists on its surface and so one can infer gravity-induced electric polarization of a body. ## 3 At the surface of the core, the electric field intensity reduces to zero. The jump in electric field intensity is accompanied at the surface of the core by the pressure jump $`\mathrm{\Delta }p(R_N)`$. It leads to the redistribution of the matter density inside a star. In a celestial body consisting of matter with an atomic structure, density and pressure grow monotonously with depth. In a celestial body consisting of electron-nuclear plasma, the pressure gradient inside the polarized core is absent and the matter density is constant. Pressure affecting the matter inside this body is equal to the pressure jump on the surface of the core $$p=\mathrm{\Delta }p(R_N)=\frac{E(R_N)^2}{8\pi }=\frac{2\pi }{9}G\gamma ^2R_N^2,$$ (22) where R<sub>N</sub> is the radius of the core. One can say that this pressure jump is due to the existence of the polarization jump or, which is the same, the existence of the bound surface charge formed by an electron pushed out from the core and making the total charge of the celestial body equal to zero. Because the electron subsystem of plasma inside a star is the relativistic Fermi gas, we can write its equation of state $$p=\frac{(3\pi ^2)^{1/3}}{4}\frac{\mathrm{}c\gamma ^{4/3}}{m_{p}^{}{}_{}{}^{4/3}\beta ^{4/3}}$$ (23) where $`\beta m_p`$ is the mass of the matter related to one electron of the Fermi gas system, and $`m_p`$ is the proton mass. Because of the electroneutrality, one proton should be related to electron of the Fermi gas of plasma. The existence of one neutron per proton is characteristic for a substance consisting of light nuclei. The quantity of neutrons grows approximately to 1.8 per proton for the heavy nuclei substance. Therefore, it is necessary to expect that inside stars $$2<\beta <2.8.$$ (24) As pressure inside a star is known (Eq.(22)), from Eq.(23) it is possible to determine a steady-state value of mass of a star $$M_{}=\zeta A_{}^{3/2}\frac{m_p}{\beta ^2}.$$ (25) This mass is expressed by dimensionless constants only $$A_{}=\left(\frac{\mathrm{}c}{Gm_{p}^{}{}_{}{}^{2}}\right)=1.5410^{38}$$ (26) $`\zeta =(1.5^5\pi )^{1/2}5`$, and the slowly varying parameter $`\beta `$ (Eq.(24)). The masses of stars can be measured with a considerable accuracy, if these stars compose a binary system. There are almost 200 double stars whose masses are known with the required accuracy . Among these stars there are giants, white dwarfs, and stars of the main sequence. Their averaged mass is described by the equality $$M_{}=\left(1.36\pm 0.05\right)M_{},$$ (27) where $`M_{}`$ is the mass of the Sun. The center of this distribution (Fig.1) corresponds to Eq.(25) at $`\beta 2.6`$. It is interesting to note that the ”biography” of such a star appears much poorer than in the Chandrasecar model. Temperature does not influence the parameters of relativistic plasma. Therefore, a star with a mass close to the steady-state value (Eq.(25)) is in a stable equilibrium not depending on temperature. It should not collapse with a temperature decreasing. The instability of a star can arise with burning out of light nuclei - deuterium and helium - and with a related increasing of $`\beta `$. This growth leads to the reduction of a steady-state value of mass (Eq.(25)) and, probably, to the distraction of stars with greater masses. ## 4 As the density of matter inside a relativistic star is constant, it is possible to assume that it equals the mean density of the Sun and to estimate a star radius $$R\left(\frac{M_{}}{\frac{4\pi }{3}\gamma _{}}\right)^{1/3},$$ (28) where $`\gamma _{}`$ is the mean density of the Sun. It allows one to calculate the momentum of a star as the momentum of a sphere with a constant density $$I=\frac{2}{5}M_{}R^2$$ (29) and at a known frequency of rotation $`\mathrm{\Omega }`$ to calculate its angular momentum $$L=\frac{2}{5}M_{}\mathrm{\Omega }R^2.$$ (30) In this model the magnetic moment of a star is created by the rotation of a star as a whole. Thus, it is composed of two parts. One is the magnetic moment of the layer of electrons placed on the external surface of a star $$\mu _{}=\frac{1}{3}\left(\frac{4\pi }{3}\rho R^3\right)\mathrm{\Omega }R^2.$$ (31) The second component of the magnetic moment is created by the positively charged core $$\mu _+=\frac{1}{5}\left(\frac{4\pi }{3}\rho R^3\right)\mathrm{\Omega }R^2.$$ (32) The summary moment is $$\mu _\mathrm{\Sigma }=\frac{2}{15}\left(\frac{4\pi }{3}\rho R^3\right)\mathrm{\Omega }R^2.$$ (33) It is remarkable that the gyromagnetic relation of a star, i.e., the relation of its magnetic moment to the angular momentum, is expressed through world constants only $$\vartheta =\frac{\mu _\mathrm{\Sigma }}{L}=\frac{\sqrt{G}}{3c}.$$ (34) The measurements permit us to define the frequency of rotation and magnetic fields for a number of stars . It appears enough to check up the considered theory, since masses of stars and their momenta are determined inside the theory (Eq.(27) and (Eq.(30))). The magnetic moments as functions of their angular momenta for all celestial objects (for which they are known today) are shown in Fig.2. The data for planets are taken from , the data for stars are taken from , and for pulsars - from . As it can be seen from this figure with the logarithmic accuracy, all celestial bodies - stars, planets, and pulsars - really have the gyromagnetic ratio close to the universal value (Eq.(34)). Only the data for the Moon fall out, because its size is too small to create an electrically polarized core. ## 5 Apparently, the considered theory is quite true for pulsars which consist, as it is supposed, from the neutron substance with an addition of electrons and protons . As this substance is a relativistic one, there is a fair definition of a steady-state value of mass Eq.(25). The astronomers measured masses of 16 radio-pulsars and 7 x-ray pulsars included in a double system . According to this data, the distribution of masses of pulsars is $$M_{pulsar}=(1.38\pm 0.03)M.$$ (35) The center of this distribution corresponds to Eq.(25) at $`\beta 2.6`$. The gyromagnetic relations are measured for three pulsars only . These values are in a quite satisfactory agreement with Eq.(34)(Fig.2). For the majority of pulsars , there are estimations of magnetic fields obtained using a number of model assumptions . It is impossible to consider these data as the data of measurements, but nevertheless they also are in some agreement with Eq.(34), (Fig.3). For planets the situation is more difficult. First, inside planets the substance forms not relativistic electron-nuclear plasma, but nonrelativistic electron-ion plasma. It has different equation of state leading to a more complex expression for the stable mass of a planet core than the expression of Eq.(25) for stars. Second, a noncharged layer at the surface of the core can take a significant part of a planet’s volume and it is impossible to neglect a role of this stratum. However, it can be seen from Fig.2 that the gyromagnetic relations of planets are also in the quite satisfactory agreement with Eg.(34). The detailed calculation for the Earth gives for the magnetic moment $`410^{25}Oecm^3`$, which is almost exactly twice smaller than the measured value $`8.0510^{25}Oecm^3`$. Thus, it is possible to assume that the basic component of the magnetic moment of planets is induced by the same mechanism which is working in stars.
warning/0002/hep-th0002222.html	ar5iv	text	# 1 Introduction ## 1 Introduction One of the most beautiful symmetries in string theory is the radius inversion symmetry $`R1/R`$ of a circle, known as T-duality . This is the symmetry which exchanges winding modes on a circle with momentum modes on the dual circle. This symmetry has been the underlying motivation for many of the subsequent dualities discovered in string theory and in quantum field theories. In particular, over a decade ago, it was conjectured in that a similar duality might exist in the context of string propagation on Calabi-Yau manifolds, where the role of the complex deformations on one manifold get exchanged with the Kahler deformations on the dual manifold. The pairs of manifolds satisfying this symmetry are known as mirror pairs, and this duality is also called mirror symmetry. There has been a lot of progress since the original formulation of this conjecture, in its support. In particular many examples of this phenomenon were found . The intermediate step in the derivation for this class of examples involved the construction of conformal field theory for certain Calabi-Yau’s and their identification with certain Landau-Ginzburg models . This connection was further elucidated in where it was shown that the linear sigma model is a powerful tool in the study of strings propagating on a Kahler manifold. It was shown in how mirror symmetry can be used very effectively to gain insight into non-perturbative effects involving worldsheet instantons. Roughly speaking this amounts to counting the number of holomorphic curves in a Calabi-Yau manifold. This made the subject also interesting for algebraic geometers in the context of enumerative geometry. Motivated by the existing examples some general class of mirror pairs were formulated by mathematicians using toric geometry . Moreover a program to prove the rational curve counting formula, predicted by mirror symmetry, from the view point of localization and virtual fundamental cycles was initiated in and was pushed to completion in . For reviews of various aspects of mirror symmetry see ; for mathematical aspects of mirror symmetry see the excellent book . The question of a proof of mirror symmetry and its relation with T-duality, which was its original motivation, was further pursued in where it was shown that for certain toroidal orbifold models mirror symmetry reduces to T-duality. More generally, by following the prediction of the map of D-branes under mirror symmetry, it was argued in that mirror symmetry should reduce to T-duality in a more general context. Furthermore, it was shown in , how the general suggestion for construction of mirror pairs proposed using toric geometry can be intuitively related to T-duality. Mirror symmetry has also been extended from the case of Calabi-Yau sigma models to more general cases. On the one hand there are proposals as to what the mirror theories are in the case of certain sigma models with positive first Chern class . On the other hand there are proposals for what the mirror of non-compact Calabi-Yau manifolds are . The aim of this paper is to present a proof of mirror symmetry for all cases proposed thus far. The proof depends crucially on establishing a dual description of $`(2,2)`$ supersymmetric gauge theories in 1+1 dimensions. The dual theory is found using the idea analogous to Polyakov’s model of confinement in quantum electrodymanics in $`2+1`$ dimensions . He considered a $`U(1)`$ gauge theory which includes magnetic monopoles playing the role of instantons. $`U(1)`$ Maxwell theory of gauge coupling constant $`e`$ in $`2+1`$ dimensions is dual to the theory of a periodic scalar field $`\sigma \sigma +2\pi `$ with the Lagrangian $`e^2\|\mathrm{d}\sigma \|^2`$. The gas of instantons and anti-instantons with a long range interaction between them generates a potential term $$U(\sigma )=\mu ^3\mathrm{cos}(\sigma )$$ (1.1) in the effective Lagrangian in terms of the dual variable $`\sigma `$, where $`\mu `$ is the mass scale determined by $`e`$ and the monopole size. One sees from this that a mass gap is generated and that an electric flux is confined into a thin tube. We note that the description in terms of the dual variable $`\sigma `$ was essential in this argument. In supersymmetric field theories, instanton computation can be used to obtain exact results for some important physical quantities. For some of the striking examples, see . Among these, and treat supersymmetric gauge theories in $`2+1`$ dimensions and the effective theory is described in terms of the dual variable as in . Also, in , duality between vector and vector in $`3+1`$ dimensions was used to solve the problem in an essential way. We apply an analogous idea to study the long distance behaviour of $`(2,2)`$ gauge theories, making use of instantons which are vortices in this case. We dualize the phase of the charged fields in the sense of $`R1/R`$ duality and describe the low energy effective theory in terms of the dual variables. We will see that a superpotential is dynamically generated by the instanton effect, as in , and we can exactly determine the (twisted) F-term part of the effective Lagrangian. To be specific, let us consider a $`(2,2)`$ supersymmetric $`U(1)`$ gauge theory with $`N`$ chiral multiplets of charge $`Q_i`$ ($`i=1,\mathrm{},N`$). In addition to the gauge coupling, the theory has two parameters: Fayet-Iliopoulos and Theta parameters. They are combined into a single complex parameter $`t`$ and appear in the twisted superpotential as $`t\mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ is the twisted chiral field which is the field strength of the gauge multiplet (and includes the scalar, the gaugino and the field strength). Each charged chiral field is sent by the duality on its phase to a twisted chiral field $`Y_i`$ which is a neutral periodic variable $`Y_iY_i+2\pi i`$ that couples to the field strength as a dynamical Theta angle $`Q_iY_i\mathrm{\Sigma }`$. The exact twisted superpotential we will find is given by $$\stackrel{~}{W}=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_iY_it\right)+\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i},$$ (1.2) where $`\mu `$ is a scale parameter. The term proportional to $`\mathrm{\Sigma }`$ is the one that appears already at the dualization process. The exponentials of $`Y_i`$’s are the ones that are generated by instanton effect. When $`_iQ_i0`$, $`\mu `$ is a scale required to renormalize the FI parameter $`t`$. In this case, a combination of $`t`$ and $`\mu `$ is a fake and only one dimensionful parameter $`\mathrm{\Lambda }=\mu \mathrm{e}^{t/_iQ_i}`$ is the real parameter of the theory. This is the standard dimensional transmutation. In the case where $`_iQ_i=0`$, $`t`$ is the dimensionless parameter of the theory and $`\mu `$ is a fake as it can be absorbed by a field redefinition. Using the connection between $`U(1)`$ gauge theories with matter and sigma models on Kahler manifolds we then relate the above result to the statement of mirror symmetry<sup>1</sup><sup>1</sup>1The idea to use the gauged linear sigma model to derive mirror symmetry was also considered in .. In particular we find that the mirror to a sigma model is a Landau-Ginzburg model. In the case of Calabi-Yau manifolds this can also be related to the sigma model on another Calabi-Yau manifold by the equivalence of sigma models and Landau-Ginzburg models. In the case of manifolds with non-zero first Chern class, however, this is not possible (we consider only manifolds with non-negative first Chern class since otherwise the sigma model would not be well-defined): the axial $`U(1)`$ R-symmetry is borken by an anomaly and therefore the vector $`U(1)`$ R-symmetry of the mirror theory must be broken by an inhomogenious superpotential. Likewise, since the vector $`U(1)`$ R-symmetry of the original non-linear sigma model is an exact symmetry, the mirror manifold (on which the Landau-Ginzburge superpotential is defined) must always be Calabi-Yau so that the axial R-symmetry is unbroken. A typical example of manifolds of positive first Chern class is $`𝐂\mathrm{P}^1`$. It has been observed that the supersymmetric $`𝐂\mathrm{P}^1`$ sigma model is mirror to the $`N=2`$ sine-Gordon theory which is a sigma model on a cylinder $`𝐂^\times `$ with a sine-Gordon superpotential . The two theories have $`U(1)\times 𝐙_4`$ vector-axial (or axial-vector) R-symmetries. Both have two massive vacua which spontaneously breaks $`𝐙_4`$ to $`𝐙_2`$. Moreover, soliton spectrum and the scattering matrix have been observed to agree . Actually, this mirror symmetry is the first non-trivial one that can be derived by our method. The linear sigma model for $`𝐂\mathrm{P}^1`$ is a $`U(1)`$ gauge theory with two chiral multiplets of charge $`Q_1=Q_2=1`$. By integrating out the $`\mathrm{\Sigma }`$ field from (1.2), we obtain the constraint $`Y_1+Y_2=t`$ which can be solved by $`Y_1=Y+t/2,Y_2=Y+t/2`$. Then, the superpotential (1.2) takes the form $$\stackrel{~}{W}(Y)=2\mathrm{\Lambda }\mathrm{cosh}(Y),$$ (1.3) which is nothing but the sine-Gordon potential! In fact, the mirror symmetry of other models, including the conformal field theories based on Calabi-Yau sigma models, can be derived in a uniform way as a natural generalization of this example. Furthermore, the mirror theory provides an effective way to classify the vacua of the theory and to identify where they flow to in the infra-red limit. For example, for a degree $`d`$ hypersurface on $`𝐂\mathrm{P}^{N1}`$ of size $`t`$, we find the mirror to be the orbifold of the Landau-Ginzburg model with the superpotential $$W=X_1^d+\mathrm{}+X_N^d+e^{t/d}X_1\mathrm{}X_N$$ (1.4) by the group $`(𝐙_d)^{N1}`$ which acts on the indivisual fields $`X_i`$ by multiplication by $`d`$-th roots of unity, in such a way that $`X_1\mathrm{}X_N`$ is invariant. Note that the special case $`d=N`$ gives one the proposed mirror of sigma models on Calabi-Yau hypersurfaces (Greene-Plesser construction ). In the case with $`d<N`$, the hypersurface has a positive first Chern class and the sigma model is asymptotic free . The mirror theory given above shows that there are $`Nd`$ massive vacua at non-zero $`X_i`$’s and another vacuum at $`X_i=0`$. For $`d>2`$, the vacuum at $`X_i=0`$ flows to a non-trivial fixed point described by the LG orbifold with the same group $`(𝐙_d)^{N1}`$ and the superpotential (1.4) with the last term being dropped as it is irrelevant at low energies. In this example, the original linear sigma model actually leads to the description of the low energy theory in terms of another LG orbifold , with the same superpotential but with a different group — a single $`𝐙_d`$ acting on $`X_i`$’s uniformly. In fact, our result reproduces the mirror symmetry of LG orbifolds . In general, however, as the $`𝐂\mathrm{P}^1`$ example shows, our method provide information which is hardly available in the original model. The organization of this paper is as follows: In section 2 we review certain aspects of mirror symmetry with emphasis on its interplay with supersymmetry. We present the $`N=2`$ supersymmetric version of T-duality which plays a crucial role for us later in the paper. In section 3 we consider dynamical aspects of $`N=2`$ gauge theories and establish their equivalence with the above mentioned LG theories. In section 4 we review aspects of linear sigma model (i.e. gauge theory/sigma model connection). In section 5 we use the results in section 3 and present a proof of mirror symmetry. Also in this section we elaborate on what we mean by “proving” mirror symmetry. In section 6 we discuss some aspects of D-branes in the context of LG theories. This elucidates the relation between mirror symmetry for local (non-compact) and global (i.e. compact) sigma models which is discussed in section 7. Also the mirror of complete intersections in toric varieties are discussed there. In section 8 we discuss some possible directions for future work. In appendix A we present a conjectured generalization of our results to the case of complete intersections in Grassmannians and flag varieties. In appendix B some aspects of supersymmetry transformations needed in this paper are summarized. ## 2 Mirror Symmetry Of $`N=2`$ Theories In Two Dimensions In this section, we review some basic facts and fix notations on $`(2,2)`$ supersymmetric field theories in $`1+1`$ dimensions. We also define the notion of mirror symmetry and present some examples. In particular, we describe in detail the standard $`R1/R`$ duality and show how it can be viewed as mirror symmetry in the case of complex torus. ### 2.1 Supersymmetry ### $`(2,2)`$ Supersymmetry Algebra $`(2,2)`$ supersymmetry algebra is generated by four supercharges $`Q_\pm ,\overline{Q}_\pm `$, space-time translations $`P`$, $`H`$ and rotation $`M`$, and the generators $`F_V`$ and $`F_A`$ of two R-symmetries $`U(1)_V`$ and $`U(1)_A`$. These obey the following (anti-)commutation relations: $`Q_+^2=Q_{}^2=\overline{Q}_+^2=\overline{Q}_{}^2=0,`$ (2.1) $`\{Q_\pm ,\overline{Q}_\pm \}=2(HP),`$ (2.2) $`\{\overline{Q}_+,\overline{Q}_{}\}=2Z,\{Q_+,Q_{}\}=2Z^{},`$ (2.3) $`\{Q_{},\overline{Q}_+\}=2\stackrel{~}{Z},\{Q_+,\overline{Q}_{}\}=2\stackrel{~}{Z}^{},`$ (2.4) $`[M,Q_\pm ]=Q_\pm ,[M,\overline{Q}_\pm ]=\overline{Q}_\pm ,`$ (2.5) $`[F_V,Q_\pm ]=Q_\pm ,[F_V,\overline{Q}_\pm ]=\overline{Q}_\pm ,`$ (2.6) $`[F_A,Q_\pm ]=Q_\pm ,[F_A,\overline{Q}_\pm ]=\pm \overline{Q}_\pm .`$ (2.7) The hermiticity of the generators is dictated by $$Q_\pm ^{}=\overline{Q}_\pm .$$ (2.8) In the above expressions, $`Z`$ and $`\stackrel{~}{Z}`$ are central charges. The algebra with $`Z=\stackrel{~}{Z}=0`$ can be obtained by dimensional reduction of the four-dimensional $`N=1`$ supersymmetry algebra. $`U(1)_V`$ comes from the R-symmetry in four dimensions and $`U(1)_A`$ is the rotation along the reduced directions. It is not always the case that the two $`U(1)`$ R-symmetries are the symmetry of a given theory, except in superconformal field theories where they both must be the symmetry. In the class of theories we will consider in this paper, at least one of them is a symmetry and the other may or may not be broken to a discrete subgroup. A central charge can be non-zero if there is a soliton that interpolates different vacua and/or if the theory has a continuous abelian symmetry. As the name suggests, it must commute with the R-symmetry generators as well. In particular, $`Z`$ ($`\stackrel{~}{Z}`$) must always be zero in a theory where $`U(1)_V`$ ($`U(1)_A`$) is unbroken. In superconformal field theory, both must be vanishing. ### Superfields and Supersymmetric Lagrangians Fields in a supermultiplet can be combined into a single function, a superfield, of the superspace coordinates $`x^0,x^1,\theta ^\pm ,\overline{\theta }^\pm `$. Supersymmetry generators $`Q_\pm `$, $`\overline{Q}_\pm `$ act on the superfields as the derivatives $$Q_\pm =\frac{}{\theta ^\pm }+i\overline{\theta }^\pm \left(\frac{}{x^0}\pm \frac{}{x^1}\right),\overline{Q}_\pm =\frac{}{\overline{\theta }^\pm }i\theta ^\pm \left(\frac{}{x^0}\pm \frac{}{x^1}\right).$$ (2.9) These commute with another set of derivatives $$D_\pm =\frac{}{\theta ^\pm }i\overline{\theta }^\pm \left(\frac{}{x^0}\pm \frac{}{x^1}\right),\overline{D}_\pm =\frac{}{\overline{\theta }^\pm }+i\theta ^\pm \left(\frac{}{x^0}\pm \frac{}{x^1}\right).$$ (2.10) R-symmetries act on a superfield $`(x,\theta ^\pm ,\overline{\theta }^\pm )`$ as $`\mathrm{e}^{i\alpha F_V}(x,\theta ^\pm ,\overline{\theta }^\pm )=\mathrm{e}^{iq_V\alpha }(x,\mathrm{e}^{i\alpha }\theta ^\pm ,\mathrm{e}^{i\alpha }\overline{\theta }^\pm ),`$ (2.11) $`\mathrm{e}^{i\alpha F_A}(x,\theta ^\pm ,\overline{\theta }^\pm )=\mathrm{e}^{iq_A\alpha }(x,\mathrm{e}^{i\alpha }\theta ^\pm ,\mathrm{e}^{\pm i\alpha }\overline{\theta }^\pm ),`$ (2.12) where $`q_V`$ and $`q_A`$ are the vector and axial R-charges of $``$. The basic representations of the supersymmetry algebra are chiral and twisted chiral multiplets which both consist of a complex scalar and a Dirac fermion. These are represented by chiral (or $`cc`$) superfield and twisted chiral (or $`ac`$) superfield respectively. A chiral superfield $`\mathrm{\Phi }`$ satisfies $$\overline{D}_\pm \mathrm{\Phi }=0,$$ (2.13) and can be expanded as $$\mathrm{\Phi }=\varphi +\sqrt{2}\theta ^+\psi _++\sqrt{2}\theta ^{}\psi _{}+2\theta ^+\theta ^{}F+\mathrm{},$$ (2.14) where $`F`$ is a complex auxiliary field and $`+\mathrm{}`$ involves only the derivatives of $`\varphi ,\psi _\pm `$. The hermitian conjugate of $`\mathrm{\Phi }`$ is an anti-chiral (or $`aa`$) superfield $`D_\pm \overline{\mathrm{\Phi }}=0`$. A twisted chiral superfield $`Y`$ satisfies $$\overline{D}_+Y=D_{}Y=0,$$ (2.15) and can be expanded as $$Y=y+\sqrt{2}\theta ^+\overline{\chi }_++\sqrt{2}\overline{\theta }^{}\chi _{}+2\theta ^+\overline{\theta }^{}G+\mathrm{}.$$ (2.16) where $`G`$ is a complex auxiliary field and $`+\mathrm{}`$ involves only the derivatives of the component fields. The hermitian conjugate of $`Y`$ is a twisted anti-chiral (or $`ca`$) superfield $`D_+\overline{Y}=\overline{D}_{}\overline{Y}=0`$. We also introduce a vector multiplet. It consists of a vector field $`v_\mu `$, Dirac fermions $`\lambda _\pm `$, $`\overline{\lambda }_\pm `$ which are conjugate to each other, and a complex scalar $`\sigma `$ in the adjoint representation of the gauge group. It is represented in a vector superfield $`V`$ which is expanded (in the Wess-Zumino gauge) as $`V`$ $`=`$ $`\theta ^{}\overline{\theta }^{}(v_0v_1)+\theta ^+\overline{\theta }^+(v_0+v_1)\theta ^{}\overline{\theta }^+\sigma \theta ^+\overline{\theta }^{}\overline{\sigma }`$ $`+\sqrt{2}i\theta ^{}\theta ^+(\overline{\theta }^{}\overline{\lambda }_{}+\overline{\theta }^+\overline{\lambda }_+)+\sqrt{2}i\overline{\theta }^+\overline{\theta }^{}(\theta ^{}\lambda _{}+\theta ^+\lambda _+)+2\theta ^{}\theta ^+\overline{\theta }^+\overline{\theta }^{}D`$ where $`D`$ is a real auxiliary field. Using the gauge covariant derivatives $`𝒟_\pm =\mathrm{e}^VD_\pm \mathrm{e}^V`$, $`\overline{𝒟}_\pm =\mathrm{e}^V\overline{D}_\pm \mathrm{e}^V`$, we can define the field strength as $`\mathrm{\Sigma }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{\overline{𝒟}_+,𝒟_{}\}`$ $`=`$ $`\sigma +i\sqrt{2}\theta ^+\overline{\lambda }_+i\sqrt{2}\overline{\theta }^{}\lambda _{}+2\theta ^+\overline{\theta }^{}(DiF_{01})+\mathrm{},`$ where $`F_{01}`$ is the curvature of $`v_\mu `$. This is a twisted chiral (covariant) superfield $`\overline{𝒟}_+\mathrm{\Sigma }=𝒟_{}\mathrm{\Sigma }=0`$. The supersymmetry transformation of the component fields in this gauge (based on ) is recorded for convenience in Appendix B. Supersymmetric Lagrangian can be obtained from integrations over suitable fermionic coordinates. There are D-term, F-term, and twisted F-term. D-term is for arbitrary superfields $`_i`$ and is given by $$\mathrm{d}^4\theta K(,\overline{})=\frac{1}{4}d\theta ^+d\theta ^{}d\overline{\theta }^{}d\overline{\theta }^+K(,\overline{}),$$ (2.19) where $`K(,\overline{})`$ is an arbitrary real function of $`_i`$’s. This is classically R-invariant under any assignment of R-charges. F-term is for chiral superfields $`\mathrm{\Phi }_i`$ and is given by $$\mathrm{d}^2\theta W(\mathrm{\Phi })+c.c.=\frac{1}{2}\mathrm{d}\theta ^{}\mathrm{d}\theta ^+W(\mathrm{\Phi })\|_{\overline{\theta }^\pm =0}+\frac{1}{2}\mathrm{d}\overline{\theta }^+\mathrm{d}\overline{\theta }^{}\overline{W}(\overline{\mathrm{\Phi }})\|_{\theta ^\pm =0}$$ (2.20) where $`W(\mathrm{\Phi })`$ is a holomorphic function of $`\mathrm{\Phi }_i`$’s and is called a superpotential. This is invariant under vector and axial R-symmetries only when it is possible to assign R-charges to $`\mathrm{\Phi }_i`$’s so that $`W(\mathrm{\Phi })`$ has vector and axial charge 2 and 0 respectively. Twisted F-term is for twisted chiral superfields $`Y_i`$ and is given by $$\mathrm{d}^2\stackrel{~}{\theta }\stackrel{~}{W}(Y)+c.c.=\frac{1}{2}\mathrm{d}\overline{\theta }^{}\mathrm{d}\theta ^+\stackrel{~}{W}(Y)\|_{\overline{\theta }^+=\theta ^{}=0}+\frac{1}{2}\mathrm{d}\overline{\theta }^+\mathrm{d}\theta ^{}\overline{\stackrel{~}{W}}(\overline{Y})\|_{\theta ^+=\overline{\theta }^{}=0}$$ (2.21) where $`\stackrel{~}{W}(Y)`$ is a holomorphic function of $`Y_i`$’s and is called a twisted superpotential. For R-invariance, it is required that R-charges can be assigned to $`Y_i`$’s so that $`\stackrel{~}{W}(Y_i)`$ has vector and axial charge 0 and 2 respectively. ### Non-linear Sigma Model and Potentials Supersymmetric non-linear sigma model on a Kahler manifold $`X`$ is one of the main subjects of the present paper. It is described by a set of chiral superfields $`\mathrm{\Phi }^i`$ $`(i=1,\mathrm{},n)`$ representing complex coordinates of $`X`$. Let $`g_{i\overline{ȷ}}=^2K/\varphi ^i\overline{\varphi }^{\overline{ȷ}}`$ be the Kahler metric of $`X`$ where $`K(\varphi ,\overline{\varphi })`$ is a Kahler potential. The Lagrangian can be given by the D-term $`\mathrm{d}^4\theta K(\mathrm{\Phi },\overline{\mathrm{\Phi }})`$ which is expressed in terms of the component fields as $`L_K`$ $`=`$ $`g_{i\overline{ȷ}}^\mu \varphi ^i_\mu \overline{\varphi }^{\overline{ȷ}}+ig_{i\overline{ȷ}}\overline{\psi }_{}^{\overline{ȷ}}(D_0+D_1)\psi _{}^i+ig_{i\overline{ȷ}}\overline{\psi }_+^{\overline{ȷ}}(D_0D_1)\psi _+^i`$ (2.22) $`+R_{i\overline{k}j\overline{l}}\psi _+^i\psi _{}^j\overline{\psi }_{}^{\overline{k}}\overline{\psi }_+^{\overline{l}},`$ after the auxiliary fields are eliminated. Here, $`R_{i\overline{k}j\overline{l}}`$ is the curvature tensor with respect to the Levi-Civita connection $`\mathrm{\Gamma }_{lj}^i=g^{i\overline{k}}_lg_{j\overline{k}}`$ of $`X`$. Also, $$D_\mu \psi _\pm ^i=_\mu \psi _\pm ^i+_\mu \varphi ^l\mathrm{\Gamma }_{lj}^i\psi _\pm ^j$$ (2.23) is the covariant derivative with respect to the connection induced on the worldsheet. If $`X`$ has a holomorphic function $`W(\varphi )`$, the non-linear sigma model can be deformed by the F-term with superpotential $`W(\mathrm{\Phi })`$, $$L=\mathrm{d}^4\theta K(\mathrm{\Phi },\overline{\mathrm{\Phi }})+\frac{1}{2}(\mathrm{d}^2\theta W(\mathrm{\Phi })+c.c.).$$ (2.24) Eliminating the auxiliary fields, we obtain the Lagrangian $`L_K+L_W`$ where $`L_K`$ is given in (2.22) and the deformation term is $$L_W=\frac{1}{4}g^{\overline{ȷ}i}_{\overline{ȷ}}\overline{W}_iW\frac{1}{2}(D_i_jW)\psi _+^i\psi _{}^j\frac{1}{2}(D_{\overline{ı}}_{\overline{ȷ}}\overline{W})\overline{\psi }_{}^{\overline{ı}}\overline{\psi }_+^{\overline{ȷ}},$$ (2.25) in which $`D_i_jW=_i_jW\mathrm{\Gamma }_{ij}^l_lW`$. Note that a non-trivial holomorphic function $`W(\varphi )`$ exists only when $`X`$ is non-compact. In some cases, the deformation discretize the energy spectrum which would be continuous without $`L_W`$ because of the non-compactness. The character of the theory therefore depends largely on the asymptotic behaviour of the superpotential. If $`X`$ has a holomorphic isometry generated by a holomorphic vector field $`V`$, the sigma model can be deformed by another kind of potential term . This is obtained by first gauging the isometry as in , taking the weak coupling limit, and freezing the vector multiplet fields at $`\sigma =\stackrel{~}{m}`$, $`v_\mu =0`$, and $`\lambda _\pm =\overline{\lambda }_\pm =0`$. A description in $`(2,2)`$ superspace was considered in . The deformation term is $$L_V=g_{i\overline{ȷ}}\|\stackrel{~}{m}\|^2V^i\overline{V}^{\overline{ȷ}}\frac{i}{2}\left(g_{i\overline{ı}}_jV^ig_{j\overline{ȷ}}_{\overline{ı}}\overline{V}^{\overline{ȷ}}\right)\left(\stackrel{~}{m}\overline{\psi }_{}^{\overline{ı}}\psi _+^j+\overline{\stackrel{~}{m}}\overline{\psi }_+^{\overline{ı}}\psi _{}^j\right).$$ (2.26) The deformed Lagrangian is invariant under the modified $`(2,2)`$ supersymmetry where the modification is given by $`\mathrm{\Delta }Q_{}\overline{\psi }_+^{\overline{ı}}=\sqrt{2}i\stackrel{~}{m}\overline{V}^{\overline{ı}}`$, $`\mathrm{\Delta }\overline{Q}_+\psi _{}^i=\sqrt{2}i\stackrel{~}{m}V^i`$, and their complex conjugates (the action on bosonic fields is not modified). The central charge of the modified supersymmetry is non-vanishing and is given by $$\stackrel{~}{Z}=i\stackrel{~}{m}_V$$ (2.27) where $`_V`$ acts on the fields as $`_V\varphi ^i=V^i`$ and $`_V\psi _\pm ^i=_jV^i\psi _\pm ^j`$. If the superpotential $`W(\varphi )`$ is invariant under the isometry $`V^i_iW=0`$, then, the sigma model can be deformed by $`L_W`$ and $`L_V`$ at the same time without breaking $`(2,2)`$ supersymmetry. One can also consider the deformation by a set of commuting holomorphic isometries $`V_1,\mathrm{},V_n`$; simply replace $`\stackrel{~}{m}V^i_{a=1}^n\stackrel{~}{m}_aV_a^i`$, $`\overline{\stackrel{~}{m}}V^i_{a=1}^n\overline{\stackrel{~}{m}}_aV_a^i`$ (the bosonic term $`\|\stackrel{~}{m}\|^2\|V\|^2`$ in (2.26) is replaced by $`\frac{1}{2}\|_{a=1}^n\stackrel{~}{m}_aV_a\|^2\frac{1}{2}\|_{a=1}^n\overline{\stackrel{~}{m}}_aV_a\|^2`$). ### R-Symmetry The R-symmetries $`U(1)_V`$ and $`U(1)_A`$ are not always symmetries of the theory (except their $`𝐙_2`$ subgroups which acts as the sign flip of spinors — $`2\pi `$-rotation). It can be broken at the classical level by potential terms or at the quantum level by anomaly. However, superconformal field theories always possess both symmetries. We illustrate this in the class of theories introduced above. We first consider the non-linear sigma model on $`X`$ without potentials. Both $`U(1)`$’s are classically unbroken. $`U(1)_V`$ remains a symmetry of the quantum theory. $`U(1)_A`$ is subject to the chiral anomaly which is proportional to the trace of the curvature of the connection (2.23) of the tangent bundle of $`X`$. Thus, $`U(1)_A`$ is anomalous if and only if the first Chern class of $`X`$ is non-vanishing; $`c_1(X)0`$. In particular, both $`U(1)`$’s are unbroken in the sigma model on a CY manifold $`X`$, which is expected to flow to a superconformal field theory of central charge $`c/3=dimX`$. If $`\varphi ^{}c_1(X)`$ is always an integer multiple of some $`p`$, $`U(1)_A`$ is broken to its discrete subgroup $`𝐙_{2p}`$ which can be further broken spontaneously. If the theory flows to a non-trivial fixed point, $`U(1)_A`$ must be recovered there. We next consider a theory with a superpotential $`W(\varphi )`$. $`U(1)_A`$ is classically unbroken but is subject to the chiral anamaly. $`U(1)_V`$ is unbroken if and only if $`W(\varphi )`$ is scale invariant in the sense that it is possible to assign the vector R-charges to $`\mathrm{\Phi }^i`$ so that $`W(\mathrm{\Phi })`$ has vector charge $`2`$. The theory has a mass gap if at all the critical points of $`W`$ the Hessian is non-degenerate, $`det_i_jW0`$. At a degenerate critical point, the theory can flow to a non-trivial fixed point where $`U(1)_V`$ is recovered. An example where both $`U(1)`$’s are unbroken is the LG model on $`𝐂^N`$ with a quasi-homogeneous superpotential; $`W(\lambda ^{2q_i}\varphi _i)=\lambda ^2W(\varphi _i)`$ where vector R-charge $`2q_i`$ is assigned to $`\varphi _i`$. It is believed that such a model flows in the IR to an $`N=2`$ superconformal field theory with central charge $`c/3=_i(12q_i)`$ . Finally, if the sigma model is perturbed by a holomorphic isometry, $`U(1)_A`$ is explicitly broken by the fermion mass term in (2.26). This is consistent with the non-vanishing of the susy central charge (2.27). This is also related to the fact that the scalar component of the vector multiplet has a canonical axial charge $`2`$. ### Supersymmetric Ground States Let us examine the ground states of the theory. We compactify the spacial direction on $`S^1`$ and put a periodic boundary condition on all fields. We also assume $`Z=\stackrel{~}{Z}=0`$. As in any supersymmetric field theory, a state annihilated by all of $`Q_\pm ,\overline{Q}_\pm `$ is a zero energy ground state, and vice versa. Let $`Q`$ be one of $`Q_{}+\overline{Q}_+`$ and $`\overline{Q}_++\overline{Q}_{}`$ or their hermitian conjugates. Then, it follows from the algebra (2.1)-(2.4) that $$\{Q,Q^{}\}=2H.$$ (2.28) It also follows from (2.1)-(2.4) that $$Q^2=0,$$ (2.29) and we can consider the cohomology of states using $`Q`$ as the coboundary operator. In the theory where $`H`$ has a discrete spectrum, (2.28) and (2.29) imply that the supersymmetric ground states are in one to one correspondence with the $`Q`$-cohomology classes. The index of the operator $`Q`$ is the Witten index $`\mathrm{Tr}(1)^F`$ which is invariant under perturbation of the theory. If the central charge in the supersymmetry algebra is non-vanishing because of a continuous abelian global symmetry group $`T`$, say, as $`\stackrel{~}{Z}=_{a=1}^n\lambda _aS_a`$ where $`S_a`$ are the generators of $`T`$, (2.28) still holds but (2.29) is modified as $`Q^2=2_{a=1}^n\lambda _aS_a`$. However, when restricted to $`T`$-invariant states, $`Q`$ is still nilpotent and we can consider $`Q`$-cohomology. This is the $`T`$-equivariant cohomology. Since a continuous symmetry cannot be broken in $`1+1`$ dimensions, the supersymmetric ground states are still in one to one correspondence with the $`T`$-equivariant $`Q`$-cohomology classes. For a sigma model on a compact Kahler manifold $`X`$, $`Q=Q_{}+\overline{Q}_+`$ reduces in the zero momentum sector to the exterior derivative $`\mathrm{d}=+\overline{}`$ acting on differential forms on $`X`$. In fact, the zero momentum approximation is exact as far as vacuum counting is concerned and the supersymmetric vacua and harmonic forms are in one to one correspondence. In particular $`\mathrm{Tr}(1)^F=\chi (X)`$. The R-charge of a vacuum corresponding to a harmonic $`(p,q)`$-form is $`q_V=p+q`$ and $`q_A=p+qdim_𝐂X`$, where the shift by $`dim_𝐂X`$ is for the invariance of the spectrum under conjugation $`q_Aq_A`$. Of course, if $`c_1(X)`$ is non-zero and $`U(1)_A`$ is anomalous, the axial charge $`q_A`$ does not make sense. However, if $`𝐙_{2p}U(1)_A`$ is non anomalous, there is a $`𝐙_{2p}`$ grading in the space of vacua. If $`X`$ is non-compact and has a superpotential $`W(\mathrm{\Phi })`$, the spectrum is discrete at sufficiently low energies if $`g^{\overline{ȷ}i}_{\overline{ȷ}}\overline{W}_iWc>0`$ at all infinity in the field space. When $`W(\mathrm{\Phi })`$ has only non-degenerate critical points, the supersymmetric vacua are in one to one correspondence with the critical points. When $`W(\mathrm{\Phi })`$ can be perturbed to such a situation, the index is the number of critical points. If the compact sigma model is deformed by a holomorphic isometry $`V`$, the susy central charge is non-vanishing and the nilpotency of $`Q=Q_{}+\overline{Q}_+`$ is modified as $`Q^2=2i\stackrel{~}{m}_V`$. Since this is a small perturbation, the index remains the same as the $`\stackrel{~}{m}=0`$ case, $`\mathrm{Tr}(1)^F=\chi (X)`$. In the zero momentum sector, $`Q`$ is proportional to $`\mathrm{d}_{\stackrel{~}{m}}=\mathrm{d}\sqrt{2}i\stackrel{~}{m}i_V`$. If $`V`$ has only non-degenerate zeroes, the supersymmetric vacua are in one to one correspondence with the zeroes of $`V`$. In particular, the index is the number of zeroes (the Hopf index theorem). See for more details. ### Chiral Ring We can also consider cohomology of local operators with respect to $`Q=\overline{Q}_++\overline{Q}_{}`$ or $`Q=Q_{}+\overline{Q}_+`$ (when the central charge is zero). We shall call a local operator commuting with $`Q_{cc}=\overline{Q}_++\overline{Q}_{}`$ (resp. $`Q_{ac}=\overline{Q}_++Q_{}`$) a chiral or $`cc`$ operator (resp. a twisted chiral or $`ac`$ operator). The lowest component of a (twisted) chiral superfield is a (twisted) chiral operator. It follows from the supersymmetry algebra that the space-time translation of a (twisted) chiral operator is $`Q`$-exact and does not change the cohomology class of the operator. A product of two (twisted) chiral operators is annihilated by $`Q`$. They commute with each other up to $`Q`$-exact operators since one can make them space-like separated. Therefore, the $`Q`$-cohomology group of local operators form a commutative ring . This is called chiral or $`cc`$ ring for $`Q_{cc}`$-cohomology and twisted chiral or $`ac`$ ring for $`Q_{ac}`$-cohomology. In general $`cc`$ and $`ac`$ rings in a given theory are different from each other. If the susy central charge is non-zero because of an abelian global symmetry, we can define equivariant chiral ring in an obviuos way. ### Twisting to Topological Field Theory It is often useful to twist $`N=2`$ theories to topological field theories . This is possible when the quantum theory possesses at least either one of $`U(1)_V`$ or $`U(1)_A`$ R-symmetries (which we call here $`U(1)_R`$, generator $`R`$) under which the R-charges are all integral. It is standard to call it A-twist for $`R=F_V`$ and B-twist for $`R=F_A`$. We start with the Euclidean version of the theory (obtained by Wick rotation $`x^0=ix^2`$ from the Minkowski theory). It has the supersymmetry with the same algebra (2.1)-(2.7) and the same hermiticity condition (2.8). Twisting is to replace the group $`U(1)_E`$ of space-time rotation generated by $`M`$ by the diagonal subgroup of $`U(1)_E\times U(1)_R`$, considering $`M^{}=M+R`$ as the new generator of the rotation group. In particular, the twisted theory on a curved worldsheet is obtained by gauging the diagonal subgroup of $`U(1)_E\times U(1)_R`$ (instead of $`U(1)_E`$) by the spin connection. The energy momentum tensor on the flat worldsheet is thus modified as $`T_{\mu \nu }^{\mathrm{twisted}}=T_{\mu \nu }+(1/4)(ϵ_\mu ^\lambda _\lambda J_\nu ^R+ϵ_\nu ^\lambda _\lambda J_\mu ^R)`$ where $`J_\mu ^R`$ is the $`U(1)_R`$ current, and the spin of fields and conserved currents are modified. An important aspect of the twisted theory is that some of the supercharges have spin zero and make sense without reference to the coordinates. These are $`Q_{}`$ and $`\overline{Q}_+`$ for A-twist while $`\overline{Q}_\pm `$ for B-twist (see (2.5)-(2.7)). As we have seen, $`Q_{ac}=Q_{}+\overline{Q}_+`$ and $`Q_{cc}=\overline{Q}_++\overline{Q}_{}`$ are nilpotent if the susy central charges are zero. Then, we can consider $`Q=Q_{ac}`$ or $`Q_{cc}`$ as a BRST operator that selects operators in the A-twisted or B-twisted model respectively. In particular, $`ac`$ ring elements are the physical operators in the A-model and $`cc`$ ring elements are the physical operators in the B-model. $`T_{\mu \nu }^{\mathrm{twisted}}`$ is $`Q`$-exact in the class of theories we consider in this paper, and the correlation functions of $`Q`$-invariant operators are independent of the choice of the worldsheet metric. In this sense the twisted theory is a topological field theory. Even if the susy central charge is non-zero because of an abelian global symmetry, the twisted theory can still be considered as topological field theory, physical operators selected by equivariant cohomology. For the non-linear sigma model possibly with a superpotential, A-twist is possible when $`W(\mathrm{\Phi })`$ is scale invariant while B-twist is possible when $`c_1(X)=0`$. A-twist of a model with $`W0`$ yields topological sigma model while B-model on a flat manifold $`X`$ is called topological LG . The fermion path-integral of a B-model is a chiral determinant and is made well-defined using the anomaly cancellation condition $`c_1(X)=0`$. The definition involves the choice of a multipicative factor which can be translated to the choice of a nowhere vanishing holomorphic $`n`$-form where $`n`$ is the complex dimension of $`X`$. Indeed, $`c_1(X)=0`$ assures the existence of such an $`n`$-form. ### Spectral Flow Let us return to the untwisted $`N=2`$ theories (where we assume $`Z=\stackrel{~}{Z}=0`$ for now). We have seen that supersymmetric ground states of the theory on a periodic circle and the (twisted) chiral ring are both characterized as the cohomology with respect to the nilpotent supercharge $`Q`$. There is in fact an intimate relation between them: In a theory which can be A-twisted (resp. B-twisted), there is a one-to-one correspondence between supersymmetric ground states and $`ac`$ ring elements (resp. $`cc`$ ring elements). In a theory where both A and B-twist are possible, the space of supersymmetric ground states, $`ac`$ ring and $`cc`$ ring are all the same as a vector space (but of course $`ac`$ and $`cc`$ rings are different as a ring in general). This can be seen conveniently using the twisted theory. We consider a theory which is A-twistable. Let us insert an $`ac`$ ring element $`𝒪`$ at the tip of a long cigar-like hemi-sphere in the A-twisted theory. The twisted theory is equivalent to the untwisted $`N=2`$ theory on the flat cylinder region, and we obtain a state of the untwisted theory at the boundary circle. Because of the twisting in the curved region, the fermions are periodic on the boundary circle. Now, the supersymmetric ground state corresponding to $`𝒪`$ is the state at the boundary circle in the limit where the cylinder region becomes infinitely long. The state is indeed a ground state because the infinitely long cylinder plays the role of projection to zero energy states. Various aspects of this relation and the geometry of vacuum states and its relation to the chiral rings have been studied in . In particular this geometry is captured by what is called the $`tt^{}`$ equations. ### 2.2 The Mirror Symmetry The $`(2,2)`$ supersymmetry algebra (2.1)-(2.7) is invariant under the outer automorphism given by the exchange of the generators $$Q_{}\overline{Q}_{},F_VF_A,Z\stackrel{~}{Z}.$$ (2.30) Mirror symmetry is an equivalence of two $`(2,2)`$ supersymmetric field theories under which the generators of supersymmetry algebra are exchanged according to (2.30). A chiral multiplet of one theory is mapped to a twisted chiral multiplet of the mirror, and vice versa. It is of course a matter of convention which to call $`Q_{}`$ or $`\overline{Q}_{}`$. Here, we are assuming the standard convention where, in the sigma model on a Kahler manifold (possibly with a superpotential), the complex coordinates are the lowest components of chiral superfields. If we flip the convention of one of a mirror pair, the two theories are equivalent without the exchange of (2.30). #### 2.2.1 Mirror Symmetry between Tori In the case of supersymmetric sigma model on a flat torus, it has been known that mirror symmetry reduces to the $`R1/R`$ duality performed on a middle dimensional torus. Below, we review this in the simplest case of sigma model into the algebraic torus $`𝐂^\times =𝐑\times S^1`$. We start with recalling the bosonic $`R1/R`$ duality. ### $`R1/R`$ Duality Let us consider the following action for a periodic scalar field $`\phi `$ of period $`2\pi `$ $$S_\phi =\frac{1}{4\pi }_WR^2h^{\mu \nu }_\mu \phi _\nu \phi \sqrt{h}\mathrm{d}^2x,$$ (2.31) where $`h_{\mu \nu }`$ is the metric on the worldsheet $`W`$ which we choose to be of Euclidean signature. This is the action for a sigma model into a circle $`S^1`$ of radius $`R`$. This action can be obtained also from the following action for $`\phi `$ and a one-form field $`B_\mu `$ $$S^{}=\frac{1}{2\pi }_W\frac{1}{2R^2}h^{\mu \nu }B_\mu B_\nu \sqrt{h}\mathrm{d}^2x+\frac{i}{2\pi }_WB\mathrm{d}\phi .$$ (2.32) Completing the square with respect to $`B_\mu `$ which is solved by $$B=iR^2\mathrm{d}\phi ,$$ (2.33) and integrating it out, we obtain the action (2.31) for the sigma model. If, changing the order of integration, we first integrate over the periodic scalar $`\phi `$, we obtain a constraint $`\mathrm{d}B=0`$. If the worldsheet $`W`$ is a genus $`g`$ surface, there is a $`2g`$-dimensional space of closed one-forms modulo exact forms<sup>1</sup><sup>1</sup>1This can be easily extended to the case of worldsheets with boundaries.. One can choose a basis $`\omega ^i`$ ($`i=1,\mathrm{},2g`$) such that each element has integral periods on one-cycles on $`W`$ and that $`_W\omega ^i\omega ^j=J^{ij}`$ is a non-degenerate matrix of integral entry. Then, a general solution to $`\mathrm{d}B=0`$ is $$B=\mathrm{d}\vartheta _0+\underset{i=1}{\overset{2g}{}}a_i\omega ^i,$$ (2.34) where $`\vartheta _0`$ is a real scalar field and $`a_i`$’s are real numbers. Integration over $`\phi `$ actually yields constraints on $`a_j`$’s as well. Recall that $`\phi `$ is a periodic variable of period $`2\pi `$. This means that $`\phi `$ does not have to come back to its original value when circling along a nontrivial one-cycles in $`W`$, but comes back to itself up to $`2\pi `$ shifts. For such a topologically nontrivial configuration, $`\mathrm{d}\phi `$ has an expansion like (2.34) with non-zero coefficient $`a_i`$ for $`\omega ^i`$ which is dual to the one-cycle. That the shift is only allowed to take integer multiples of $`2\pi `$ means that such $`a_i`$ is constrained to be $`2\pi n_i`$ where $`n_i`$ is an integer. Thus, for a general configuration of $`\phi `$ we have $$\mathrm{d}\phi =\mathrm{d}\phi _0+\underset{i=1}{\overset{2g}{}}2\pi n_i\omega ^i,$$ (2.35) where $`\phi _0`$ is a single valued function on $`W`$. Now, integration over $`\phi `$ means integration over the function $`\phi _0`$ and summation over the integers $`n_i`$’s. Integration over $`\phi _0`$ yields the constraint $`\mathrm{d}B=0`$ which is solved by (2.34). What about the summation over $`n_i`$’s? To see this we substitute in $`B\mathrm{d}\phi `$ for $`B`$ from (2.34); $$_WB\mathrm{d}\phi =2\pi \underset{i,j}{}a_iJ^{ij}n_j.$$ (2.36) Now, noting that $`J^{ij}`$ is a non-degenerate integral matrix and using the fact that $`_n\mathrm{e}^{ian}=2\pi _m\delta (a2\pi m)`$, we see that summation over $`n_i`$ constrains $`a_i`$’s to be an integer multiples of $`2\pi `$; $$a_i=2\pi m_i,m_i𝐙.$$ (2.37) Inserting this into (2.34), we see that $`B`$ can be written as $$B=\mathrm{d}\vartheta ,$$ (2.38) where now $`\vartheta `$ is a periodic variable of period $`2\pi `$. Now, inserting this to the original action we obtain $$S_\vartheta =\frac{1}{4\pi }_W\frac{1}{R^2}h^{\mu \nu }_\mu \vartheta _\nu \vartheta \sqrt{h}\mathrm{d}^2x$$ (2.39) which is an action for a sigma model into $`S^1`$ of radius $`1/R`$. Thus, we have shown that the sigma model into $`S^1`$ of radius $`R`$ is equivalent to the model with radius $`1/R`$. This is the $`R1/R`$ duality (which is called target space duality or T-duality in string theory). Comparing (2.33) and (2.38), we obtain the relation $$R\mathrm{d}\phi =i\frac{1}{R}\mathrm{d}\vartheta .$$ (2.40) Since $`R\mathrm{d}\phi `$ and $`iR\mathrm{d}\phi `$ are the conserved currents in the original system that count momentum and winding number respectively, the relation (2.40) means that momentum and winding number are exchanged under the $`R1/R`$ duality. In particular, the vertex operator $$\mathrm{exp}(i\vartheta )$$ (2.41) that creates a unit momentum in the dual theory must be equivalent to an operator that creates a unit winding number in the original theory. This can be confirmed by the following path integral manipulation. Let us consider the insertion of $$\mathrm{exp}\left(i_p^qB\right)$$ (2.42) in the system with the action (2.32), where the integration is along a path $`\tau `$ emanating from $`p`$ and ending on $`q`$. Then, using (2.38) we see that $$\mathrm{exp}\left(i_p^qB\right)=\mathrm{e}^{i\vartheta (q)}\mathrm{e}^{i\vartheta (p)}.$$ (2.43) On the other hand, the insertion of $`\mathrm{e}^{i_p^qB}`$ changes the $`B`$-linear term in (2.32). We note that $`_p^qB`$ can be expressed as $`_WB\omega `$, where $`\omega `$ is a one-form with delta function support along the path $`\tau `$. This $`\omega `$ can be written as $`\omega =\mathrm{d}\theta _\tau `$ where $`\theta _\tau `$ is a multi-valued function on $`W`$ that jumps by one when crossing the path $`\tau `$. Now, the modification of the action (2.32) can be written as $$\frac{i}{2\pi }_WB\mathrm{d}\phi \frac{i}{2\pi }_WB\mathrm{d}\phi +i_p^qB=\frac{i}{2\pi }_WB\mathrm{d}(\phi +2\pi \theta _\tau ).$$ (2.44) Integrating out $`B_\mu `$, we obtain the action (2.31) with $`\phi `$ replaced by $`\phi ^{}=\phi +2\pi \theta _\tau `$. Note that $`\phi ^{}`$ jumps by $`2\pi `$ when crossing the path $`\tau `$ which starts and ends on $`p`$ and $`q`$. In particular, it has winding number $`1`$ and $`1`$ around $`p`$ and $`q`$ respectively. Thus, the insertion of $`\mathrm{e}^{i\vartheta }`$ creates the unit winding number in the original system. ### Mirror Symmetry as $`R1/R`$ Duality We now proceed to a supersymmetric sigma model on the algebraic torus, or the cylinder $`𝐂^\times =𝐑\times S^1`$. We show that $`R1/R`$ duality performed on the $`S^1`$ factor is indeed a mirror symmetry. We work now in Minkowski signature. We denote the complex coordinate of the cylinder $`𝐑\times S^1`$ as $$\varphi =\varrho +i\phi $$ (2.45) where $`\varrho `$ is the coordinates of $`𝐑`$ and $`\phi `$ is the periodic coordinate of $`S^1`$ of period $`2\pi `$. The Lagrangian of the system is $$L=\mathrm{d}^4\theta \frac{R^2}{2}\|\mathrm{\Phi }\|^2=\frac{R^2}{2}\left(\eta ^{\mu \nu }_\mu \overline{\varphi }_\nu \varphi +i\overline{\psi }_{}(_0+_1)\psi _{}+i\overline{\psi }_+(_0_1)\psi _+\right),$$ (2.46) where $`\mathrm{\Phi }`$ is the chiral superfield whose lowest component is $`\varphi `$. The Kahler metric for $`\varphi `$ is $`\mathrm{d}s^2=R^2\|\mathrm{d}\varphi \|^2=R^2(\mathrm{d}\varrho ^2+\mathrm{d}\phi ^2)`$ so that $`S^1`$ has radius $`R`$.<sup>2</sup><sup>2</sup>2 In this paper, we take the convention $`S=\frac{1}{2\pi }\mathrm{d}^2xL`$ as the relation of the action and the Lagrangian. Thus, the weight factor in Path-Integral is $`\mathrm{exp}(\frac{i}{2\pi }\mathrm{d}^2xL)`$ (in Minkowski signature). We perform the duality transformation on $`\phi `$. As we have seen, this yields another periodic variable $`\vartheta `$ of period $`2\pi `$ with the Kinetic term $`(1/2R^2)\eta ^{\mu \nu }_\mu \vartheta _\nu \vartheta `$. Thus, the dual theory is also a sigma model into a cylinder, but with a metric $$\mathrm{d}\stackrel{~}{s}^2=\mathrm{d}\varrho ^2+\frac{1}{R^2}\mathrm{d}\vartheta ^2=\frac{1}{R^2}\left(R^4\mathrm{d}\varrho ^2+\mathrm{d}\vartheta ^2\right).$$ (2.47) Thus, either $`R^2\varrho +i\vartheta `$ or $`R^2\varrho i\vartheta `$ is the complex coordinates of the new cylinder. What is the superpartner of this (anti-)holomorphic variable? We note the supersymmetry transformations $`\delta \psi _\pm =i\sqrt{2}(_0\pm _1)\varphi \overline{ϵ}^\pm `$ and $`\delta \overline{\psi }_\pm =i\sqrt{2}(_0\pm _1)\varphi ϵ^\pm `$. From (2.40), we see (after continuation back to Minkowski signature by $`x^2=ix^0`$) that $`R^2(_0\pm _1)\phi =(_0\pm _1)\vartheta `$ and therefore $`R^2(_0+_1)\varphi =(_0+_1)\eta `$ and $`R^2(_0_1)\varphi =(_0_1)\overline{\eta }`$ where $$\eta =R^2\varrho i\vartheta .$$ (2.48) Thus, the supersymmetry transformation is expressed as $`R^2\delta \psi _+=i\sqrt{2}(_0+_1)\eta \overline{ϵ}^+,R^2\delta \overline{\psi }_+=i\sqrt{2}(_0+_1)\overline{\eta }ϵ^+,`$ (2.49) $`R^2\delta \psi _{}=i\sqrt{2}(_0_1)\overline{\eta }\overline{ϵ}^+,R^2\delta \overline{\psi }_{}=i\sqrt{2}(_0+_1)\eta ϵ^{}.`$ (2.50) This is not a supersymmetry transformation for a chiral multiplet, but that for a twisted chiral multiplet. Indeed, renaming the fermions as $$R^2\psi _\pm =\pm \overline{\chi }_\pm ,R^2\overline{\psi }_\pm =\pm \chi _\pm ,$$ (2.51) the Lagrangian for the dual theory becomes $`\mathrm{d}^4\theta (\frac{1}{2R^2}\|\mathrm{\Theta }\|^2)`$ for a twisted chiral superfield $`\mathrm{\Theta }=\eta +\sqrt{2}(\theta ^+\overline{\chi }_++\theta ^{}\chi _{})+\mathrm{}`$. Thus, we have seen that $`R1/R`$ duality on $`S^1`$ transforms a theory of a chiral multiplet to another theory of a twisted chiral multiplet. Thus, this is a mirror symmetry. The above manipulation can be simplified by performing the dualization in superspace. We follow the procedure developed in . We start with the following Lagrangian for a real superfield $`B`$ and a twisted chiral superfield $`\mathrm{\Theta }`$. $$L^{}=\mathrm{d}^4\theta \left(\frac{R^2}{4}B^2\frac{1}{2}(\mathrm{\Theta }+\overline{\mathrm{\Theta }})B\right)$$ (2.52) We first integrate over the twisted chiral field $`\mathrm{\Theta }`$, $`\overline{\mathrm{\Theta }}`$. This yields the following constraint on $`B`$ $$\overline{D}_+D_{}B=D_+\overline{D}_{}B=0,$$ (2.53) which is solved by $$B=\mathrm{\Phi }+\overline{\mathrm{\Phi }},$$ (2.54) where $`\mathrm{\Phi }`$ is a chiral superfield. Now, inserting this into the original Lagrangian we obtain the Lagrangian (2.46) $$L=\mathrm{d}^4\theta \frac{R^2}{4}(\mathrm{\Phi }+\overline{\mathrm{\Phi }})^2=\mathrm{d}^4\theta \frac{R^2}{2}\overline{\mathrm{\Phi }}\mathrm{\Phi }.$$ (2.55) for the sigma model into the cylinder with radius $`R`$ on $`S^1`$. Now, reversing the order of integration, we consider integrating out $`B`$ first. Then, $`B`$ is solved by $$B=\frac{1}{R^2}(\mathrm{\Theta }+\overline{\mathrm{\Theta }}).$$ (2.56) Inserting this into $`L^{}`$ we obtain $$\stackrel{~}{L}=\mathrm{d}^4\theta \left(\frac{1}{2R^2}\overline{\mathrm{\Theta }}\mathrm{\Theta }\right),$$ (2.57) which is again the Lagrangian for supersymmetric sigma model on the cylinder. This time, the radius of $`S^1`$ is $`1/R`$ and the complex coordinate is described by the twisted chiral superfield $`\mathrm{\Theta }`$. From (2.54) and (2.56), we obtain $`R^2(\mathrm{\Phi }+\overline{\mathrm{\Phi }})=\mathrm{\Theta }+\overline{\mathrm{\Theta }}`$ which reproduces the relation between the component fields obtained above (e.g. (2.51)). #### 2.2.2 Examples Here we present three classes of examples of (conjectural) mirror symmetry. They are mirror symmetry between LG model and LG model, sigma model and sigma model, and sigma model and LG model. ### Minimal Models and Orbifolds $`N=2`$ minimal models are the simplest class of exactly solvable $`(2,2)`$ SCFTs. It has been argued that the $`(d2)`$-th minimal model arises as the IR fixed point of the LG model of chiral superfield $`X`$ with the superpotential $$W=X^d.$$ (2.58) If we orbifold the model by the $`𝐙_d`$ symmetry generated by $`X\mathrm{e}^{2\pi i/d}X`$, we obtain again the $`(d2)`$-th minimal model $$\stackrel{~}{W}=\stackrel{~}{X}^d.$$ (2.59) This time, however, $`\stackrel{~}{X}`$ is a twisted chiral superfield and $`\stackrel{~}{W}`$ is a twisted superpotential. This means that the minimal model and its $`𝐙_d`$-orbifold can be considered as mirror to each other. If a CFT $`𝒞`$ has a discrete abelian symmetry group $`\mathrm{\Gamma }`$, the orbifold CFT $`𝒞^{}=𝒞/\mathrm{\Gamma }`$ has a symmetry group $`\mathrm{\Gamma }^{}`$ isomorphic to $`\mathrm{\Gamma }`$ and the orbifold $`𝒞^{}/\mathrm{\Gamma }^{}`$ is identical to the original CFT $`𝒞`$. We apply this general fact to the sum of $`N`$ copies of minimal models $$W=X_1^d+\mathrm{}+X_N^d,$$ (2.60) modulo its $`𝐙_d`$ symmetry group generated by $`X_iX_i\mathrm{e}^{2\pi i/d}X_i`$ for all $`i`$. This LG orbifold $`𝒞=\left(W=_iX_i^d\right)/𝐙_d`$ has a symmetry group $`\mathrm{\Gamma }(𝐙_d)^{N1}`$ generated by $`X_i\mathrm{e}^{2\pi i\alpha _i/d}X_i`$. Orbifold of $`𝒞`$ by this $`\mathrm{\Gamma }`$ is $`𝒞^{}=\left(W=_iX_i^d\right)/(𝐙_d)^N=_i\left(W=X_i^d\right)/𝐙_d`$. By the mirror symmetry of (2.58) mod $`𝐙_d`$ and (2.59), this is identical to the sum of $`N`$ copies of the minimal model given by the twisted chiral superpotential $$\stackrel{~}{W}=\stackrel{~}{X}_1^d+\mathrm{}+\stackrel{~}{X}_N^d.$$ (2.61) This indeed has a symmetry $`\mathrm{\Gamma }^{}`$ isomorphic to $`(𝐙_d)^{N1}`$ generated by $`\stackrel{~}{X_i}\mathrm{e}^{2\pi i\alpha _i/d}X_i`$ with $`_i\alpha _i=0`$ (mod $`d`$). By the general fact on the orbifold, we see that the orbifold $`\left(\stackrel{~}{W}=_i\stackrel{~}{X}_i^d\right)/(𝐙_d)^{N1}`$ is identical to the original SCFT $`\left(W=_iX_i^d\right)/𝐙_d`$. In other words, the $`𝐙_d`$ orbifold and the $`(𝐙_d)^{N1}`$ orbifold of the sum of $`N`$ copies of the minimal model are mirror to each other. ### Quintic Applying the above construction of mirror pair to the orbifold minimal models corresponding to a CY sigma model (at the special point of its moduli space), Greene and Plesser constructed pairs of CY manifold whose sigma models are mirror to each other . The connection between Landau-Ginzburg models and Calabi-Yau sigma models was first discussed in . The derivation of this connection in involved a change of variables in field space, as if one were dealing with an ordinary integral. This heuristic derivation of the relation between LG models and Calabi-Yau sigma models was made precise by Cecotti who showed the arguments are precise in the context of computing periods associated to special geometry of the LG model, and that the derivation of can be viewed as showing the equivalence of periods and special geometry of the Calabi-Yau with an LG model (in modern terminology as far as the middle dimensional D-brane masses are concerned). In fact this identification of LG models and special geometry associated to the vacuum geometry of the sigma model is crucial in our derivation of mirror symmetry for the case of complete intersections in toric varieties. The connection between LG models and Calabi-Yau sigma models was further elucidated in using the linear sigma model description, which we will review in section 4 in this paper. ### $`𝐂\mathrm{P}^{N1}`$ and Affine Toda Thoery Less known class of mirror symmetry is between non-liner sigma models on manifolds of positive first Chern class and LG models without scale invariance. The typical example is the mirror symmetry of the $`𝐂\mathrm{P}^{N1}`$ sigma model and supersymmetric $`A_{N1}`$ affine Toda theory. The $`𝐂\mathrm{P}^{N1}`$ model is asymptotic free and generates a dynamical scale $`\mathrm{\Lambda }`$. It has $`N`$ vacua with mass gap. $`U(1)_V`$ is unbroken but $`U(1)_A`$ is anomalously broken to $`𝐙_{2N}`$ which is spontaneously broken to $`𝐙_2`$. The $`A_{N1}`$ affine Toda theory is an LG model of $`N1`$ periodic variables $`X_i`$ having superpotential $$W=\mathrm{\Lambda }\left(\mathrm{e}^{X_1}+\mathrm{}+\mathrm{e}^{X_{N1}}+\underset{i=1}{\overset{N1}{}}\mathrm{e}^{X_i}\right).$$ (2.62) This theory has the same properties of the $`𝐂\mathrm{P}^{N1}`$ model mentioned above, except that $`U(1)_V`$ and $`U(1)_A`$ are exchanged, the mass scale is explicitly introduced and $`U(1)_V𝐙_{2N}`$ is an explicit breaking. This duality is in many ways an $`N=2`$ generalization of the duality of the (bosonic) sine-Gordon theory and the massive Thirrling model . In particular, solitons of the affine Toda theory are mapped to the fundamental fields in the $`𝐂\mathrm{P}^{N1}`$ model (in the linear sigma model realization) . The equivalence of the two theories has been observed from various points of view. The agreement of BPS soliton spectrum and their scattering matrix , of correlation functions of topologically twisted theories (coupled to topological gravity) . This example is extended to a more general class of manifolds in . ## 3 The Dynamics of $`N=2`$ Gauge Thoeries in Two Dimensions In this section, we study the dynamics of $`(2,2)`$ supersymmetric gauge theories in $`1+1`$ dimensions. We consider $`U(1)`$ gauge theory with charged chiral multiplets and assume that there is no superpotential for the charged fields. The theory is described by vector superfield $`V`$ with field strength $`\mathrm{\Sigma }`$ which is a twisted chiral superfield. We consider $`N`$ chiral superfields $`\mathrm{\Phi }_i`$ of charge $`Q_i`$. For earlier studies of this class of theories, see . The classical theory is parametrized by the gauge coupling $`e`$, Fayet-Iliopoulos (FI) parameter $`r`$ and Theta angle $`\theta `$ where $`e`$ has dimension of mass but $`r`$ and $`\theta `$ are dimensionless. The Lagrangian of the theory is given by $$L=\mathrm{d}^4\theta (\underset{i=1}{\overset{N}{}}\overline{\mathrm{\Phi }}_i\mathrm{e}^{2Q_iV}\mathrm{\Phi }_i\frac{1}{2e^2}\overline{\mathrm{\Sigma }}\mathrm{\Sigma })+\frac{1}{2}(\mathrm{d}^2\stackrel{~}{\theta }t\mathrm{\Sigma }+c.c.),$$ (3.1) where $`t`$ is the complex combination $$t=ri\theta .$$ (3.2) The theory is super-renormalizable with respect to the gauge coupling $`e`$. However, the FI parameter $`r`$ is renormalized to cancell a one-loop divergence unless $`_iQ_i=0`$. The dependence of the bare parameter $`r_0`$ on the cut-off $`\mathrm{\Lambda }_{UV}`$ is $$r_0=\underset{i=1}{\overset{N}{}}Q_i\mathrm{log}\left(\frac{\mathrm{\Lambda }_{UV}}{\mathrm{\Lambda }}\right)$$ (3.3) where $`\mathrm{\Lambda }`$ is a scale parameter that replaces $`r`$ as the parameter of the theory if $`_iQ_i0`$. The classical theory is invariant under $`U(1)_V\times U(1)_A`$ R-symmetry where $`\mathrm{\Sigma }`$ is assigned an axial charge $`2`$ and zero vector charge. $`U(1)_V`$ is an exact symmetry of the theory but $`U(1)_A`$ is subject to the chiral anomaly. The axial rotation by $`\mathrm{e}^{i\alpha }`$ shifts the theta angle by $$\theta \theta 2\underset{i=1}{\overset{N}{}}Q_i\alpha .$$ (3.4) Thus $`U(1)_A`$ is unbroken if $`_iQ_i=0`$ but otherwise is broken to the discrete subgroup $`𝐙_{2p}`$ with $`p=_iQ_i`$. In addition, the theory has other global symmetries. There are at least $`N1`$ $`U(1)`$ symmetries which are the phase rotation of the $`N`$ chiral superfield modulo $`U(1)`$ gauge transformations. This will be important in our study. Of course there could be larger symmetry if some of the $`U(1)`$ charges $`Q_i`$ coincide. These global symmetries are non-anomalous and are the symmetries of the quantum theory. In what follows, we study the dynamics of this gauge theory. We first dualize each of the charged chiral fields $`\mathrm{\Phi }_i`$ using the phase rotation symmetry. We then study the effective theory described in terms of the dual variables. The goal of this section is to show that the twisted superpotential as mentioned in the introduction of the paper is dynamically generated. ### 3.1 Abelian Duality Let us consider a complex scalar field $`\varphi `$ which is minimally coupled to a gauge field $`A_\mu `$. In terms of the polar variables $`(\rho ,\phi )`$ defined by $`\varphi =\rho \mathrm{e}^{i\phi }`$, the kinetic term $`\eta ^{\mu \nu }D_\mu \varphi ^{}D_\nu \varphi `$ is written as the sum of $`(_\mu \rho )^2`$ and $$L_\phi =\rho ^2(_\mu \phi +QA_\mu )^2,$$ (3.5) where $`Q`$ is the charge of $`\varphi `$. This Lagrangian is invariant under the shift $`\phi \phi +`$constant, and therefore we can consider dualizing $`\phi `$ as we have done when we discussed $`R1/R`$ duality. What is new here is that we have a gauge field coupled to $`\phi `$ as in (3.5). The dualization procedure start with the following Lagrangian for the vector field $`B_\mu `$, the angle variable $`\phi `$, plus the gauge field $`A_\mu `$: $$L^{}=\frac{1}{4\rho ^2}(B_\mu )^2+ϵ^{\mu \nu }B_\mu (_\nu \phi +QA_\nu ).$$ (3.6) Integration over $`B`$ yields the Lagrangian (3.5). If, instead, we first integrate over $`\phi `$, we obtain the constraint $`B_\mu =_\mu \vartheta `$ where $`\vartheta `$ is an angle variable of period $`2\pi `$. Plugging this into $`L^{}`$, we obtain the new Lagrangian $$L_\vartheta =\frac{1}{4\rho ^2}(_\mu \vartheta )^2+Qϵ^{\mu \nu }_\mu \vartheta A_\nu =\frac{1}{4\rho ^2}(_\mu \vartheta )^2Q\vartheta ϵ^{\mu \nu }\frac{1}{2}F_{\mu \nu },$$ (3.7) where a partial integration is used. Thus, the dual variable $`\vartheta `$ is coupled to the gauge field $`A_\mu `$ as a dynamical Theta angle. ### Supersymmetric Case It is straightforward to repeat this dualization in the supersymmetric theory with a chiral superfield $`\mathrm{\Phi }`$ of charge $`Q`$. In fact, we only have to dualize the phase of $`\mathrm{\Phi }`$ and suitably rename other fields. As we have seen above, the dual variable couples to the gauge field as a dynamical Theta angle. Such a variable must be in a twisted chiral multiplet $`Y`$ that couples to the field strength $`\mathrm{\Sigma }`$ in the twisted superpotential as $$QY\mathrm{\Sigma }.$$ (3.8) One can see this explicitly by performing the duality transformation in the superspace, as we now show. We start with the following Lagrangian for a vector superfield $`V`$, a real superfield $`B`$ and a twisted chiral superfield $`Y`$ whose imaginary part is periodic with period $`2\pi `$. $$L^{}=\mathrm{d}^4\theta \left(\mathrm{e}^{2QV+B}\frac{1}{2}(Y+\overline{Y})B\right),$$ (3.9) where $`Q`$ is an integer. We first integrate over $`Y`$. This yields the constraint $`\overline{D}_+D_{}B=D_+\overline{D}_{}B=0`$ on $`B`$ which is solved by $$B=\mathrm{\Psi }+\overline{\mathrm{\Psi }},$$ (3.10) where $`\mathrm{\Psi }`$ is a chiral superfield. Since the imaginary part of $`Y`$ is an angular variable of period $`2\pi `$, so is the imaginary part of $`\mathrm{\Psi }`$. Now, inserting this into the original Lagrangian we obtain $$L=\mathrm{d}^4\theta \mathrm{e}^{2QV+\mathrm{\Psi }+\overline{\mathrm{\Psi }}}$$ (3.11) which is nothing but the Lagrangian for the chiral superfield $`\mathrm{\Phi }=\mathrm{e}^\mathrm{\Psi }`$ of charge $`Q`$. Now, reversing the order of integration, we consider integrating out $`B`$ first. Then, $`B`$ is solved by $$B=2QV+\mathrm{log}\left(\frac{Y+\overline{Y}}{2}\right).$$ (3.12) Inserting this into $`L^{}`$ we obtain $$\stackrel{~}{L}=\mathrm{d}^4\theta \left(QV(Y+\overline{Y})\frac{1}{2}(Y+\overline{Y})\mathrm{log}(Y+\overline{Y})\right)$$ (3.13) Using the fact that $`Y`$ is a twisted chiral superfield, $`\overline{D}_+Y=D_{}Y=0`$, the term proportional to $`V`$ can be written as $`{\displaystyle \mathrm{d}^4\theta VY}`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle d\theta ^+d\overline{\theta }^{}\overline{D}_+D_{}VY}`$ (3.14) $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \mathrm{d}^2\stackrel{~}{\theta }\mathrm{\Sigma }Y}`$ (3.15) where we have used $`\mathrm{\Sigma }=\overline{D}_+D_{}V`$ which holds for abelian gauge group. Together with the gauge kinetic term and the classical FI-Theta terms, we obtain the following Lagrangian <sup>1</sup><sup>1</sup>1In the usual T-duality the dilaton shifts, proportional to the volume of the space. In the case at hand, since translations in the dualizing circle is gauged this shift in dilaton does not arise.: $$\stackrel{~}{L}=\mathrm{d}^4\theta \{\frac{1}{2e^2}\overline{\mathrm{\Sigma }}\mathrm{\Sigma }\frac{1}{2}(Y+\overline{Y})\mathrm{log}(Y+\overline{Y})\}+\frac{1}{2}(\mathrm{d}^2\stackrel{~}{\theta }\mathrm{\Sigma }(QYt)+c.c.)$$ (3.16) We indeed see that the charged chiral superfield $`\mathrm{\Phi }`$ has turned into a neutral twisted chiral superfield $`Y`$ which couples to the field strength $`\mathrm{\Sigma }`$ as the dynamical Theta angle. It follows from (3.10) and (3.12) that the original chiral field $`\mathrm{\Phi }`$ and the twisted chiral field $`Y`$ are related by $$Y+\overline{Y}=2\overline{\mathrm{\Phi }}\mathrm{e}^{2QV}\mathrm{\Phi }.$$ (3.17) We see that the dual field $`Y`$ is a gauge invariant composite of the original field $`\mathrm{\Phi }`$. Using the expression (2) of $`V`$ in the Wess-Zumino gauge, we can write down the relation between the components fields $`y=\varrho i\vartheta ,\overline{\chi }_+,\chi _{}`$ of $`Y`$ and those $`\varphi =\rho \mathrm{e}^{i\phi },\psi _+,\psi _{}`$ of $`\mathrm{\Phi }`$: $`\varrho `$ $`=`$ $`\rho ^2,`$ (3.18) $`_\pm \vartheta `$ $`=`$ $`\pm 2\left(\rho ^2(_\pm \phi +QA_\pm )+\overline{\psi }_\pm \psi _\pm \right)`$ (3.19) where $`_\pm =_0\pm _1`$ etc, and $`\overline{\chi }_+=2\varphi ^{}\psi _+,`$ $`\chi _{}=2\overline{\psi }_{}\varphi ,`$ (3.20) $`\chi _+=2\overline{\psi }_+\varphi ,`$ $`\overline{\chi }_{}=2\varphi ^{}\psi _{}.`$ (3.21) We note that the term $`\pm 2\overline{\psi }_\pm \psi _\pm `$ in (3.19) reflects the fact that we are dualizing on the phase of the whole superfield $`\mathrm{\Phi }`$. The Kahler metric of the field $`Y`$ is given by $$\mathrm{d}s^2=\frac{\|\mathrm{d}y\|^2}{2(y+\overline{y})}=\frac{1}{4\varrho }(\mathrm{d}\varrho ^2+\mathrm{d}\vartheta ^2)=\mathrm{d}\rho ^2+\frac{1}{4\rho ^2}\mathrm{d}\vartheta ^2,$$ (3.22) as can also be seen from the bosonic treatment (e.g. (3.7)). We note that the relation (3.18) implies a condition that the real part of $`y`$ is allowed to take only non-negative values, $`\mathrm{Re}(y)0`$. The boundary $`\mathrm{Re}(y)=0`$ corresponds to $`\|\mathrm{\Phi }\|=0`$ where the circle on which we are dualizing shrinks to zero size. With respect to the metric (3.22), the boundary is at finite distance from any point with finite $`\mathrm{Re}(y)`$. Naively we expect a lot of singularities coming from such end points. However, it will be argued below that the physically relevant region is infinitely far away from such a boundary region, once the renormalization is taken into account. ### Renormalization We recall that we had to renormalize the FI parameter of the theory as (3.3). In order for the coupling $`\mathrm{\Sigma }(QYt)`$ to be finite, we have to renormalize also the field $`Y`$. This can be done by letting the bare dual field $`Y_0`$ to depend on the UV cut-off as $$Y_0=\mathrm{log}(\mathrm{\Lambda }_{UV}/\mu )+Y,$$ (3.23) where $`\mu `$ is the scale at which the field is renormalized. Note that this resolves the issue of the bound $`\mathrm{Re}(y)0`$. In fact, the correct condition is $`\mathrm{Re}(y_0)0`$ and therefore it means $$\mathrm{Re}(y)\mathrm{log}(\mathrm{\Lambda }_{UV}/\mu )$$ (3.24) for the renormalized variable. In particular, in the continuum limit $`\mathrm{\Lambda }_{UV}/\mu \mathrm{}`$, there is no bound on the renormalized field. Also, we note that the Kahler metric for the renormalized field is $$\mathrm{d}s^2=\frac{\|\mathrm{d}y\|^2}{2(2\mathrm{log}(\mathrm{\Lambda }_{UV}/\mu )+y+\overline{y})}\frac{\|\mathrm{d}y\|^2}{4\mathrm{log}(\mathrm{\Lambda }_{UV}/\mu )}$$ (3.25) which becomes flat in the continuum limit. ### R-Symmetry We would like to know the transformation property of the new field $`Y`$ under the vector and axial R-symmetries. It appears from the relation (3.17) that the superfield $`Y`$ transforms as a charge zero field under both $`U(1)_V`$ and $`U(1)_A`$. However, it cannot be directly seen from (3.17) how the imaginary part $`\vartheta `$ of $`Y`$ transforms. Nevertheless, one can read the transformation of $`\vartheta `$ from (3.17) or (3.19) in a way similar to . Let us note that the conserved currents of the vector and axial R-symmetries are given by $$J_\pm ^V=\overline{\psi }_\pm \psi _\pm +\mathrm{},J_\pm ^A=\pm \overline{\psi }_\pm \psi _\pm +\mathrm{},$$ (3.26) where $`+\mathrm{}`$ are contributions from the vector multiplet fields. The operator product of these with $`_\pm \vartheta `$ expressed as (3.19) has the following singularity: $`J_\pm ^V(x)_\pm \vartheta (y){\displaystyle \frac{\pm 2}{(x^\pm y^\pm )^2}},`$ (3.27) $`J_\pm ^A(x)_\pm \vartheta (y){\displaystyle \frac{2}{(x^\pm y^\pm )^2}}.`$ (3.28) These show that $`\vartheta `$ is invariant under vector R-symmetry but is shifted by $`2`$ by the axial R-symmetry. Therefore the superfield $`Y`$ transforms under $`U(1)_V\times U(1)_A`$ as $`\mathrm{e}^{i\alpha F_V}Y(\theta ^\pm ,\overline{\theta }^\pm )\mathrm{e}^{i\alpha F_V}=Y_i(\mathrm{e}^{i\alpha }\theta ^\pm ,\mathrm{e}^{i\alpha }\overline{\theta }^\pm ),`$ (3.29) $`\mathrm{e}^{i\alpha F_A}Y(\theta ^\pm ,\overline{\theta }^\pm )\mathrm{e}^{i\alpha F_A}=Y(\mathrm{e}^{i\alpha }\theta ^\pm ,\mathrm{e}^{\pm i\alpha }\overline{\theta }^\pm )2i\alpha .`$ (3.30) Indeed, the Lagrangian (3.16) exhibits the $`U(1)_V`$ invariance and the correct $`U(1)_A`$ anomaly under this action. ### Multi-Flavor Case It is straightforward to extend the dualization considered above to the case where there are several charged chiral fields. Dualizing each of the chiral fields $`\mathrm{\Phi }_i`$, we obtain the following twisted superpotential for the twisted chiral fields $`Y_{i0}`$ $$\stackrel{~}{W}=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_iY_{i0}t_0\right).$$ (3.31) The relation between $`\mathrm{\Phi }_i`$ and $`Y_{i0}`$ are the same as in (3.17)-(3.21). We renormalize the fields as $$Y_{i0}=\mathrm{log}(\mathrm{\Lambda }_{UV}/\mu )+Y_i.$$ (3.32) so that (3.31) is finite in the continuum limit $`\mathrm{\Lambda }_{UV}\mathrm{}`$ in the case $`_iQ_i0`$. In the case $`_iQ_i=0`$, (3.31) is invariant under an $`i`$-independent shift of $`Y_{i0}`$’s and we also do this field redefinition (3.32). In any case, the bound $`\mathrm{Re}(y_{i0})0`$ is eliminated from $`y_i`$’s. With respect to the renormalized fields the twisted superpotential can be written as $$\stackrel{~}{W}=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_iY_it(\mu )\right).$$ (3.33) where $`t(\mu )`$ is the effective FI-Theta parameter $$t(\mu )=\{\begin{array}{cc}_{i=1}^NQ_i\mathrm{log}(\mu /\mathrm{\Lambda })i\theta \hfill & \text{if }_{i=1}^NQ_i0,\hfill \\ ri\theta \hfill & \text{if }_{i=1}^NQ_i=0.\hfill \end{array}$$ (3.34) As in the single flavor case, one can see that the R-symmetry group $`U(1)_V\times U(1)_A`$ acts on the fields $`Y_i`$ as $`\mathrm{e}^{i\alpha F_V}Y_i(\theta ^\pm ,\overline{\theta }^\pm )\mathrm{e}^{i\alpha F_V}=Y_i(\mathrm{e}^{i\alpha }\theta ^\pm ,\mathrm{e}^{i\alpha }\overline{\theta }^\pm ),`$ (3.35) $`\mathrm{e}^{i\alpha F_A}Y_i(\theta ^\pm ,\overline{\theta }^\pm )\mathrm{e}^{i\alpha F_A}=Y_i(\mathrm{e}^{i\alpha }\theta ^\pm ,\mathrm{e}^{\pm i\alpha }\overline{\theta }^\pm )2i\alpha .`$ (3.36) We see that the superpotential (3.33) is invariant under $`U(1)_V`$. Note also that (3.33) exhibits the axial anomaly (3.4) for $`p=_iQ_i0`$ and the breaking of $`U(1)_A`$ down to $`𝐙_{2p}`$. It may appear that there are extra symmetries $`Y_iY_i+c_i`$ with $`c_ic_j`$ and $`_{i=1}^NQ_ic_i=0`$ that makes the transformation (3.36) ambiguous. However, there is no corresponding symmetry in the original system; $`Y_i`$’s can be shifted only by axial symmetries and the only axial symmetry in the system is the $`U(1)_A`$ with current $`J_\pm ^A=\pm _{i=1}^N\overline{\psi }_{i\pm }\psi _{i\pm }+\mathrm{}`$ as long as $`Q_i`$’s are all non-zero. Thus, there is no room for ambiguity in the R-transformation (3.36). In fact, new terms in $`\stackrel{~}{W}`$ that violates the invariance under the extra shifts with $`c_ic_j`$ is generated, as we will show next. ### 3.2 Dynamical Generation Of Superpotential The twisted superpotential (3.33) obtained from dualization is an exact expression in perturbation theory with resepect to $`1/r`$; it is simply impossible to write down a perturbative correction that respects the R-symmetry and/or anomaly, holomorphy in $`t`$, and periodicity of Theta angle. The D-term is of course subject to perturbative corrections. However, the twisted superpotential is possibly corrected by non-perturbative effects. A typical non-perturbative effect in quantum field theory is by the presence of instantons. The bosonic part of our theory is an abelian Higgs model which can have an instanton configuration — vortex. It has been known that in an abelian Higgs model a Theta dependent vacuum energy density is generated by the effect of the gas of vortices and anti-vortices . As in that case, and also as in Polyakov’s model of confinement where a bosonic potential for the dual field is generated from the gas of monopoles and anti-monopoles, we expect that a superpotential for $`Y_i`$’s can be generated by the gas of vortices and anti-vortices. Around the vortex for a charged scalar $`\varphi _i`$, the phase of $`\varphi _i`$ has winding number one. As we have seen in $`R1/R`$ duality, a winding configuration is dual to the insertion of the vertex operator $`\mathrm{e}^{i\vartheta _i}`$. The supersymmetric completion of this operator is the twisted chiral superfield $$\mathrm{e}^{Y_i}.$$ (3.37) These exponentials have vector R-charge $`0`$ and axial R-charge $`2`$, as can be seen from (3.35) and (3.36). Thus, we can add these to the twisted superpotential without violating the $`U(1)_V`$ R-symmetry, and maintaining the correct anomaly of $`U(1)_A`$. In what follows we shall show that a correction of the form (3.37) is indeed generated. In fact we will show that the correction is simply the sum of them and the exact superpotential is given by $$\stackrel{~}{W}=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_iY_it(\mu )\right)+\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}.$$ (3.38) This is one of the main results of this paper. Note that the change of the renormalization scale $`\mu `$ can be absorbed by the shift of $`Y_i`$’s dictated by (3.32). The actual parameter of the theory is still the dynamical scale $`\mathrm{\Lambda }`$ for $`_iQ_i0`$ and the FI-Theta parameter $`t`$ for $`_iQ_i=0`$. Before embarking on the computation to show that $`\mathrm{e}^{Y_i}`$’s are indeed generated, we make a simple consistency check of (3.38). Let us consider integrating out $`Y_i`$’s for a fixed configuration of $`\mathrm{\Sigma }`$ whose lowest component of $`\mathrm{\Sigma }`$ is large and slowly varying. The variation with respect to $`Y_i`$’s yields the relation $`Q_i\mathrm{\Sigma }\mu \mathrm{e}^{Y_i}=0`$ or $`Y_i=\mathrm{log}(Q_i\mathrm{\Sigma }/\mu )`$. Inserting this into (3.38) we obtain the effective superpotential for the $`\mathrm{\Sigma }`$ field $$\stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}(\mathrm{\Sigma })=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_i\left(\mathrm{log}\left(Q_i\mathrm{\Sigma }/\mu \right)1\right)t(\mu )\right).$$ (3.39) This is nothing but what we would obtain when we integrate out the chiral superfield $`\mathrm{\Phi }_i`$’s in the original gauge theory for a fixed configuration of $`\mathrm{\Sigma }`$ . #### 3.2.1 Localization We first establish that the validity of (3.38) for general $`N`$ and $`Q_i`$ is a consequence of the case with $`N=1`$. This can be seen by considering a theory with a larger gauge symmetry where the $`U(1)^{N1}`$ global symmetries are gauged and recovering the original theory in the weak coupling limit of the extra gauge interactions. Sometimes we will refer to this procedure as localization. We start with the $`U(1)^N`$ gauge theory with a single charged matter for each $`U(1)`$ which is described by the following classical Lagrangian $$L=\underset{i,j}{}\mathrm{d}^4\theta \frac{1}{2e_{\overline{ı}j}^2}\overline{\mathrm{\Sigma }}_i\mathrm{\Sigma }_j+\underset{i=1}{\overset{N}{}}\{\mathrm{d}^4\theta \overline{\mathrm{\Phi }}_i\mathrm{e}^{2Q_iV_i}\mathrm{\Phi }_i+\frac{1}{2}(\mathrm{d}^2\stackrel{~}{\theta }t_i\mathrm{\Sigma }_i+c.c.)\}.$$ (3.40) If the gauge coupling matrix is diagonal $$\frac{1}{e_{\overline{ı}j}^2}=\delta _{i,j}\frac{1}{e_i^2},$$ (3.41) the $`N`$ single flavor theories are decoupled from each other. In such a case, the exact twisted superpotential is given by the sum of those for single flavor theories. If we assume that (3.38) holds for the single flavor cases, it is $$\stackrel{~}{W}=\underset{i=1}{\overset{N}{}}\left(\mathrm{\Sigma }_i\left(Q_iY_it_i(\mu )\right)+\mu \mathrm{e}^{Y_i}\right),$$ (3.42) where we have chosen the scale $`\mu `$ common to all $`i`$, and $`t_i(\mu )`$ is the effective FI-Theta parameter at $`\mu `$ $$t_i(\mu )=Q_i\mathrm{log}(\mathrm{\Lambda }_i/\mu )i\theta _i,$$ (3.43) for the $`i`$-th theory. Now the essential point here is that the (twisted) superpotential is independent of variation of the D-term. Thus, (3.42) is valid for all values of $`1/e_{\overline{ı}j}^2`$ as long as there is no singularity in going from the diagonal one (3.41). In particular, let us consider the coupling matrix of the following form $$\underset{i,j}{}\frac{1}{2e_{\overline{ı}j}^2}\overline{\mathrm{\Sigma }}_i\mathrm{\Sigma }_j=\frac{1}{2e^2}\left\|\frac{1}{N}\underset{i=1}{\overset{N}{}}\mathrm{\Sigma }_i\right\|^2+\frac{1}{ϵe^2}\underset{i=1}{\overset{N1}{}}\|\mathrm{\Sigma }_i\mathrm{\Sigma }_{i+1}\|^2.$$ (3.44) In the limit $$ϵ0,$$ (3.45) the only dynamical gauge symmetry is the diagonal $`U(1)`$ subgroup of $`U(1)^N`$ with the field strength given by $$\mathrm{\Sigma }_1+\mathrm{\Sigma }_2+\mathrm{}+\mathrm{\Sigma }_N=:N\mathrm{\Sigma },$$ (3.46) and we recover the $`U(1)`$ gauge theory with $`N`$ matter fields $`\mathrm{\Phi }_i`$ with charge $`Q_i`$. To be precise, the limit (3.45) itself does not completely fix $`\mathrm{\Sigma }_i\mathrm{\Sigma }_j`$, but rather sets it to be a constant. In other words we have the choice $`\mathrm{\Sigma }_i=\mathrm{\Sigma }+\mathrm{\Delta }_i`$ with $`\mathrm{\Delta }_i`$ being a constant. However, such a shift $`\mathrm{\Delta }_i`$ would yield a twisted superpotential whose perturbative part does not match with what we have (3.33) for the single $`U(1)`$ gauge theory. Therefore we must have $`\mathrm{\Delta }_i=0`$ in the present theory. Nevertheless as we will discuss later a non-zero $`\mathrm{\Delta }_i`$ is indeed allowed when we consider a perturbation of the theory with “twisted masses”. Therefore we will in general take into account the above possible deformation of the $`U(1)`$ gauge theory. The bare FI and Theta parameter of the $`U(1)`$ gauge theory is related to those of the $`U(1)^N`$ theories by $$t_0=\underset{i=1}{\overset{N}{}}(Q_i\mathrm{log}(\mathrm{\Lambda }_{UV}/\mathrm{\Lambda }_i)i\theta _i).$$ (3.47) In particular, for the theory with $`_{i=1}^NQ_i0`$, the dynamical scale $`\mathrm{\Lambda }`$ is given by $`_i\mathrm{\Lambda }^{Q_i}=_i\mathrm{\Lambda }_i^{Q_i}`$, while the FI parameter of the theory with $`_{i=1}^NQ_i=0`$ is $`r=_{i=1}^NQ_i\mathrm{log}(\mathrm{\Lambda }_i)`$. Then, the effective coupling $`t(\mu )`$ at energy $`\mu `$ is simply the sum $`_{i=1}^Nt_i(\mu )`$. Thus, the twisted superpotential (3.42) becomes (3.38) in the limit (3.45). This shows that (3.38) for general $`N`$ and $`Q_i`$ follows from the $`N=1`$ case. Thus, to show (3.38) we only have to show the single flavor case. We note also that in the single flavor case we only have to show it in the case of unit charge $`Q=1`$. Other cases just follow from that case by a redefinition of the gauge field and the FI-Theta parameter; $`Q\mathrm{\Sigma }\mathrm{\Sigma }`$, $`Qtt`$. #### 3.2.2 The Generation Of Superpotential Now, let us consider the single flavor case with $`Q=1`$. We first determine the possible form of the non-perturbative correction $`\mathrm{\Delta }\stackrel{~}{W}`$ from the general requirements — holomorphy in $`t`$, periodicity in Theta angle, R-symmetry, and asymptotic behaviour. Let $`\mathrm{\Lambda }=\mu \mathrm{e}^t`$ be the dynamical scale and let us put $`\stackrel{~}{Y}=Yt`$. Since $`t`$ and $`Y`$ are periodic with period $`2\pi i`$, $`\mathrm{\Delta }\stackrel{~}{W}`$ must be a holomorphic function of $`\mathrm{\Sigma }`$, $`\mathrm{\Lambda }`$ and $`\mathrm{e}^{\stackrel{~}{Y}}`$. The anomaly (3.4) of the axial R-symmetry is absorbed by the shift of the Theta angle $`\theta \theta +2\alpha `$, or $`tt2i\alpha `$. This modified $`U(1)_A`$ symmetry transforms the variables as $$\mathrm{\Sigma }\mathrm{e}^{2i\alpha }\mathrm{\Sigma },\mathrm{\Lambda }\mathrm{e}^{2i\alpha }\mathrm{\Lambda },\mathrm{e}^{\stackrel{~}{Y}}\mathrm{e}^{\stackrel{~}{Y}}.$$ (3.48) Since the twisted superpotential must have charge $`2`$ under this transformation, $`\mathrm{\Delta }\stackrel{~}{W}`$ must be of the form $`\mathrm{\Sigma }f(\mathrm{\Lambda }/\mathrm{\Sigma },\mathrm{e}^{\stackrel{~}{Y}})`$ which is expanded in a Laurent series as $$\mathrm{\Delta }\stackrel{~}{W}=\mathrm{\Sigma }\underset{n,m}{}c_{n,m}(\mathrm{\Lambda }/\mathrm{\Sigma })^n\mathrm{e}^{m\stackrel{~}{Y}}=\underset{n,m}{}c_{n,m}\mathrm{\Sigma }^{1n}\mu ^n\mathrm{e}^{(nm)tmY}.$$ (3.49) Now, we recall that the field $`Y`$ was introduced by dualizing the circle of radius $`\|\varphi \|`$. Therefore, in the description in terms of $`Y`$ and $`\mathrm{\Sigma }`$, we are looking at the region $`\varphi 0`$ of the field space where the gauge symmetry is broken and $`\mathrm{\Sigma }`$ is therefore massive. Thus, the twisted superpotential must be analytic at $`\mathrm{\Sigma }=0`$. This means that only terms of $`1n0`$ is non-vanishing in (3.49). On the other hand, in the semi-classical limit where $`r`$ is very large, the correction must be small compared to the perturbative term $`\mathrm{\Sigma }(Yt)`$. This requires $`nm`$. Finally, since $`Y`$ is unbounded in real positive direction (but is bounded from below as $`\mathrm{Re}Y\|\mathrm{\Phi }\|^20`$ in the semi-classical description), $`m0`$ is also required for the correction to be small. To summarize, only $`1nm0`$ is allowed. $`n=m=0`$ is of the same order as the perturbative term and is excluded. $`n=1,m=0`$ is just a constant term. The final candidate $`n=m=1`$ is a non-trivial term which is $`\mathrm{e}^Y`$. Thus, the twisted superpotential must be of the form $$\stackrel{~}{W}=\mathrm{\Sigma }(Yt(\mu ))+c\mu \mathrm{e}^Y,$$ (3.50) where $`c`$ is a dimensionless constant. The question is therefore whether the coefficient $`c`$ is zero or not. We have already observed an evidence that supports non-vanishing of $`c`$; Integration over $`Y`$ yields the correct effective superpotential for $`\mathrm{\Sigma }`$, $`\stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}(\mathrm{\Sigma })=\mathrm{\Sigma }\mathrm{log}(\mathrm{\Sigma }/\mathrm{\Lambda })`$, if $`c`$ is non-zero. In what follows we show that the term $`\mu \mathrm{e}^Y`$ is indeed generated by an instanton effect. ### The Vortex We continue our gauge theory to the Euclidean signature by Wick rotation $`x^0=ix^2`$. We choose the orientation so that $`z=x^1+ix^2`$ is the complex coordinates. This leads to $`F_{01}iF_{12}`$, $`D_0+D_12D_{\overline{z}}`$ and $`D_0D_12D_z`$. After solving for the auxiliary field as $`D=e^2(\|\varphi \|^2r_0)`$ and $`F=0`$, the Euclidean action is given by $`S_E={\displaystyle \frac{1}{2\pi }}{\displaystyle \mathrm{d}^2x}`$ $`(\|D_\mu \varphi \|^2+\|\sigma \varphi \|^2+{\displaystyle \frac{1}{2e^2}}\|_\mu \sigma \|^2+{\displaystyle \frac{1}{2e^2}}(F_{12}^2+D^2)+i\theta F_{12}`$ (3.51) $`2i\overline{\psi }_{}D_{\overline{z}}\psi _{}+2i\overline{\psi }_+D_z\psi _++\overline{\psi }_{}\sigma \psi _++\overline{\psi }_+\overline{\sigma }\psi _{}`$ $`+{\displaystyle \frac{1}{e^2}}\left(i\overline{\lambda }_{}_{\overline{z}}\lambda _{}+i\overline{\lambda }_+_z\lambda _+\right)`$ $`+i(\varphi ^{}\lambda _{}\psi _+\varphi ^{}\lambda _+\psi _{}\overline{\psi }_+\overline{\lambda }_{}\varphi +\overline{\psi }_{}\overline{\lambda }_+\varphi )).`$ An instanton is a topologically non-trivial configuration that minimizes the bosonic part of this action. An instanton can contribute to the (twisted) superpotential only when it carries two fermionic zero modes of the right kind. Since a twisted F-term is obtained by the integration over two fermionic coordinates other than $`\theta ^{}`$ and $`\overline{\theta }^+`$, a relevant configuration must be invariant under the supercharges $`Q_{}`$ and $`\overline{Q}_+`$. The invariance of the fermions under these supercharges requires (see Appendix B) $`\sigma =0,`$ (3.52) $`D_{\overline{z}}\varphi =0,`$ (3.53) $`F_{12}=e^2(\|\varphi \|^2r_0).`$ (3.54) The bosonic part of the action is $`\frac{1}{2\pi }\mathrm{d}^2x(\frac{1}{2e^2}\|_\mu \sigma \|^2+\|\sigma \varphi \|^2)`$ plus $`{\displaystyle \frac{1}{2\pi }}{\displaystyle \mathrm{d}^2x\left(\|D_\mu \varphi \|^2+\frac{1}{2e^2}(F_{12}^2+D^2)+i\theta F_{12}\right)}`$ (3.55) $`={\displaystyle \frac{1}{2\pi }}{\displaystyle \mathrm{d}^2x\left(\|2D_{\overline{z}}\varphi \|^2F_{12}\|\varphi \|^2+\frac{1}{2e^2}(F_{12}+D)^2\frac{1}{e^2}DF_{12}+i\theta F_{12}\right)}`$ $`={\displaystyle \frac{1}{2\pi }}{\displaystyle \mathrm{d}^2x\left(\|2D_{\overline{z}}\varphi \|^2+\frac{1}{2e^2}(F_{12}+D)^2\right)}{\displaystyle \frac{t_0}{2\pi }}{\displaystyle F_{12}\mathrm{d}^2x},`$ where $`D=e^2(\|\varphi \|^2r_0)`$ and $`t_0=r_0i\theta `$. For a given topological number $$k=\frac{1}{2\pi }F_{12}\mathrm{d}^2x,$$ (3.56) the real part of the action is bounded by $`kr_0`$, and the minimun is indeed attained by a solution to the equations (3.52)-(3.54). The value of the action for such an instanton is $$S_E=kt_0.$$ (3.57) Under the axial rotation by $`\mathrm{e}^{i\alpha }`$, the path-integral measure in this topological sector changes by the phase $`\mathrm{e}^{2ik\alpha }`$. Since the twisted superpotential has axial R-charge 2, we see that the relevant configurations are those with $`k=1`$. A solution to (3.53) and (3.54) is the vortex. For each vortex with $`k=1`$, $`\varphi `$ has a single simple zero. The moduli space of gauge equivalence classes of $`k=1`$ vortices is complex one-dimensional and is parametrized by the location of the zero of $`\varphi `$. To see this, we note the following well-known fact. The orbit of a solution to (3.53) under the complexified gauge transformations contains vortex solutions in one gauge equivalence class, and conversely, any gauge equivalence class of vortex solutions is contained in one such orbit. Here, a complexified gauge transformation is a rotation $`(iA_{\overline{z}},\varphi )(iA_{\overline{z}}+h_{\overline{z}}h^1,h\varphi )`$ by a function $`h`$ with values in $`𝐂^\times `$. Thus, we only have to find solutions to the equation $`D_{\overline{z}}\varphi =0`$ modulo the complexified gauge transformations. In other words, we only have to find pairs of a holomorphic line bundle with a holomorphic section. Here, it is convenient to compactify our Euclidean 2-plane to a Riemann sphere. Then, there is a unique holomorphic line bundle of $`k=1`$. Such a bundle has two-dimensional space of holomorphic sections, where each section has a single simple zero. The residual complexified gauge symmetry is a multiplication by a constant and it does not change the location of the zero of a section. Thus, the space of equivalence classes is complex one-dimensional and is parametrized by the zero locus of $`\varphi `$. Let us examine in more detail the behaviour of a vortex solution. We consider the vortex at $`z=0`$. For the finiteness of the action, $`\|\varphi \|`$ must approach the vacuum value $`\|\varphi \|^2=r_0`$ at infinity. By rescaling $`\varphi =\sqrt{r_0}\widehat{\varphi }`$ where $`\|\widehat{\varphi }\|1`$ at infinity, it becomes clear that the equations depend only on one length parameter $`1/e\sqrt{r_0}`$. This characterizes the size of the vortex. Thus, we expect that the gauge field is nearly flat on $`\|z\|1/e\sqrt{r_0}`$ and $`\widehat{\varphi }`$ is nearly covariantly constant there. Since $`(1/2\pi )F_{12}=1`$, we have $`A_\mu =_\mu \mathrm{arg}(z)`$ and $`\varphi =\sqrt{r_0}z/\|z\|`$ at infinity. The exact solution takes the form $`iA`$ $`=`$ $`{\displaystyle \frac{1f}{2}}\left({\displaystyle \frac{\mathrm{d}z}{z}}{\displaystyle \frac{\mathrm{d}\overline{z}}{\overline{z}}}\right),`$ (3.58) $`\varphi `$ $`=`$ $`\sqrt{r_0}\mathrm{exp}\left({\displaystyle _{\|z\|^2}^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}wf(w)}{2w}}\right){\displaystyle \frac{z}{\|z\|}},`$ (3.59) where $`f`$ is a function of $`w=\|z\|^2`$ which satisfies the equation $`wf^{\prime \prime }=\frac{e^2r_0}{2}f+ff^{}`$ and the boundary condition $`f(0)=1,f(+\mathrm{})=0`$. The asymptotic behaviour of the function $`f(\|z\|^2)`$ at $`\|z\|1/e\sqrt{r_0}`$ is $$f=\mathrm{𝑐𝑜𝑛𝑠𝑡}\sqrt{m\|z\|}\mathrm{e}^{m\|z\|}+\mathrm{},$$ (3.60) where const is a numerical constant and $$m=e\sqrt{2r_0}.$$ (3.61) The vortex at $`z=z_0`$ is obtained from the above solution simply by the replacement $`zzz_0`$. ### Fermionic Zero Modes Now let us examine the fermionic zero modes in the instanton background. Since $`\sigma =0`$ in this background, the fermionic part of the action (3.51) decomposes into two parts as $$i(\overline{\psi }_{},\lambda _+)\left(\begin{array}{cc}2D_{\overline{z}}& \varphi \\ \varphi ^{}& \frac{1}{e^2}_z\end{array}\right)\left(\genfrac{}{}{0pt}{}{\psi _{}}{\overline{\lambda }_+}\right)+i(\overline{\psi }_+,\lambda _{})\left(\begin{array}{cc}2D_z& \varphi \\ \varphi ^{}& \frac{1}{e^2}_{\overline{z}}\end{array}\right)\left(\genfrac{}{}{0pt}{}{\psi _+}{\overline{\lambda }_{}}\right)$$ (3.62) Using the index theorem of , one can see that the operators of the first and the second terms have index $`1`$ and $`1`$ respectively in our background. Furthermore, there is no normalizable zero modes for $`(\overline{\psi }_{},\lambda _+)`$ nor $`(\psi _+,\overline{\lambda }_{})`$. To see this we note that $`{\displaystyle \mathrm{d}^2z\left(\left\|2D_z\psi _+\varphi \overline{\lambda }_{}\right\|^2+2e^2\left\|\varphi ^{}\psi _+\frac{1}{e^2}_{\overline{z}}\overline{\lambda }_{}\right\|^2\right)}`$ (3.63) $`=`$ $`{\displaystyle \mathrm{d}^2z\left(\|2D_z\psi _+\|^2+2e^2\|\varphi \|^2\|\psi _+\|^2+\frac{2}{e^2}\|_{\overline{z}}\overline{\lambda }_{}\|^2+\|\varphi \|^2\|\overline{\lambda }_{}\|^2\right)}`$ (3.64) in the vortex background, where $`\psi _+`$ and $`\overline{\lambda }_{}`$ are considered here as commuting spinors. The left hand side vanishes for a zero mode of $`(\psi _+,\overline{\lambda }_{})`$, but then vanishing of the right hand side requires $`\psi _+=\overline{\lambda }_{}=0`$. The argument for $`(\overline{\psi }_{},\lambda _+)`$ is the same. Thus, each of $`(\psi _{},\overline{\lambda }_+)`$ and $`(\overline{\psi }_+,\lambda _{})`$ has exactly one zero mode. Actually, the expression for the zero mode in terms of the vortex solution is available; $$\left(\begin{array}{c}\psi _{}^{(0)}\\ \overline{\lambda }_{}^{(0)}\end{array}\right)=\left(\begin{array}{c}D_z\varphi \\ F_{12}\end{array}\right),\left(\begin{array}{c}\overline{\psi }_+^{(0)}\\ \lambda _+^{(0)}\end{array}\right)=\left(\begin{array}{c}D_{\overline{z}}\varphi ^{}\\ F_{12}\end{array}\right).$$ (3.65) It is easy to verify using the vortex equations (3.53) and (3.54) that these are indeed the zero modes. These come from the supersymmetry transformation under the broken supercharges $`\overline{Q}_{}`$ and $`Q_+`$. The asymptotic behaviour of $`\varphi ^{}\psi _{}^{(0)}`$ and $`\overline{\psi }_+^{(0)}\varphi `$ are $`\varphi ^{}\psi _{}^{(0)}`$ $`=`$ $`_z\|\varphi \|^2=\mathrm{𝑐𝑜𝑛𝑠𝑡}\times r_0{\displaystyle \frac{\overline{z}}{\|z\|}}\sqrt{{\displaystyle \frac{m}{\|z\|}}}\mathrm{e}^{m\|z\|}+\mathrm{}`$ (3.66) $`\overline{\psi }_+^{(0)}\varphi `$ $`=`$ $`_{\overline{z}}\|\varphi \|^2=\mathrm{𝑐𝑜𝑛𝑠𝑡}\times r_0{\displaystyle \frac{z}{\|z\|}}\sqrt{{\displaystyle \frac{m}{\|z\|}}}\mathrm{e}^{m\|z\|}+\mathrm{},`$ (3.67) for the vortex at $`z=0`$. These agree with the asymptotic behaviour of the propagators for a Dirac fermion of mass $`m=e\sqrt{2r_0}`$; $$\varphi ^{}\psi _{}^{(0)}r_0S_{}^F(z),\overline{\psi }_+^{(0)}\varphi r_0S_{++}^F(z).$$ (3.68) This is not an accident; After a suitable similarity transformation, the operators in (3.62) near infinity (where $`\varphi =\sqrt{r}z/\|z\|`$ and $`A_\mu =_\mu \mathrm{arg}(z)`$) become nothing but the Dirac operator for a free fermion of mass $`m=e\sqrt{2r}`$. ### The Computation In order to prove the generation of $`\mathrm{e}^Y`$ term in the twisted superpotential, we would like to compute a quantity that is non-vanishing only when such a term is generated. To find out what is the appropriate object to look at, let us return to our theory of $`\mathrm{\Sigma }`$ and $`Y`$. We recall that the Lagrangian is $$\stackrel{~}{L}=\mathrm{d}^4\theta (\frac{1}{2e^2}\overline{\mathrm{\Sigma }}\mathrm{\Sigma }\frac{1}{4r_0}\overline{Y}Y)+\frac{1}{2}(\mathrm{d}^2\stackrel{~}{\theta }(\mathrm{\Sigma }(Yt)+c\mu \mathrm{e}^Y)+c.c.),$$ (3.69) and we were asking whether $`c`$ is zero or not. Here we have approximated the Kahler potential for $`Y`$ by the one in the continuum limit (3.25). This theory involves a $`U(1)`$ gauge field. However, there is no charged field and the only appearance of the gauge field is in the kinetic term and in the Theta term (with the Theta angle being $`\vartheta \theta `$ where $`\vartheta =\mathrm{Im}(y)`$). As is well known , the effect of the gauge field is to generate a mass term for $`\vartheta `$: $$U=\frac{e^2}{2}\left(\stackrel{~}{\vartheta \theta }\right)^2,$$ (3.70) where $`(\stackrel{~}{\alpha })^2=\mathrm{min}\{(\alpha +2\pi n)^2\|n𝐙\}`$. Thus, we can treat $`\mathrm{\Sigma }`$ as an ordinary twisted chiral superfield which has a complex auxiliary field. In particular, the theory without $`\mathrm{e}^Y`$ term is a free theory of two twisted chiral multiplets. It is easy to diagonalize the $`\mathrm{\Sigma }Y`$ mixing and it turns out that the combinations $`X^{(\pm )}=\pm \mathrm{\Sigma }/(2e)+(Yt)/(2\sqrt{2r_0})`$ are superfields of mass $$\pm m=\pm e\sqrt{2r_0}.$$ (3.71) Now, it is easy to see that the fermionic components $`\chi _+`$ and $`\overline{\chi }_{}`$ of $`\overline{Y}`$ has vanishing two point function in the free theory; $$\chi _+(x)\overline{\chi }_{}(y)=0,\text{if }c=0.$$ (3.72) However, if $`\mathrm{e}^Y`$ is generated, the twisted F-term would contain a term $`\mathrm{e}^y\overline{\chi }_+\chi _{}`$ in the Lagrangian. This would contribute to the two point function as $$\chi _+(x)\overline{\chi }_{}(y)=c\mu r_0^2\mathrm{d}^2z\mathrm{e}^tS_{++}^F(xz)S_{}^F(zy)$$ (3.73) where $`S_{\alpha \beta }^F(xy)`$ is the Dirac propagator for the fermions of mass $`m=e\sqrt{2r_0}`$. Now, let us compute the two point function $`\chi _+(x)\overline{\chi }_{}(y)`$ in the original gauge theory. We recall from (3.21) that $$\chi _+=2\overline{\psi }_+\varphi ,\overline{\chi }_{}=2\varphi ^{}\psi _{}.$$ (3.74) Since the product of these carries an axial R-charge 2, only the vortex backgrounds with $`k=1`$ can contribute to the two point function. The contribution is expressed as an integration over the location $`z`$ of the vortex $$\chi _+(x)\overline{\chi }_{}(y)=4\mathrm{\Lambda }_{UV}\mathrm{d}^2z\mathrm{e}^{t_0}\left(\overline{\psi }_+^{(0)}\varphi _{(z)}\right)(x)\left(\varphi _{(z)}^{}\psi _{}^{(0)}\right)(y),$$ (3.75) where $`\varphi _{(z)}`$ is the vortex solution at $`z`$. The factor of $`\mathrm{\Lambda }_{UV}`$ comes from the measure of bosonic zero modes ($`\mathrm{\Lambda }_{UV}^2`$) and that for the fermionic zero modes ($`\mathrm{\Lambda }_{UV}^1`$). As we have seen, the fermionic zero mode multiplied by $`\varphi `$ is proportional to the Dirac propagator (3.68). Thus, we obtain $$\chi _+(x)\overline{\chi }_{}(y)=\mathrm{𝑐𝑜𝑛𝑠𝑡}\times r_0^2\mathrm{\Lambda }_{UV}\mathrm{d}^2z\mathrm{e}^{t_0}S_{++}^F(xz)S_{}^F(zy),$$ (3.76) which agrees with (3.73) considering the relation $`\mathrm{\Lambda }_{UV}\mathrm{e}^{t_0}=\mu \mathrm{e}^t`$. Thus, we have shown that our dual theory correctly reproduces the gauge theory result if and only if $`c0`$. #### 3.2.3 Solitons and Dualization In the $`R1/R`$ duality, as reviewed before, the momentum and winding modes get exchanged. This view provides us with another way to interpret the generation of superpotential (3.50). Let us turn off the gauge interaction and consider $`\mathrm{\Sigma }`$ as a non-dynamical parameter. Before dualization, we have a field $`\mathrm{\Phi }`$ of mass $`\mathrm{\Sigma }`$. In the dual description $`\mathrm{\Phi }`$ should arise as a winding mode. Indeed if we consider the superpotential $$W=\mathrm{\Sigma }(Yt)+e^Y.$$ (3.77) Viewing $`\mathrm{\Sigma }`$ as non-dynamical, the vacua are labeled by $`_YW=0`$ and we obtain $$_YW=0e^Y=\mathrm{\Sigma }$$ (3.78) and the critical points are given by $$Y_n=Y_0+2n\pi i$$ (3.79) where $`Y_0`$ is a special solution of (3.78). The vacua are indexed by an integer $`\|n`$, corresponding to winding number in the $`Y`$ plane. This should correspond to momentum mode in the $`\mathrm{\Phi }`$ variable. In other words the $`\mathrm{\Phi }`$ field should have $`Y`$-winding number charge $`+1`$ and acts on vacua $$\mathrm{\Phi }:\|n\|n+1$$ (3.80) In other words $`\mathrm{\Phi }`$ should be identified with the soliton interpolating between the vacua. The BPS mass in this sector is given by $$\|m\|=\frac{1}{2\pi }\|W(Y_n)W(Y_{n1})\|=\|\mathrm{\Sigma }\|,$$ (3.81) in agreement with the expected mass of $`\mathrm{\Phi }`$. This reasoning provides another view point on the superpotential in the actual gauge system. Related ideas were recently discussed in . ### 3.3 A Few Generalizations It is straightforward to extend the description using the dual fields to more general gauge theories. We consider here two generalizations; the case with $`U(1)^k`$ gauge group and the case with twisted masses. ### Many $`U(1)`$ Gauge Groups The first example is $`U(1)^k`$ gauge group with $`N`$ matter fields. We denote the field strength superfield for the $`a`$-th gauge group by $`\mathrm{\Sigma }_a`$ ($`a=1,\mathrm{},k`$), and the chiral superfield for the $`i`$-th charged matter by $`\mathrm{\Phi }_i`$ ($`i=1,\mathrm{},N`$). We denote by $`Q_{ia}`$ the charge of the $`i`$-th matter under the $`a`$-th gauge group. The exact twisted superpotential can be obtained by the localization argument which reduces the problem to the sum of copies of the single flavor case. This time, instead of keeping only one gauge coupling we keep $`k`$ of them. We start with the sum of $`N`$ copies of $`U(1)`$ theory with charge $`1`$ matter, and take the weak coupling limit except for the $`U(1)^k`$ gauge group embedded in $`U(1)^N`$ according to the charge matrix $`Q_{ia}`$. This constrains the $`N`$ gauge fields as $$\mathrm{\Sigma }_i=\underset{a=1}{\overset{k}{}}Q_{ia}\mathrm{\Sigma }_a.$$ (3.82) The exact superpotential for $`\mathrm{\Sigma }_a`$’s and the dual $`Y_i`$ of $`\mathrm{\Phi }_i`$ is then given by $$\stackrel{~}{W}=\underset{a=1}{\overset{k}{}}\mathrm{\Sigma }_a\left(\underset{i=1}{\overset{N}{}}Q_{ia}Y_it_a\right)+\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}.$$ (3.83) Integrating over $`Y_i`$’s, we obtain the following effective twisted superpotential for $`\mathrm{\Sigma }_a`$’s: $$\stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}(\mathrm{\Sigma }_a)=\underset{a=1}{\overset{k}{}}\mathrm{\Sigma }_a\left(\underset{i=1}{\overset{N}{}}Q_{ia}\left(\mathrm{log}\left(\underset{b=1}{\overset{k}{}}Q_{ib}\mathrm{\Sigma }_b/\mu \right)1\right)+t_a\right).$$ (3.84) This is the effective superpotential that we would obtain if we integrate out $`\mathrm{\Phi }_i`$’s in the original gauge theory . ### Twisted Masses The next example is $`U(1)^k`$ gauge theory with $`N`$ charged matter fields as above, but the twisted masses $`\stackrel{~}{m}_i`$ for the matter fields are turned on. The twisted mass can be considered as the lowest components of the field strength superfields of the $`U(1)^N/U(1)^k`$ flavor symmetry group. This is a non-trivial deformation of our gauge theory. The (anomalous) axial R-symmetry is explicitly broken by this perturbation but is restored if $`\stackrel{~}{m}_i`$ are rotated as $`\stackrel{~}{m}_i\mathrm{e}^{2i\alpha }\stackrel{~}{m}_i`$. As noted before, the dualization for this case follows from extending the $`U(1)^k`$ gauge group to $`U(1)^N`$ and taking the suitable decoupling limit. With the twisted masses, the constraint on the field strengths (3.82) is shifted as $$\mathrm{\Sigma }_i=\underset{a=1}{\overset{k}{}}Q_{ia}\mathrm{\Sigma }_a\stackrel{~}{m}_i.$$ (3.85) The exact twisted superpotential is thus $$\stackrel{~}{W}=\underset{a=1}{\overset{k}{}}\mathrm{\Sigma }_a\left(\underset{i=1}{\overset{N}{}}Q_{ia}Y_it_a\right)+\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}\underset{i=1}{\overset{N}{}}\stackrel{~}{m}_iY_i.$$ (3.86) Note that the net number of deformation parameters is $`(Nk)`$ from the flavor group $`U(1)^N/U(1)^k`$; $`\delta \stackrel{~}{m}_i=_{a=1}^kQ_{ia}c_a`$ is absorbed by a shift of the origin of $`\mathrm{\Sigma }_a`$’s. Integration over $`Y_i`$’s yields an effective superpotential for $`\mathrm{\Sigma }`$ that we would obtain if we integrate out $`\mathrm{\Phi }_i`$’s in the original gauge theory (as is done in for $`k=1`$ case with $`Q_i=\pm 1`$). ## 4 Sigma Models From Gauge Theories The gauge theory studied in the previous section reduces at low enough energies to the non-linear sigma model on a certain manifold. To see this, we examine the space of classical vacua of the theory. This can be read by looking at the potential energy for the scalar fields $`\sigma ,\varphi _i`$ $$U=\frac{e^2}{2}\left(\underset{i=1}{\overset{N}{}}Q_i\|\varphi _i\|^2r_0\right)^2+\underset{i=1}{\overset{N}{}}Q_i^2\|\sigma \|^2\|\varphi _i\|^2.$$ (4.1) We notice two branches of solutions to the vacuum equation $`U=0`$; Higgs branch where $`\sigma =0`$ and $`_{i=1}^NQ_i\|\varphi _i\|^2=0`$, and Coulomb branch where $`\sigma `$ is free and all $`\varphi _i=0`$. If $`\pm _{i=1}^NQ_i>0`$, $`\pm r_0`$ is bound to be large and there is only a Higgs branch. If $`_{i=1}^NQ_i=0`$, there is again only a Higgs branch except $`r_00`$ where a Coulomb branch develops. The degrees of freedom transverse to the Higgs branch have masses of order $`e\sqrt{\|r_0\|}`$, as can be seen from the Lagrangian $`_i\|D\varphi _i\|^2+(1/2e^2)\|\mathrm{d}\sigma \|^2+U`$. Thus, for a generic value of the parameter $`r_0`$, the theory at energies much smaller than $`e\sqrt{\|r_0\|}`$ describes the non-linear sigma model on the Higgs branch. This in particular means that all aspects of the sigma model at finite energies (including the BPS soliton spectra for massive sigma models) can be seen by studying the corresponding gauge theory. To be more precise, the Higgs branch $`X`$ is the space of solutions to $$\underset{i=1}{\overset{N}{}}Q_i\|\varphi _i\|^2=r_0$$ (4.2) modulo $`U(1)`$ gauge transformations $`\varphi _i\mathrm{e}^{iQ_i\gamma }\varphi _i`$. This space has complex dimension $`N1`$ and inherits a structure of Kahler manifold from that of flat $`𝐂^N`$ of $`\varphi _i`$’s. It is a standard fact that this is equivalent as a complex manifold to the quotient of $`𝐂^N𝒫`$ by the $`𝐂^\times `$ action $`\varphi _i\lambda ^{Q_i}\varphi _i`$, where $`𝒫`$ is some subset of codimension $`1`$. In particular, complex coordinates of $`X`$ are represented in the sigma model by the lowest components of chiral superfields. The parameter $`t_0=r_0i\theta `$ is identified as the complexified Kahler class. The first Chern class of this space is proportional to $`\|_{i=1}^NQ_i\|`$ and the cut-off dependence (3.3) of $`r_0`$ corresponds to the renormalization of the sigma model metric . The sigma model limit is thus $$e\mathrm{\Lambda }$$ (4.3) for $`_{i=1}^NQ_i0`$. The space $`X=(𝐂^N𝒫)/𝐂^\times `$ has an algebraic torus $`(𝐂^\times )^{N1}`$ as a group of holomorphic automorphisms; the group $`(𝐂^\times )^N`$ acting on $`𝐂^N`$ in a standard way modulo the comlpexified gauge group $`𝐂^\times `$. There is an open subset $`\{\varphi _i0,i\}`$ of $`X`$ on which $`(𝐂^\times )^{N1}`$ acts freely and transitively. Such a space $`X`$ is called a toric variety, or toric manifold if it is smooth. As is clear from the equation (4.2), if $`Q_i`$ are all positive (or all negative), the manifold $`X`$ is compact but if there is a mixture of positive and negative $`Q_i`$’s $`X`$ is non-compact. For $`_{i=1}^NQ_i=0`$, the FI parametr $`r=r_0`$ does not run and $`t=ri\theta `$ is a free parameter of the theory. At $`r=0`$, $`X`$ contains a singular point $`\varphi _i=0`$ where the $`U(1)`$ gauge group is unbroken. A new flat direction (Coulomb branch) develops there and the sigma model on $`X`$ becomes singular. The actual singularity is determined by the vanishing of the quantum effective potential at large values of $`\sigma `$. Since there is an effective superpotential $`\stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}(\sigma )`$ (3.39) which is valid at large $`\sigma `$, this is equivalent to the condition that $`_\sigma \stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}(\sigma )=0`$ at large values of $`\sigma `$. Thus, singularity of the quantum theory is located at $$t=\underset{i=1}{\overset{N}{}}Q_i\mathrm{log}(Q_i).$$ (4.4) Note that the theory is singular only for a particular value of $`\theta `$ (either $`0`$ or $`\pi `$) and $`r`$. In particular, the theories with $`r0`$ and $`r0`$ are smoothly connected to each other , even though the space $`X`$ at $`r0`$ differs from $`X`$ at $`r0`$. One can also consider abelian gauge theory with gauge group $`U(1)^k`$ and $`N`$ matter fields with $`N\times k`$ charge matrix $`Q_{ia}`$. For a generic value of $`t_a`$ in a certain range, the vacuum manifold is a Kahler manifold $`X`$ of dimension $`Nk`$. $`X`$ as a complex manifold is of the type $`(𝐂^N𝒫)/(𝐂^\times )^k`$ where $`𝒫`$ is a union of certain planes in $`𝐂^N`$. The algebraic torus $`(𝐂^\times )^{Nk}=(𝐂^\times )^N/(𝐂^\times )^k`$ acts on $`X`$ in an obvious way as holomorphic automorphisms and $`X`$ contains an open subset on which $`(𝐂^\times )^{Nk}`$ acts freely and transitively. Thus, $`X`$ is a toric variety. (In fact, any normal toric variety is obtained this way). We refer the physics reader to for more precise definition of a toric variety and its general properties. See also for more detail. If $`_{i=1}^NQ_{ia}=0`$ for all $`a`$, the theory is parametrized by $`k`$ dimensionless parameters $`t_a=r_ai\theta _a`$. Otherwise, there is a single scale parameter and $`k1`$ dimensionless parameters. Unlike in the single $`U(1)`$ case, the theory can possibly become singular at some locus even if $`_{i=1}^NQ_{ia}0`$. Singular locus in the parameter space is determined by finding a flat direction in the $`\sigma _a`$ space; $`_{\sigma _a}\stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}=0`$ where $`\stackrel{~}{W}_{\mathrm{𝑒𝑓𝑓}}`$ is given in (3.84) (more precisely, we must consider all possible “mixed branches” and do the same computation in the reduced theory). This was studied in detail in . However, as in the single $`U(1)`$ case, two generic points in the parameter space can be smoothly connected to each other without meeting a singularity. The first Chern class $`c_1(X)`$ of $`X`$ is not necessarily positive semi-definite. In the present paper, we only consider $`X`$ with $`c_1(X)0`$. In such a case, the running of the FI parameter of the gauge theory matches the running of the sigma model coupling at one-loop level, and our gauge theory indeed describes the non-linear sigma model on $`X`$ at energies well below the gauge coupling constants. To be more precise, the precise relation of the parameters is possibly complicated when $`c_1(X)`$ is close to zero and the one-loop running is not dominant. This can also be understood by comparing the size of the moduli spaces of gauge theory instantons and instantons of the non-linear sigma model (as shown in in the case of projective hypersurfaces). It was also noted in that the formal parameters corresponding to irrelevant deformations are in complicated relation even when $`c_1`$ is large. The precise relation can be determined by finding the so called flat coordinates with the expansion point at infinity. That would lead to the natural coordinates used in the large volume expansion of the non-linear sigma model . (This in particular applies to the mirror QFTs that we will obtain later in this paper in finding the map between the parameters of the non-linear sigma model and the mirror.) If there is a negative component in $`c_1(X)`$, the sigma model is not asymptotic free and is not well-defined. In such a case, our gauge theory has little to do with the manifold $`X`$. If $`c_1(X)`$ is negative definite, it is infra-red free and the sigma model (defined as a cut-off theory) flows to a free theory of $`c/3=dimX`$. #### 4.0.1 Examples of Toric Varieties Let us consider some examples of toric varieties. If we consider a $`U(1)`$ gauge theory with $`N`$ fields with charges $`+1`$, this gives a linear sigma model realization of $`𝐂\mathrm{P}^{N1}`$. More generally, if the charges of the matter fields are positive but not necessarily equal, it gives a realization of weighted projective space, with weights determined by the charges. If we consider a $`U(1)`$ gauge theory with $`N`$ fields with charge $`+1`$ and one with charge $`d`$, this gives a realization of the total space of the $`𝒪(d)`$ line bundle over $`𝐂\mathrm{P}^{N1}`$. Hirzebruch surface $`F_a`$ ($`a=0,1,\mathrm{}`$) is a toric manifold of dimension two which is realized as the vacuum manifold of the $`U(1)\times U(1)`$ gauge theory with four chiral fields with charges $`(1,0)`$, $`(1,0)`$, $`(0,1)`$, and $`(a,1)`$. $`F_0`$ is the product $`𝐂\mathrm{P}^1\times 𝐂\mathrm{P}^1`$ and $`F_1`$ is the blow up of $`𝐂\mathrm{P}^2`$ at one point. The first Chern class is positive for these two cases and it is positive-semi-definite for $`F_2`$. In all these cases, the gauge theory describes the non-linear sigma model. For $`a3`$, however, $`c_1(F_a)`$ is not even positive semi-definite and the sigma model is not well-defined. #### 4.0.2 Holomorphic Vector Fields and Deformations of Linear Sigma Model The unitary subgroup $`U(1)^{Nk}`$ of the holomorphic automorphism group $`(𝐂^\times )^{Nk}`$ is actually an isometry group of the Kahler manifold $`X`$. Since it is an abelian group, as explained in section 2, we can use this to deform the sigma model on $`X`$ by a potential term like (2.26). One may be interested in whether the deformation can be realized in the gauge theory. In fact, this $`U(1)^{Nk}`$ is a commutative subgroup of the flavor symmetry group of the gauge theory, and we can consider the deformation by twisted masses. By construction, the potential deformation of the sigma model is naturally identified as the twisted mass deformation of the gauge group. One can also confirm this by computing the bosonic potential in gauge theory and by checking that it agrees with the potential $`\frac{1}{2}\|\stackrel{~}{m}_AV_A\|^2+\frac{1}{2}\|\overline{\stackrel{~}{m}}_AV_A\|^2`$ from the holomorphic isometry (where $`V_A`$ are the generators of $`U(1)^{Nk}`$). We exhibit this in the simplest case $`X=𝐂\mathrm{P}^1`$ and leave the general case as an exercise. The $`𝐂\mathrm{P}^1`$ sigma model is realized by a $`U(1)`$ gauge theory with two fields $`\mathrm{\Phi }_1`$ and $`\mathrm{\Phi }_2`$ of unit charge. The potential of the gauge theory with the twisted mass is given by $$U=\frac{e^2}{2}(\|\varphi _1\|^2+\|\varphi _2\|^2r)^2+\|\sigma \stackrel{~}{m}\|^2\|\varphi _1\|^2+\|\sigma \|^2\|\varphi _2\|^2.$$ (4.5) We now fix $`\varphi _1`$ and $`\varphi _2`$ to lie in the $`𝐂\mathrm{P}^1`$ before perturbation ($`\|\varphi _1\|^2+\|\varphi _2\|^2=r`$) and then extremize the potential with respect to $`\sigma `$. Plugging the result back into (4.5), we obtain the potential $$U_{\stackrel{~}{m}}=\frac{\|\stackrel{~}{m}\|^2}{r}\|\varphi _1\|^2\|\varphi _2\|^2.$$ (4.6) On the other hand, the $`U(1)`$ isometry generates a holomorphic vector field with $`V^z=iz/z`$ where $`z`$ is the coordinate of $`𝐂\mathrm{P}^1`$ given by $`z=\varphi _1/\varphi _2`$. The metric of $`𝐂\mathrm{P}^1`$ determined by the quotient is the standard Fubini-Study metric $`\mathrm{d}s^2=r\|\mathrm{d}z\|^2/(1+\|z\|^2)^2`$. Measuring $`V`$ by this, we obtain the potential $$\|\stackrel{~}{m}\|^2V^2=\|\stackrel{~}{m}\|^2\frac{r\|z\|^2}{(1+\|z\|^2)^2}=\|\stackrel{~}{m}\|^2\frac{\|\varphi _1\|^2\|\varphi _2\|^2}{r}$$ (4.7) which is nothing but (4.6). ### 4.1 Hypersurfaces and Complete Intersections The non-linear sigma models of hypersurface or complete intersections in a compact toric manifold can also be realized as a gauge theory . Let $`X`$ be the compact toric manifold realized as the vacuum manifold of $`U(1)^k`$ gauge theory with $`N`$ matter fields $`\mathrm{\Phi }_i`$ of charge matrix $`Q_{ia}`$. We shall consider the submanifold $`M`$ of $`X`$ defined by the equations $$G_\beta =0,\beta =1,\mathrm{},l,$$ (4.8) where $`G_\beta `$ are polynomials of $`\mathrm{\Phi }_i`$ of charge $`d_{\beta a}`$ for the $`a`$-th $`U(1)`$ gauge group. Let us add $`l`$ matter fields $`P_\beta `$ of charge matrix $`d_{\beta a}`$ to the $`U(1)^k`$ gauge theory. This theory by itself realizes a non-linear sigma model on a non-compact toric manifold $`V`$. $`V`$ is the total space of the sum of $`l`$ line bundles on $`X`$; The new coordinates $`p_\beta `$ parametrize the fibre directions and $`X`$ is embedded in $`V`$ as the zero section $`p_\beta =0`$. Now let us consider the gauge theory, with the same gauge group and the same matter content, but having a gauge invariant superpotential of the form $$W=\underset{\beta =1}{\overset{l}{}}P_\beta G_\beta (\mathrm{\Phi }).$$ (4.9) Then, the vacuum equation requires $`G_\beta =0`$ and $`_\beta p_\beta _iG_\beta =0`$. If $`M`$ is a smooth complete intersection in $`X`$, this means that $`p_\beta =0`$ for all $`\beta `$ and the vacuum manifold is $`M`$ itself. Thus, the gauge theory with the superpotential (4.9) realizes the non-linear sigma model on the complete intersection $`M`$ in $`X`$. Thus the sigma model on the compact manifold $`M`$ is closely related to the sigma model on the non-compact toric manifold $`V`$. In fact they both have the same gauge field and matter content. The only difference between them is that, in the compact theory, there is an F-term involving a gauge invariant superpotential $`W`$ which yields $`cc`$ ring deformation. In this sense the compact theory can be embedded in the non-compact theory. The non-compact theory is obtained by considering $`WϵW`$ in the limit of setting $`ϵ0`$. Of course this limit changes drastically the behavior of the theory and in particular the theory has $`2l`$ more complex dimensions in the UV. However, we can ask if there are any quantities which are unaffected by this deformation. The answer is that if we are considering quantities (such as $`ac`$ ring) which are sensitive only to the twisted F-terms such as Kahler parameters (the FI terms of the gauge theory), then they should not depend on the F-term deformations. Thus for those questions it should be irrelevant whether we are considering the compact theory or the non-compact theory. All that we have to do is to find out how the states and the operators of the compact theory are embedded in that of the non-compact theory and compute the protected quantities that way. This is in fact analogous to embedding questions of vacuum geometry from one theory in a theory of higher central charge, discussed in . It is also similar to the computation of the elliptic genus of minimal models using a free theory with higher central charge . In the present case this issue has been studied in where they also obtain the embedding of the chiral field of the sigma model on $`M`$ theory with that of $`V`$. This result is stated as follows. Let $`\delta _\beta `$ ($`\beta =1,\mathrm{},l`$) be the $`ac`$ ring element of the theory on $`V`$ defined by $$\delta _\beta =\underset{a=1}{\overset{k}{}}d_{\beta a}\mathrm{\Sigma }_a,$$ (4.10) where $`\mathrm{\Sigma }_a`$ is the field strength of the $`a`$-th gauge group. Then, the correlation functions $`\mathrm{}_M`$ of the A-twisted model on $`M`$ are obtained from those $`\mathrm{}_V`$ on $`V`$ by $$𝒪_1\mathrm{}𝒪_s_M=𝒪_1\mathrm{}𝒪_s(\delta _1^2)\mathrm{}(\delta _l^2)_V.$$ (4.11) An intuitive understanding of this relation is available in the quantum mechanics obtained by dimensional reduction. Let $`L_a`$ be a line bundle over $`V`$ defined by the equivalence relation $`(p_\beta ,\varphi _i;c)(_a\lambda _a^{d_{\beta a}}p_\beta ,_a\lambda _a^{Q_{ia}}\varphi _i;\lambda _ac)`$ where $`(\lambda _a)`$ is the element of the complexified gauge group $`(𝐂^\times )^k`$ and $`c`$ is the fibre coordinate. Then, one can show that the $`ac`$ ring element $`\mathrm{\Sigma }_a`$ corresponds to the first Chern class $`c_1(L_a)`$ of $`L_a`$. Now, let us consider the tensor product $$L_\beta =\underset{a=1}{\overset{k}{}}L_a^{d_{\beta a}},$$ (4.12) whose first Chern class corresponds to the $`ac`$ ring element $`\delta _\beta `$. The line bundle $`L_\beta ^1`$ has a section proportional to $`p_\beta `$ and thus $`\delta _\beta `$ represents the divisor class of $`p_\beta =0`$ in $`V`$. In particular, the product $`(\delta _1)\mathrm{}(\delta _l)`$ represents the class $`X`$ in $`V`$. On the other hand, the bundle $`L_\beta `$ has a section proportional to $`G_\beta `$ and thus $`\delta _\beta `$ represents the divisor class of $`G_\beta =0`$. Therefore, $`(\delta _1^2)\mathrm{}(\delta _l^2)`$ represents a delta function supported on $`M`$. What is shown in is basically that this quantum mechanical interpretation remains true for the full 2d QFT as long as topological correlators are concerned. It turns out that we need a stronger version of the relation (4.11). We thus proceed to a physical derivation of this relation, which uses the fact that the compact and non-compact theories are embedded in the same underlying physical theory, and yields the stronger version that we need. The basic idea is that with $`ϵ0`$ we have turned on a superpotential and we can follow the states from the non-compact theory to the compact theory. In this sense the non-compact theory flows in the IR limit (for non-vanishing $`ϵ`$) to the compact theory. We can thus follow the states in the non-compact theory and ask which ones survive in the IR limit. By the nature of the RG flow, we will be losing some states as we take the IR limit. However, if we concentrate on the ground states in the Ramond sector which correspond to normalizable states in the non-compact theory, then in the compact theory they are bound to survive as a ground state (the same cannot be said of the non-normalizable ground states in the non-compact theory, which might disappear from the spectrum of the compact theory). In the sigma model on the non-compact manifold $`V`$, the ground state corresponding to the operator $`\delta =\delta _1\mathrm{}\delta _l`$ is a normalizable state which is the product of Kahler forms that control the sizes of the compact part of the geometry. We denote this state by $$\|\delta _V.$$ (4.13) In the large volume limit, this state has the axial R-charge $`dim_𝐂V+2l=dim_𝐂M`$ which is the lowest among normalizable ground states. Now, let us turn on the superpotential $`ϵW`$ where $`W`$ is the one given in (4.9). The state $`\|\delta _V`$ which is the unique normalizable ground state in the Ramond sector with axial charge $`dim_𝐂V+2l`$ will be deformed to a state which we denote by $`\|\delta _ϵ`$. By standard supersymmetry arguments, this state is a normalizable ground state of the Ramond sector. For large $`ϵ`$ (or for any finite $`ϵ`$ in the IR limit) the theory corresponds to the sigma model on the compact manifold $`M`$ for which there is a unique ground state with axial charge $`dim_𝐂M`$, which is also sometimes denoted by $`\|1_M`$. Thus we can identify $`\|\delta _ϵ=\|1_M`$ as long as axial charge is conserved. Thus, we have the following correspondence of states as we increase $`ϵ`$; $$\|\delta _V\stackrel{ϵ0}{}\|\delta _ϵ\stackrel{ϵ\mathrm{}}{}\|1_M.$$ (4.14) Here we have used the axial R-charge to identify the state to which $`\|\delta _V`$ is deformed. For the more general case, where the axial R-sysmmetry is broken by an anomaly, the result still remains true, as we will now argue. To see this, note that at large Kahler parameters, i.e. as $`t_a\mathrm{}`$ axial R-charge is a good symmetry. So at least to leading order it is correct. To show it is true for all $`t_a`$ we proceed as follows: According to the topologically twisted theory picks a section of the vacuum bundle which is holomorphic in the sense defined by the topological twisting. In particular we can choose the pertubed state $`\|\delta _ϵ`$ so that $$\frac{}{\overline{t}_a}\|\delta _ϵ=0.$$ (4.15) For infinitely large $`t_a`$, since the axial charge is conserved the state $`\|\delta _V`$ is deformed for finite $`ϵ`$ to $`\|1_M`$. For finite but large $`t_a`$ the fact that we can choose a holomorphic section of the vacuum bundle shows that the difference between $`\|\delta _ϵ`$ and $`\|1_M`$ can only be given by states with coefficients involving some powers of $`q_a=e^{t_a}`$. Here we note that the axial R-charge can be made conserved by shifting the Theta angle to cancell (3.4). This in particular means that we assign non-negative R-charges to $`q_a`$ as long as the first Chern class of $`M`$ is positive-semi-definite. It thus follows that $`q_a`$’s would be accompanied by states which will have too small an R-charge to correspond to a normalizable state, since $`\|\delta _V`$ is the unique normalizable ground state with minimum R-charge. This establishes the relation (4.14) in general. The relation (4.14) can be used to yield another derivation of the relation (4.11): The topological correlations for sigma model on $`M`$ can be written as $$𝒪_1\mathrm{}𝒪_s_M=1\|𝒪_1\mathrm{}𝒪_s\|1_M=\delta \|𝒪_1\mathrm{}𝒪_s\|\delta _V=(\delta _1^2)\mathrm{}(\delta _l^2)𝒪_1\mathrm{}𝒪_s_V.$$ (4.16) where we have used the $`ϵ`$ independence in the topologial A-model computations, as $`ϵW`$ does not affect the $`ac`$ correlation functions (note that the overall proportionality factor is not fixed, and depends on the choice of normalization of the topologically twisted theory). #### 4.1.1 Examples of Hypersurfaces Consider a $`U(1)`$ gauge theory with $`N`$ fields $`\mathrm{\Phi }_i`$ of charge $`+1`$ and one field $`P`$ with charge $`d`$. So far this is the same as the linear sigma model description of the total space of the line bundle $`𝒪(d)`$ over $`𝐂\mathrm{P}^{N1}`$. However, now consider adding to the action the F-term with superpotential $$W=PG(\mathrm{\Phi }_i)$$ (4.17) where $`G`$ is a polynomial of degree $`d`$ in $`\mathrm{\Phi }_i`$. This gauge theory describes, in the infrared, the sigma model on a hypersurface of degree $`d`$ in $`𝐂\mathrm{P}^{N1}`$. The hypersurface for $`d=N`$ is a Calabi-Yau manifold of complex dimension $`N2`$. For $`d<n`$ it is a manifold of positive $`c_1`$ and the sigma model is asymptotically free. For $`d>n`$ it is a manifold with negative $`c_1`$ and the sigma model is not asymptotically free. The geometric regime corresponds to when the FI parameter is very large $`r0`$ (which is the UV limit for $`d<N`$). If we consider the limit $`r\mathrm{}`$ (which is the IR limit for $`d<N`$), we end up with an LG theory : $`P`$ picks up an expectation value and breaks the $`U(1)`$ gauge symmetry to a discrete subgroup $`𝐙_d`$ and the effective theory is given by the LG model with $`W=G(\mathrm{\Phi }_i)`$ divided by a $`𝐙_d`$ action, which corresponds to multiplying the fields $`\mathrm{\Phi }_i`$ by a $`d`$-th root of unity. If we turn off the superpotential, i.e., back in the non-compact toric model corresponding to the total space of the line bundle $`𝒪(d)`$ over $`𝐂\mathrm{P}^{N1}`$ the limit $`r\mathrm{}`$ would correspond to the orbifold of the free theory of $`N`$ fields $`\mathrm{\Phi }_i`$ modded out by $`𝐙_d`$, acting as the $`d`$-th root of unity on $`\mathrm{\Phi }_i`$. The generalization of the above to the case of hypersurfaces of weighted projective space are straightforward. ## 5 A Proof Of Mirror Symmetry In this section we show how the results of section 3 on the dynamics of supersymmetric gauge theories naturally lead to a proof of mirror symmetry. Before embarking on the proof, we discuss what we mean by the “proof”. ### 5.1 What We Mean by Proof of Mirror Symmetry As we have discussed before the action for a $`(2,2)`$ supersymmetric field theory has three types of terms: D-terms, F-terms and twisted F-terms. In the case of a supersymmetric sigma model, these three types of terms have the following interpretations: The D-terms correspond to changing the metric on the target manifold without changing the Kahler or complex parameters. For example for $`𝐂\mathrm{P}^1`$ we can consider a fixed total area, but deform the metric from the round metric to any other metric with the same total area by varying the D-term. Variations of F-terms and twisted F-terms correspond to deformations of complex and Kahler parameters of the target space. As discussed before, F-terms and twisted F-terms control many important aspects of the $`(2,2)`$ theories. In particular the $`cc`$ and $`ac`$ rings depend only on F-terms and twisted F-terms respectively and are independent of D-terms. These in particular encode the instanton corrections of the sigma model to the $`ac`$ ring. Also BPS structure of the massive $`(2,2)`$ theories are completely determined by the F-terms and twisted F-terms . One can define the notion of D-branes for $`(2,2)`$ theories (as is familiar in the conformal case, and can be easily extended to the non-conformal case as we will discuss in section 6). One can also show that the overlap of the corresponding boundary states with vacua will only depend on the F-terms (or twisted F-terms depending on which combination of supercharge the D-brane preserves). These indicate the importance of F-terms and twisted F-terms for the $`(2,2)`$ theories. In fact they become even more prominent in the conformal case, for example for the case of Calabi-Yau manifolds. Namely in that case the sigma model flows in the infrared limit to a conformal field theory which is determined entirely by the complex and Kahler parameters of the manifold. In other words the metric on the Calabi-Yau manifold adjusts itself to the unique form consistent with conformal invariance for a given Kahler and complex structure. In particular the D-terms are entirely determined by the F-terms and twisted F-terms in this case, as far as the IR behavior is concerned. As a consequence, a deformation of the conformal field theory corresponds to a deformation of the Kahler or complex structure of the Calabi-Yau manifold, which in turn is realized through variations of F-terms and twisted F-terms. From these facts, and from the fact that the $`cc`$ and $`ac`$ rings are determined in terms of F-terms and twisted F-terms one wonders whether the opposite is true (i.e. whether the $`cc`$ and $`ac`$ rings determine the F-terms and twisted F-terms and thus the full theory). This is, however, not completely true; theories with the same chiral ring can differ in the “integral structure” . More precisely, in order to completely specify the theory we need to know the structure of the allowed D-branes. We prove mirror symmetry in two different senses: In the strong sense, we find a dual theory, for which we prove it is equivalent to the original theory, up to deformations involving D-terms. In the weak sense, we propose a dual theory for which we can only show that the $`cc`$, $`ac`$, BPS structure of solitons and the $`D`$-brane structure are the same. The weak and strong senses become equivalent if we assume that the $`cc`$ and $`ac`$ ring, together with the integral structure determine the theory up to D-term variations. We are now ready to discuss the proof for mirror symmetry. We divide the discussion into several cases and indicate in each case what is the sense of the proof (strong/weak) that we shall present. ### 5.2 The Proof The cases of interest naturally divide into three different classes: (i) $`X=𝐂\mathrm{P}^{N1}`$ or a more general compact toric manifold with $`c_1(X)0`$, (ii) $`X`$ is a non-compact Calabi-Yau or a more general non-compact toric manifold. And finally, (iii) Sigma models on hypersurfaces or complete intersection in $`X`$. It turns out that the mirror for the cases (i) and (ii) are rather straightforward to derive, using the tools we have developed so far and we shall prove them in the strong sense. For the derivation of the case (iii) we need an additional tool, which we will discuss in section 6, and we postpone a complete discussion of (iii) to section 7. This will lead to a proof of case (iii) in the weak sense. In this section after deriving the general mirror for the cases (i) and (ii) we briefly mention how it effectively gives an answer of the case (iii) (to be more fully developed in section 7). We present some examples for cases (i) and (ii) and how it relates to case (iii) to illustrate the meaning of the results. We first consider the sigma model on a toric manifold which is described by a gauge theory with a single $`U(1)`$ gauge group with chiral fields of charge $`Q_i`$ ($`i=1,\mathrm{},N`$). We recall that the dual of the gauge theory is described by a vector multiplet with field strength $`\mathrm{\Sigma }`$ and $`N`$ periodic variable $`Y_i`$ dual to the charged matter fields. It has an approximate Kahler potential $`(1/2e^2)\|\mathrm{\Sigma }\|^2+_i\|Y_i\|^2/4r_0`$ and an exact twisted superpotential $$\stackrel{~}{W}=\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Q_iY_it(\mu )\right)+\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}.$$ (5.1) The fields $`\mathrm{\Sigma }`$ and $`_iQ_iY_i`$ have mass of order $`e\sqrt{\|r_0\|}`$ while the mass scale for the modes tangent to $`_iQ_iY_i=t`$ is of order $`\mu \sqrt{\|r_0\|}`$ where $`\mu \mathrm{\Lambda }`$ for $`_iQ_i0`$. In the sigma model limit $`e\mu `$, these mass scales are well-separated and it becomes appropriate to integrate out $`\mathrm{\Sigma }`$. This yields the constraint $$\underset{i=1}{\overset{N}{}}Q_iY_i=t.$$ (5.2) Clearly this can be solved by $`N1`$ periodic variables. Now, the dual theory becomes the theory of such $`N1`$ variables solving (5.2) with the twisted superpotential $$\stackrel{~}{W}=\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}.$$ (5.3) This can be interpreted as a sigma model on $`(𝐂^\times )^{N1}`$ with a twisted superpotential. Since the complex coordinates are the lowest components of twisted chiral superfields, this theory can be identified as the mirror of the non-linear sigma model on $`X`$. It is straightforward to extend this to the more general case where we start with $`k`$ $`U(1)`$ gauge groups and $`N`$ matter fields of charge $`Q_{ia}`$ ($`i=1,\mathrm{},N`$, $`a=1,\mathrm{},k`$). The target space $`X`$ of the sigma model is the quotient of $$\underset{i=1}{\overset{N}{}}Q_{ia}\|\varphi _i\|^2=r_a$$ (5.4) by the $`U(1)^k`$ action $`\varphi _i\mathrm{e}^{iQ_{ia}\gamma _a}\varphi _i`$, which is a toric variety of dimension $`Nk`$. The dual description of the gauge theory is obtained in (3.83). Integrating out the vector multiplet, we obtain the algebraic torus $`(𝐂^\times )^{Nk}`$ as the solutions to $$\underset{i=1}{\overset{N}{}}Q_{ia}Y_i=t_a.$$ (5.5) The dual theory is a sigma model on this $`(𝐂^\times )^{Nk}`$ with the twisted superpotential $$\stackrel{~}{W}=\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}.$$ (5.6) This can be identified as the mirror of the sigma model on $`X`$. We thus have established the mirror symmetry of the non-linear sigma model on a toric variety $`X^{Nk}`$ and the theory on the algebraic torus $`(𝐂^\times )^{Nk}`$ with a superpotential. We can also consider deforming the sigma model using holomorphic vector fields, as discussed in section 2. There are $`(Nk)`$ such parameters corresponding to the $`U(1)^N/U(1)^k`$ holomorphic isometry group for the above toric variety. We parameterize these deformations by $`N`$ parameters $`\stackrel{~}{m}_i`$ (as we will see below $`k`$ of them are redundant). The corresponding gauge theory is the one with the twisted masses $`\stackrel{~}{m}_i`$ and the twisted superpotential for the dual theory is obtained in (3.86). Adding the twisted masses does not affect the elimination of the field strength $`\mathrm{\Sigma }_a`$ and we obtain the same constraints (5.5). All it does is to shift the superpotential (5.6) as $$\stackrel{~}{W}=\mu \underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}\underset{i=1}{\overset{N}{}}\stackrel{~}{m}_iY_i.$$ (5.7) This is not a single valued function on $`(𝐂^\times )^{Nk}`$, but the multivaluedness is a constant shift and therefore the Lagrangian itself is well-defined. It is easy to see, using (5.5) that $`k`$ of the parameters $`\stackrel{~}{m}_i`$ are redundant and do not affect the superpotential. We will now illustrate these ideas in two concrete cases: One is in the context of compact toric varieties which we exemplify by using the case of $`𝐂\mathrm{P}^{N1}`$. The second one is for the case of non-compact toric varieties which we exemplify by considering the total space of $`𝒪(d)`$ line bundle over $`𝐂\mathrm{P}^{N1}`$. In the context of the latter example we also explain briefly how the mirror for the hypersurface case arises. We complete the discussion for the hypersurface case in section 7. ### 5.3 Compact Toric Manifold ### The $`𝐂\mathrm{P}^{N1}`$ Model The linear sigma model for $`X=𝐂\mathrm{P}^{N1}`$ is the $`U(1)`$ gauge theory with $`N`$ matter fields of charge $`1`$. The constraint $`_{i=1}^NY_i=t`$ is solved by $`Y_i=t/N\mathrm{\Theta }_i`$ ($`i=1,\mathrm{},N1`$) and $`Y_N=t/N+_{i=1}^{N1}\mathrm{\Theta }_i`$, where $`\mathrm{\Theta }_i`$ are periodic variables of period $`2\pi i`$ and can be considered as coordinates of $`(𝐂^\times )^{N1}`$. The superpotential is $$\stackrel{~}{W}=\mathrm{\Lambda }\left(\mathrm{e}^{\mathrm{\Theta }_1}+\mathrm{}+\mathrm{e}^{\mathrm{\Theta }_{N1}}+\underset{i=1}{\overset{N1}{}}\mathrm{e}^{\mathrm{\Theta }_i}\right),$$ (5.8) where $`\mathrm{\Lambda }=\mu \mathrm{e}^{t/N}`$ is the dynamical scale of the theory. This is the superpotential for supersymmetric affine $`A_{N1}`$ Toda field theory. Thus, we have derived the mirror symmetry of $`𝐂\mathrm{P}^{N1}`$ model and affine Toda theory which was observed in from various points of view. Having no F-term, the theory is invariant under $`U(1)_V`$ R-symmetry. The twisted superpotential (5.8) explicitly breaks $`U(1)_A`$ but its $`𝐙_{2N}`$ subgroup remains unbroken; for $`\mathrm{\Theta }_j\mathrm{\Theta }_j+2\pi i/N`$, $`\stackrel{~}{W}\mathrm{e}^{2\pi i/N}\stackrel{~}{W}`$. The vacua of the theory are given by the critical points of $`\stackrel{~}{W}`$, $`_{\mathrm{\Theta }_i}\stackrel{~}{W}=0`$. It is solved by $`\mathrm{e}^{\mathrm{\Theta }_1}=\mathrm{}=\mathrm{e}^{\mathrm{\Theta }_{N1}}=:X`$ where $`X=X^{(N1)}`$. Namely, there are $`N`$ vacua at $`\mathrm{e}^{\mathrm{\Theta }_j}=\mathrm{e}^{2\pi i\mathrm{}/N}`$ (all $`j`$) with the critical value $$\stackrel{~}{W}=N\mathrm{\Lambda }\mathrm{e}^{2\pi i\mathrm{}/N},$$ (5.9) ($`\mathrm{}=0,\mathrm{},N1`$). Each vacuum is massive and breaks spontaneously the axial R-symmetry $`𝐙_{2N}`$ to $`𝐙_2`$. All these are indeed the properties which are possessed by the $`𝐂\mathrm{P}^{N1}`$ model. An advantage of LG type description is that the BPS soliton spectrum can be exactly analyzed using the superpotential. A BPS soliton corresponds to a trajectory connecting two vacua which projects onto a straight line in the $`\stackrel{~}{W}`$ space. The mass is equal to the absolute value of the susy central charge which is given by the difference of the two critical values of $`\stackrel{~}{W}`$. From (5.9) we see that the BPS soliton connecting the $`0`$-th and the $`\mathrm{}`$-th vacua has central charge $$\stackrel{~}{Z}_0\mathrm{}=N\mathrm{\Lambda }(\mathrm{e}^{2\pi i\mathrm{}/N}1).$$ (5.10) One can also see that there are $`\left(\genfrac{}{}{0pt}{}{N}{\mathrm{}}\right)`$ such solitons . For each of them, $`\mathrm{}`$ of $`\mathrm{e}^{Y_i}`$ are equal to a trajectory $`f_{\mathrm{}}`$ in $`𝐂^\times `$ while the remaining $`(N\mathrm{})`$ of them to another $`f_N\mathrm{}`$. They both starts from $`\mathrm{e}^{t/N}`$ and ends at $`\mathrm{e}^{t/N+2\pi i\mathrm{}/N}`$ but $`f_{\mathrm{}}`$ has a relative winding number $`(1)`$ compared to $`f_N\mathrm{}`$. Since the winding in $`Y_i`$ is dual to the charge for phase rotation of $`\mathrm{\Phi }_i`$, the soliton has the same quantum number as the product $`\mathrm{\Phi }_{i_1}\mathrm{}\mathrm{\Phi }_i_{\mathrm{}}`$. Indeed, the soliton spectrum of the $`𝐂\mathrm{P}^{N1}`$ model has been studied in and it was found that the $`\mathrm{}=1`$ solitons are the elementary electrons $`\mathrm{\Phi }_i`$ and the higher-$`\mathrm{}`$ solitons are their $`\mathrm{}`$-th antisymmetric products. This spectrum for the solitons was also recovered from the $`tt^{}`$ geometry in . We can also consider deforming the $`𝐂\mathrm{P}^{N1}`$ sigma model by the addition of a combination of the holomorphic vector fields $`U(1)^{N1}`$. In such a case we obtain the mirror (5.7): $$\stackrel{~}{W}=\mathrm{\Lambda }\left(\mathrm{e}^{\mathrm{\Theta }_1}+\mathrm{}+\mathrm{e}^{\mathrm{\Theta }_{N1}}+\underset{i=1}{\overset{N1}{}}\mathrm{e}^{\mathrm{\Theta }_i}\right)+\underset{i=1}{\overset{N1}{}}(\stackrel{~}{m}_i\stackrel{~}{m}_N)\mathrm{\Theta }_i,$$ (5.11) As has been noted, $`\stackrel{~}{m}_i`$’s changes the central charge of the supersymmetry algebra by a term proportional to the charges $`S_i`$ of the global abelian symmetry $`U(1)^{N1}`$, $$\delta \stackrel{~}{Z}=\underset{i=1}{\overset{N1}{}}(\stackrel{~}{m}_i\stackrel{~}{m}_N)S_i.$$ (5.12) The charges $`S_i`$ can be identified with the weights of the $`SU(N)`$ global symmetry. In this way we can recover not only the soliton spectrum, but also their quantum numbers under the $`SU(N)`$ global symmetry, and obtain the anticipated result noted above . Note also, that the above deformation, deforms the quantum cohomology ring. If we denote by $`x=d\stackrel{~}{W}/dt`$ the generator of the chiral ring (corresponding to the Kahler class of $`𝐂\mathrm{P}^{N1}`$), from the above superpotential one obtains $$\underset{i=1}{\overset{N}{}}(x\stackrel{~}{m}_i)=\mathrm{\Lambda }^N.$$ (5.13) Note that the twisted masses deforms the cohomologry ring from the simple form $`x^N=\mathrm{\Lambda }^N`$ to an arbitrary polynomial of degree $`N`$. The result for the case of $`𝐂\mathrm{P}^1`$ was first derived through other arguments in . The general case was conjectured in from brane construction of the theory. As noted before, we do not attempt to specify the D-terms which are subject to perturbative corrections in general. However it is expected that the sine-Gordon theory ($`A_1`$ affine Toda theory) is integrable and the D-term is protected from quantum correction . Thus, one may expect that a more detailed statement of the mirror symmetry can be made. The $`𝐂\mathrm{P}^1`$ model realized as the linear sigma model possesses the $`SU(2)`$ isometry group as the global symmetry. Recall that we have seen that the Kahler potential for $`Y_i`$’s is approximately $`\|Y_i\|^2/4r_0`$ at the classical level. Recall also that $`r_0\mathrm{}`$ in the continuum limit. This suggests that the equivalence of the $`SU(2)`$ invariant $`𝐂\mathrm{P}^1`$ model and the sine-Gordon theory holds only in the limit where the Kahler potential of the latter vanishes. This is actually consistent with the observation that the $`N=2`$ sine-Gordon theory possesses $`SU(2)`$ global symmetry, rather than its q-deformation, in the limit of vanishing Kahler potential. Also, it was observed in that the scattering matrix of the BPS solitons of the $`𝐂\mathrm{P}^1`$ model and the sine-Gordon theory agree with each other in such a limit. It would be interesting to investigate the integrability and the protection of the D-term in more general cases. ### More General Cases We now study a general aspects of the mirror theory for more general toric manifold $`X`$. Let us consider the equations $`_{i=1}^Nv_iQ_{ia}=0`$ for integers $`v_i`$. The space of solutions form a lattice of rank $`Nk`$. Let $`v_i^\mathrm{A}`$ be the integral basis of this lattice; $$\underset{i=1}{\overset{N}{}}v_i^\mathrm{A}Q_{ia}=0,\{\begin{array}{c}\mathrm{A}=1,\mathrm{},Nk,\hfill \\ a=1,\mathrm{},k.\hfill \end{array}$$ (5.14) The constraints (5.5) on $`Y_i`$ can be solved by $`Nk`$ periodic variables $`\mathrm{\Theta }_\mathrm{A}`$ as $$Y_i=\underset{\mathrm{A}=1}{\overset{Nk}{}}v_i^\mathrm{A}\mathrm{\Theta }_\mathrm{A}+t_i$$ (5.15) where $`(t_i)`$ is a solution to $`_{i=1}^NQ_{ia}t_i=t_a`$ (an arbitrary choice will do; another choice is related by a shift of $`\mathrm{\Theta }_\mathrm{A}`$’s). Now, the superpotential (5.6) of the mirror theory can be expressed as $$\stackrel{~}{W}=\underset{i=1}{\overset{N}{}}\mathrm{exp}\left(t_i\mathrm{\Theta },v_i\right),$$ (5.16) where $`\mathrm{\Theta },v_i`$ is the short hand notation for $`_{\mathrm{A}=1}^{Nk}v_i^\mathrm{A}\mathrm{\Theta }_\mathrm{A}`$. Note that the expression (5.16) is the same as the function in which determines the quantum cohomology of toric manifolds. Namely, we have derived the result of as a straightforward consequence of our dual description. It is useful to consider $`v_i=(v_i^\mathrm{A})`$ as (generically linearly dependent) $`N`$ vectors in the lattice $`\mathrm{N}=𝐙^{Nk}`$ and $`\mathrm{\Theta }_\mathrm{A}`$ as the coordinates on $`\mathrm{M}_𝐂=\mathrm{M}𝐂`$ where $`\mathrm{M}`$ is the dual lattice of $`\mathrm{N}`$. Then, $`\mathrm{\Theta },v_i`$ that appears in (5.16) can be identified as the natural pairing. It is well-known in toric geometry that $`X`$ is compact if and only if the cone generated by $`v_i`$ covers the whole $`\mathrm{N}_𝐑=\mathrm{N}𝐑`$ (i.e. any element of $`\mathrm{N}_𝐑`$ is expressed as a linear combination of $`v_i`$’s with non-negative coefficients). This in particular means that there is no value of $`\mathrm{\Theta }\mathrm{M}_𝐂`$ such that $`\mathrm{Re}\mathrm{\Theta },v_i0`$ for all $`i`$. In other words, for any non-zero $`\mathrm{Re}\mathrm{\Theta }`$, $`\mathrm{Re}\mathrm{\Theta },v_i<0`$ for some $`i`$. Thus, there is no obvious run-away direction of the superpotential (5.16) and we generically expect a discrete spectrum. <sup>1</sup><sup>1</sup>1There can be an accidental situation (which appears in a sublocus of the parameter space) where there is a run-away direction. This comes from the singularity of the sigma model that is associated with the development of the Coulomb or mixed branch, which is studied in . This is of course an expected property for a mirror of a compact sigma model. Since the Witten index of the sigma model is equal to the Euler number $$\mathrm{Tr}(1)^F=\chi (X),$$ (5.17) the number of critical points of the superpotential (5.16) must agree with $`\chi (X)`$. To check this, we note another well-known fact in toric geometry: Our algebraic torus $`(𝐂^\times )^{Nk}`$ (with coordinates $`\mathrm{e}^{\mathrm{\Theta }_\mathrm{A}}`$) is embedded as an open subset of a “dual” toric variety $`Y`$ and each term in (5.16) extends to a section of the anti-canonical bundle $`K_Y^1`$ of $`Y`$.<sup>2</sup><sup>2</sup>2The “dual” toric variety is constructed as follows. Take the convex hull of $`v_i`$’s in $`\mathrm{N}_𝐑`$; this makes a convex polytope $`\mathrm{\Delta }`$ in $`\mathrm{N}_𝐑`$. Consider the dual $`\mathrm{\Delta }^{}\mathrm{M}_𝐑`$ of $`\mathrm{\Delta }\mathrm{N}_𝐑`$ defined as the set of points in $`\mathrm{M}_𝐑`$ whose values at $`\mathrm{\Delta }`$ are $`1`$. Then, the vertices $`\{u_I\}`$ of $`\mathrm{\Delta }^{}`$ determines the “dual” toric variety $`Y`$. Thus, each partial derivative $`_\mathrm{A}\stackrel{~}{W}=\stackrel{~}{W}/\mathrm{\Theta }_\mathrm{A}`$ also extends to a section $`s_\mathrm{A}`$ of $`K_Y^1`$. Since a critical point is a common zero of all $`_\mathrm{A}\stackrel{~}{W}`$’s, the number of critical points is the number of intersection points of the divisors $`_\mathrm{A}\stackrel{~}{W}=0`$ in $`(𝐂^\times )^{Nk}`$. Since there is no run-away behaviour of the potential for generic values of $`t_a`$, we do not expect a common zero of $`s_\mathrm{A}`$’s at infinity $`Y(𝐂^\times )^{Nk}`$. Thus, the number of critical points must be the same as the topological intersection number of the $`Nk`$ divisors $`s_\mathrm{A}=0`$ in $`Y`$ which is counted as $$\mathrm{Tr}_{\mathrm{mirror}}(1)^F=[Y],c_1(K_Y^1)^{Nk}.$$ (5.18) Therefore the number (5.18) must agree with the Euler number of $`X`$. It appears that this is not known in general. However, this certainly holds in every example one can check as long as $`c_1(X)0`$. Since the sigma model is well-defined and the gauge theory agrees with it only when the first Chern class of $`X`$ is positive semi-definite, the above result makes sense only in such cases $`c_1(X)0`$. To illustrate this, we consider the Hirzebruch surface $`F_a`$. As noted before, the charge matrix for $`F_a`$ is $`Q_{i1}={}_{}{}^{t}(1,0,1,a)`$ and $`Q_{i2}={}_{}{}^{t}(0,1,0,1)`$ and therefore the vectors $`v_i^\mathrm{A}`$ are given by $`v_1=(1,0)`$, $`v_2=(0,1)`$, $`v_3=(1,a)`$ and $`v_4=(0,1)`$. Then the superpotential is given by $$\stackrel{~}{W}=\mathrm{e}^{\mathrm{\Theta }_1}+\mathrm{e}^{t_2}\mathrm{e}^{\mathrm{\Theta }_2}+\mathrm{e}^{t_1}\mathrm{e}^{\mathrm{\Theta }_1a\mathrm{\Theta }_2}+\mathrm{e}^{\mathrm{\Theta }_2}.$$ (5.19) It is easy to see that the number of critical points are $`4`$ for $`a=0,1`$ and $`2a`$ for $`a2`$. Since the Euler number of $`F_a`$ is always $`4`$, the result is “correct” only for $`a=0,1,2`$. This is consistent because $`c_1(F_a)`$ is positive semi-definite only for $`a=0,1,2`$. ### 5.4 Non-Compact Case In the case where $`X`$ is non-compact there is an obvious run-away direction of the superpotential $`\stackrel{~}{W}`$. To see this, we note that $`X`$ is non-compact if and only if $`v_i`$’s generate a proper convex cone in $`\mathrm{N}_𝐑`$. This means that there is a point $`\mathrm{\Theta }_0`$ in $`\mathrm{M}_𝐑`$ such that $`\mathrm{\Theta }_0,v_i0`$ for all $`i`$. Now, consider the behaviour of the superpotential $`\stackrel{~}{W}`$ at $`\mathrm{\Theta }=t\mathrm{\Theta }_0`$ in the limit $`t+\mathrm{}`$. In the case where $`\mathrm{\Theta }_0,v_i>0`$ for all $`i`$, each term of the superpotential vanishes in the limit and this is the run-away direction. If there is some $`i`$ such that $`\mathrm{\Theta }_0,v_i=0`$ the superpotential stays finite but can be extremized by choosing an appropriate $`\mathrm{\Theta }_0`$. In any case, there is a run-away direction of the superpotential. We thus expect a continuous spectrum, which is indeed a property of the sigma model on a non-compact manifold. Below, we present two basic examples of non-compact toric manifolds. The first one, the total space of $`𝒪(d)`$ over $`𝐂\mathrm{P}^{N1}`$, provides a starting point of the discussion of the mirror for hypersurfaces in $`𝐂\mathrm{P}^{N1}`$ and more general toric complete intersections. The second one, the total space of $`𝒪(1)𝒪(1)`$ over $`𝐂\mathrm{P}^1`$, is important for the study of phase transition in the sense of . ### $`𝒪(d)`$ over $`𝐂\mathrm{P}^{N1}`$ The linear sigma model for the non-compact space given by the total space $`𝒪(d)`$ over $`𝐂\mathrm{P}^{N1}`$ is given by a $`U(1)`$ gauge theory with $`N`$ fields with charge $`+1`$ and one field with charge $`d`$. As discussed before we find that the mirror for this theory is an LG theory given by an LG theory with superpotential $$W=\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}+\mathrm{e}^{Y_P}$$ where $$dY_P+\underset{i=1}{\overset{N}{}}Y_i=t$$ (5.20) and $`Y_P`$ is the dual field to the charge $`d`$ matter and $`Y_i`$ for $`i=1,\mathrm{},N`$ are dual to the charge $`+1`$ matter fields. It is natural to use (5.20) to solve for $`Y_P`$. Let us define $$X_i=\mathrm{e}^{Y_i/d}$$ (5.21) for $`i=1,\mathrm{},N`$. Then we have $$\mathrm{e}^{Y_P}=\mathrm{e}^{t/d}X_1X_2\mathrm{}X_N$$ (5.22) and therefore $$W=X_1^d+X_2^d+\mathrm{}+X_N^d+\mathrm{e}^{t/d}X_1X_2\mathrm{}X_N$$ (5.23) However, we have to note that the field redefinition (5.21) is not single valued. In fact as we shift $`Y_iY_i+2\pi i`$ we transform $`X_i`$ by a primitive $`d`$-th root of unity $$X_iX_i\mathrm{e}^{2\pi i/d}.$$ However, from (5.22) we see that the product of the $`X_i`$ is well defined. Thus the LG theory we get is given by the above superpotential $`W`$ modded out by $`G=(𝐙_d)^{N1}`$ where $`G`$ is the group which acts on $`X_i`$ by all $`d`$-th roots of unity which preserve the product $`X_1X_2\mathrm{}X_N`$. As we discussed in section 4, the linear sigma model for the above non-compact theory and the one given by a hypersurface of degree $`d`$ in $`𝐂\mathrm{P}^{N1}`$ differ by deformations involving only the $`cc`$ ring elements, which thus do not affect the $`ac`$ ring, which we are studying through mirror symmetry. Thus we expect the same LG with superpotential (5.23) (modded out by $`G`$) to describe the mirror of hypersurface of degree $`d`$ in $`𝐂\mathrm{P}^{N1}`$. Note that for the case of Calabi-Yau manifolds $`d=N`$, and this gives the mirror construction of Greene and Plesser . However clearly there should be a difference between non-compact case and the compact theory. For example their conformal central charges (in $`N=2`$ units) differ by 2 in the Calabi-Yau case. What we will find is that indeed there is a difference and this is reflected in what is the appropriate field variable. In the compact case (i.e. the hypersurface case) we will find that the fundamental fields are $`X_i`$, whereas in the non-compact case the fundamental fields are $`Y_i`$ which are related to $`X_i`$ as defined in (5.21). This will be discussed after our discussion of the realization of D-branes in LG theories and their associated “BPS masses”. ### $`𝒪(1)+𝒪(1)`$ over $`𝐂\mathrm{P}^1`$ As a special case of interest, which has been studied in let us consider the mirror for $`𝒪(1)+𝒪(1)`$ over $`𝐂\mathrm{P}^1`$. This is given by the linear sigma model with a single $`U(1)`$ with four fields with charges $`(1,1,+1,+1)`$. For non-vanishing $`t`$ we can eliminate $`\mathrm{\Sigma }`$ as before and we obtain, in the dual formulation the superpotential $$W=X_1+X_2+X_3+X_4$$ with $`X_i=\mathrm{e}^{Y_i}`$ and $$X_1X_2=X_3X_4\mathrm{e}^t.$$ We can eliminate, say $`X_1`$ and write $$W=X_2+X_3+X_4+\mathrm{e}^tX_3X_4/X_2$$ Later in this paper we will return to this example and show how the prepotential for this model (i.e., the famous tri-logarithmic structure of ) can be obtained using the mirror potential we have found here. ## 6 D-branes, BPS Mass and LG Models In the case of Calabi-Yau manifolds one can consider special n-dimensional Lagrangian submanifolds, and they represent a class in $`H_n(M)`$. We can imagine a D-brane wrapping that class. In the sigma model we consider worldsheet with boundaries where the boundary can end on these manifolds. Such boundaries preserve the A-model supercharges. Given the fact that the neutral observables of the $`B`$ model correspond to $`n`$-form (after contraction with the holomorphic $`n`$-form) we have a natural pairing between $`B`$-model chiral fields, and $`A`$-model boundary states. The pairing is simply given by integrating the corresponding n-form on the corresponding n-cycle. Moreover, this integration can also be interpreted as the inner product of the boundary state defined by the D-brane and state defined (through topological twisting) by the $`B`$-model observable . In other words $$\gamma _i\|\varphi _\alpha =_{\gamma _i}\varphi _\alpha $$ where $`\gamma _i`$ represents an n-dimensional cycle giving rise to boundary state $`\gamma _i\|`$ and $`\varphi _\alpha `$ represents an element of the B-model chiral field corresponding to an $`n`$-form. There is one distinguished element for the chiral ring, the identity operator, which corresponds to the holomorphic $`n`$-form $`\mathrm{\Omega }`$ on the Calabi-Yau. In particular we have $$Z_i=\gamma _i\|1=_{\gamma _i}\mathrm{\Omega }$$ For type IIB compactifications on Calabi-Yau 3-folds, the $`Z_i`$ also represents central extension of supersymmetry algebra, in a sector with D-brane wrapping the $`\gamma _i`$ cycle. In the context of LG models, the natural topological theory corresponds to the B-model one . It is thus natural to ask the analog of this pairing in the context of the LG models. This is important in our applications because we have found an equivalence between the linear sigma models and LG models. Consider a Landau Ginzburg theory with superpotential $`W(x_i)`$ where $`x_i`$ are chiral fields, with $`i=1,\mathrm{},n`$. Then it is known that the chiral ring for this theory is given by : $$=𝐂[x_i]/dW$$ i.e., the ring generated by $`x_i`$ subject to setting $`dW=0`$. It is also known that there is a corresponding homology group which is equal to the dimension of the ring, namely consider the homology group $$H_n(𝐂^n,B)$$ where $`B`$ denotes the asymptotic region in $`𝐂^n`$ where $`ReW+\mathrm{}`$, i.e. n-cycles in the $`x_i`$ space with boundaries ending on $`B`$. Then the dimension of this space is equal to the dimension of the ring $``$. In fact as we will now discuss, there is a natural pairing between them and this identifies them as dual spaces. This is completely parallel to the case of the Calabi-Yau case and the pairing between mid-dimensional cycles and B-model observables (with equal left/right $`U(1)`$ charge). As is shown in , in fact we can construct analogs of Lagrangian submanifolds in the context of LG models as well, representing a basis for $`H_n(𝐂^n,B)`$ and preserving the A-combination of supercharges $`Q_{ac}=Q_{}+\overline{Q}_+`$. Moreover the image of these cycles on the W-plane is a straight line extending from the critical values of W to infinity along the positive real axis. Let $`\gamma _i`$ represent one such cycle. We would like now to discuss the analog of the pairing discussed above for the case of Calabi-Yau manifolds, between the B-model observables and A-model boundary states. Consider the pairing $$\mathrm{\Pi }_{i,a}=\gamma _i\|\varphi _a$$ where $`\varphi ^a`$ denotes a B-model chiral field observable and $`\gamma _i\|`$ denotes the boundary state corresponding to $`\gamma _i`$. It has been shown in that the $`\mathrm{\Pi }_b`$ satisfy the flatness equation in the context of $`tt^{}`$: $$_a\mathrm{\Pi }_b=(D_a+C_a)\mathrm{\Pi }_b=0\overline{}_a\overline{\mathrm{\Pi }}_b=(\overline{D}_a+\overline{C}_a)\overline{\mathrm{\Pi }}_b=0$$ where $`C_a`$ denotes the multiplication by the chiral field corresponding to $`\varphi _a`$ on the chiral fields and $`D_a`$ is a covariant derivative defined in . This was formulated in by showing that the above quantity can be computed by reducing to the case of quantum mechanics. The same result can also be derived more directly using the gymnastics leading to $`tt^{}`$ equation. In fact for the conformal case this has already been shown in and the same argument also applies to the massive cases (by considering the geometry of semi-infinite cigar, as in the derivation of $`tt^{}`$, with topological B-twisting on the tip of the cigar, as will be discussed in ). A special case of this overlap, is given by $`\mathrm{\Pi }_{i,1}`$, i.e. overlap with the state corresponding to the identity operator of the B-model. It has been shown in that in the conformal case $`\mathrm{\Pi }_{i,1}`$ is a holomorphic function of the couplings, i.e. it is independent of the coefficients of the superpotential $`\overline{W}`$ and depends only on the holomorphic couplings appearing in $`W`$. In the non-conformal case, this is no longer true. However even in the massive case one can obtain a holomorphic object by expanding near the conformal limit. This means formally considering $`\overline{W}\overline{\lambda }\overline{W}`$ and taking the limit $`\overline{\lambda }0`$. From this point on, when we refer to $`\mathrm{\Pi }_{i,a}`$ we have in mind this limit. It has been shown in that in such a case $`\mathrm{\Pi }_{i,1}`$ has a simple integral expression: $$\mathrm{\Pi }_{i,1}=_{\gamma _i}\mathrm{exp}(W)𝑑x_1\mathrm{}𝑑x_n.$$ (6.1) This is the Landau-Ginzburg analog of the BPS mass for the D-brane given by $`\gamma _i`$. Moreover, with a good choice of basis for chiral fields of the B-model (called the topological or flat coordinates) one has $`_{t^a}\mathrm{\Pi }_{i,1}=\mathrm{\Pi }_{i,a}={\displaystyle _{\gamma _i}}\varphi _a\mathrm{exp}(W)𝑑x_1\mathrm{}𝑑x_n,`$ (6.2) $`_{t^a}_{t^b}\mathrm{\Pi }_{i,1}=_{t^b}\mathrm{\Pi }_{i,a}=C_{ab}^c\mathrm{\Pi }_{i,c},`$ (6.3) where $`\varphi _a`$ denotes the chiral field corresponding to the deformation given by $`t^a`$. Note that this implies that the periods $`\mathrm{\Pi }_{i,a}`$ have the full information about the chiral ring structure constants. Note that in the limit $`W=0`$ the above expression for the period is the usual integral of the holomorphic $`n`$-form (which for $`𝐂^n`$ is $`\mathrm{\Omega }=dx_1\mathrm{}dx_n`$) on a cycle. The appearance of $`e^W`$ reflects the fact that in the presence of superpotential the supercharges get modified and field which corresponds to D-brane boundary conditions in addition to the delta function forcing $`x_r`$’s to the subspace $`\gamma _i`$ includes an additional factor $`e^W`$, i.e. $`e^W\delta (x\gamma )`$ is needed to be $`Q_{ac}=Q_{}+\overline{Q}_+`$ invariant. ### 6.1 Picard-Fuchs Equations for Periods Applying the previous discussion to the mirror of the non-compact toric varieties we have discussed, we can compute the periods of the branes in the mirror LG theory. We have the periods being given as $$\mathrm{\Pi }=\underset{i,b}{}dY_i\underset{b}{}\delta (\underset{i}{}Q_{ib}Y_it_b)\mathrm{exp}(\underset{i}{}e^{Y_i}),$$ (6.4) where we have suppressed the indices for cycles $`\gamma _i`$ and the identity operator associated with $`\mathrm{\Pi }`$. Consider instead $$\mathrm{\Pi }(\mu _j,t_b)=\underset{i,b}{}dY_i\underset{b}{}\delta (\underset{i}{}Q_{ib}Y_it_b)\mathrm{exp}(\underset{i}{}\mu _ie^{Y_i}),$$ (6.5) which satisfies for each $`b`$ the equations $$[\underset{Q_{ib}>0}{}(\frac{}{\mu _i})^{Q_{ib}}]\mathrm{\Pi }(\mu _j,t_b)=e^{t_b}[\underset{Q_{ib}<0}{}(\frac{}{\mu _i})^{Q_{ib}}]\mathrm{\Pi }(\mu _j,t_b).$$ (6.6) On the other hand by a shift in $`Y_i`$ by $`\mathrm{log}\mu _i`$ we can get rid of the $`\mu _i`$ dependence above, except for a shift in the delta function constraint. In other words we have $$\mathrm{\Pi }(\mu _i,t_b)=\mathrm{\Pi }_{i,1}(1,t_b\mathrm{log}\underset{i}{}\mu _i^{Q_{ib}}),$$ (6.7) which is the original periods we are interested in computing. This in particular means that (6.6) can be rewritten as differential operators involving the $`t_b`$ parameters. Together with the boundary conditions for large $`t_b`$’s these give an effective way of computing the periods. ### $`𝒪(3)`$ over $`𝐂\mathrm{P}^2`$ Just as an example of the previous equations, we can consider the equations satisfied by the periods of the mirror of sigma model on $`𝒪(3)`$ over $`𝐂\mathrm{P}^2`$. This is given by a $`U(1)`$ gauge theory with 4 matter fields with charges $`(3,1,1,1)`$. Using (6.6) we have $$\frac{}{\mu _2}\frac{}{\mu _3}\frac{}{\mu _4}\mathrm{\Pi }=e^t\frac{^3}{\mu _1^3}\mathrm{\Pi }.$$ (6.8) Defining $`\theta =d/dt`$ and noting that $`\mathrm{\Pi }`$ depends on $`\mu _i`$ in the combination $`t\mathrm{log}[\mu _2\mu _3\mu _4/\mu _1^3]`$ we obtain $$\theta ^3\mathrm{\Pi }=e^t(3\theta +2)(3\theta +1)(3\theta )\mathrm{\Pi }.$$ (6.9) ### 6.2 Prepotential for $`𝒪(1)+𝒪(1)`$ bundle over $`𝐂\mathrm{P}^1`$ We will now use the periods of the mirror to compute the prepotential for the non-compact Calabi-Yau threefold given by $`𝒪(1)+𝒪(1)`$ bundle over $`𝐂\mathrm{P}^1`$. The prepotential for this model is known to be given by trilogarithm function $$F(t)=P_3(t)+\underset{n>0}{}e^{nt}/n^3$$ (6.10) where $`P_3(t)`$ is a polynomial of degree 3 in $`t`$ (some of whose coefficients are ambiguous). The physical meaning of $`F(t)`$ as far as D-branes are concerned is as follows: $`\mathrm{\Pi }_2=t`$ denotes the (complexified) volume of D2-brane wrapping $`𝐂\mathrm{P}^1`$, whereas $$\mathrm{\Pi }_4=\frac{F}{t},$$ (6.11) which is the dual period, corresponds to the complexified quantum volume of a non-compact D4 brane intersecting the $`𝐂\mathrm{P}^1`$ at one point. The ambiguity in the coefficients of $`P_3(t)`$ reflects the infinity of the volume of this D4 brane. From the discussion of the BPS masses of D-branes, and the mirror for this model, which we found in the form of the LG model $$W=X_2+X_3+X_4+e^tX_3X_4/X_2$$ (6.12) we are led to consider the appropriate periods given by integrals $$\mathrm{\Pi }=\frac{dX_2dX_3dX_4}{X_2X_3X_4}e^W$$ (6.13) where the measure is obtained by noting that the correct variables are the $`Y_i`$’s. It is more convenient to consider $`\mathrm{\Pi }/t`$ and integrating over $`X_4`$: $$\frac{\mathrm{\Pi }}{t}=e^t\frac{dX_2dX_3dX_4}{X_2^2}e^{(X_2+X_3+X_4+e^tX_3X_4/X_2)}=e^t\frac{dX_2dX_3}{X_2^2}\delta (1+e^t\frac{X_3}{X_2})e^{(X_2X_3)}$$ (6.14) performing the integral over $`X_3`$ we have $$\frac{\mathrm{\Pi }}{t}=\frac{dX_2}{X_2}e^{X_2(1e^t)}$$ (6.15) There are two cycle one can consider here: Integrating around the origin of $`X_2`$ we obtain 1. This corresponds to the D2 brane BPS mass, namely $`\mathrm{\Pi }_2/t=1`$. The other period corresponds to integrating from $`0`$ to $`\mathrm{}`$ on the $`X_2`$ plane, which is the dual cycle corresponding to D4 brane. To do that consider taking another derivative of $`\mathrm{\Pi }`$: $$\frac{^2\mathrm{\Pi }}{t^2}=\frac{^3F}{t^3}=e^t𝑑X_2e^{X_2(1e^t)}=\frac{e^t}{(1e^t)}$$ (6.16) In agreement with the known result for $`F(t)`$ (6.10). ### 6.3 Compact versus Non-Compact As we discussed in section 4 compact complete intersections in toric varieties are closely related to the corresponding non-compact ones. Moreover all the quantities that are unaffected by the $`cc`$ deformations in the compact theory, should be computable in the non-compact theory. This is true at the level of embedding of twisted chiral operators of the compact theory in the non-compact one and the example is the relation (4.11) by presented in section 4. The same should be true about the overlap of states with B-type boundary condition with the ground state vacua, which according to is independent of $`cc`$ deformations. From the above discussion of periods, one natural state to consider is the one corresponding to the identity operator in the context of the compact theory. As discussed in section 4 this corresponds to the state $`\|1_{Compact}\|\delta _{NonCompact}`$ in the non-compact theory. We now consider the pairing between B-type boundary states and the ground states (represented using A-model) in the sigma model. The B-type boundary states in the non-compact theory $`\stackrel{~}{\gamma }_i\|`$ will flow to some boundary states in the compact theory $`\gamma _i\|`$. By the $`ϵ`$ independence we thus have $$\gamma _i\|\mathrm{}_{Compact}=\stackrel{~}{\gamma }_i\|\mathrm{}\|\delta _1\mathrm{}\delta _k_{NonCompact}$$ (6.17) In particular we learn that the BPS mass (which corresponds to the overlap of the state associated to identity operator with D-brane boundary state) for the compact theory is given by $$\mathrm{\Pi }_{i,1}=\gamma _i\|1_{Compact}=\stackrel{~}{\gamma }_i\|\delta _1\mathrm{}\delta _k_{NonCompact}$$ (6.18) Since we have found the mirror of non-compact theory in terms of an LG model, from our previous discussion it follows that $$\mathrm{\Pi }_{i,1}=_{\stackrel{~}{\gamma }_i}\delta _1\mathrm{}\delta _k\mathrm{exp}(W)$$ (6.19) Note that this result, for the case of linear sigma models corresponding to degree $`d`$ hypersurfaces in toric varieties described by a single $`U(1)`$ gauge theory can also be written as $$\mathrm{\Pi }_{Compact}=d\frac{}{t}\mathrm{\Pi }_{NonCompact}$$ (6.20) This result is well known in the context of local mirror symmetry for Calabi-Yau manifolds and should be viewed as a generalization of it. The equations (6.19) and its specialization (6.20) will be the fundamental results we need in completing our discussion for deriving the mirror of hypersurfaces (and complete intersections) in toric varieties. Note that using the fact that the periods of the non-compact theory satisfy appropriate Picard-Fuchs equations, as derived in (6.6) we can use the above result (6.20) to write a Picard-Fuchs equation satisfied by the periods of the mirror for the compact case. ## 7 Mirror Symmetry for Complete Intersections In this section we complete our derivation of the mirror theory corresponding to complete intersections in toric varieties. First we discuss the simple case of degree $`d`$ hypersurfaces in $`𝐂\mathrm{P}^{N1}`$ to illustrate how the idea works. Then we show how a subset of complete intersection sigma models have a similar mirror in the form of a simple Landau-Ginzburg orbifolds in flat space. We then show that the most general construction can also be described in terms of LG theory on non-compact Calabi-Yau manifolds. ### 7.1 Hypersurfaces in $`𝐂\mathrm{P}^{N1}`$ Consider a degree $`d`$ hypersurface in $`𝐂\mathrm{P}^{N1}`$. As discussed before, all the ring structure and BPS masses can be computed in the associated non-compact theory. We have $`N+1`$ dual matter fields $`Y_i`$ with $`i=1,..,N`$ and $`Y_P`$ and one $`\mathrm{\Sigma }`$ field. From (6.20) we can compute the BPS masses for the compact theory in the form $`\mathrm{\Pi }`$ $`=`$ $`d{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle d\mathrm{\Sigma }dY_P\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\mathrm{exp}\left(\stackrel{~}{W}\right)}`$ (7.1) $`=`$ $`d{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle dY_P\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\delta \left(\underset{i=1}{\overset{N}{}}Y_idY_Pt\right)\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}\mathrm{e}^{Y_P}\right)}.`$ Constraint $`_iY_idY_P=t`$ can be solved, as discussed earlier, by $`\mathrm{e}^{Y_i}=X_i^d,`$ (7.2) $`\mathrm{e}^{Y_P}=\mathrm{e}^{t/d}X_1\mathrm{}X_N.`$ (7.3) The map from $`X_i`$ to $`\mathrm{e}^{Y_i}`$ and $`\mathrm{e}^{Y_P}`$ is one-to-one up to the action of $`(𝐙_d)^{N1}`$ on $`X_i`$ defined by $$X_i\omega _iX_i,\omega _i^d=1,\omega _1\mathrm{}\omega _N=1.$$ (7.4) Then, we obtain $`\mathrm{\Pi }`$ $`=`$ $`d{\displaystyle \frac{\mathrm{d}}{\mathrm{d}t}}{\displaystyle \underset{i=1}{\overset{N}{}}\frac{\mathrm{d}X_i}{X_i}\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}X_i^d\mathrm{e}^{t/d}\underset{i=1}{\overset{N}{}}X_i\right)}`$ (7.5) $`=`$ $`\mathrm{e}^{t/d}{\displaystyle \underset{i=1}{\overset{N}{}}\mathrm{d}X_i\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}X_i^d\mathrm{e}^{t/d}\underset{i=1}{\overset{N}{}}X_i\right)}.`$ This is the period for Landau-Ginzburg model with superpotential $$W_{LG}=X_1^d+\mathrm{}+X_N^d+\mathrm{e}^{t/d}X_1\mathrm{}X_N,$$ (7.6) or more precisely the LG orbifold by the $`(𝐙_d)^{N1}`$ action (7.4). What we see here is that the compact theory and non-compact theory both have the same Landau-Ginzburg potential, but the measure for the fundamental fields, which determine the measure are different in the two cases. In other words the LG orbifold with fields $`X_i`$ as the fundamental fields (rather than $`Y_i`$ as in the non-compact case) has the same ring structure and BPS masses as the degree $`d`$ hypersurfaces in $`𝐂\mathrm{P}^{N1}`$. Let us examine the vacua of this LG orbifold. The equation $`\mathrm{d}W_{LG}=0`$ has $`(Nd)`$ solutions at $$X_1^d=\mathrm{}=X_N^d=\frac{\mathrm{e}^{t/d}}{d}X_1\mathrm{}X_N=:S,S^{Nd}=(d)^d\mathrm{e}^t$$ (7.7) and, for $`d2`$, one solution at $$X_1=\mathrm{}=X_N=0.$$ (7.8) The $`(Nd)`$ critical points (7.7) are all non-degenerate and correspond to massive vacua that break spontaneously the $`𝐙_{2(Nd)}`$ axial R-symmetry to $`𝐙_2`$. The critical point at $`X_i=0`$ is degenerate for $`d>2`$ and corresponds to a non-trivial fixed point. The behaviour of the theory in the IR limit depends on the relation of $`d`$ and $`N`$, as follows. $``$ For $`d=1`$ where $`M=𝐂\mathrm{P}^{N2}`$, there are only $`N1`$ massive vacua. If we integrate out $`X_N`$, we obtain the constraint $`X_{N1}=\mathrm{e}^t/(X_1\mathrm{}X_{N2})`$ and the effective superpotential for the remaining fields is $$W=X_1+\mathrm{}+X_{N2}\mathrm{e}^t/(X_1\mathrm{}X_{N2}).$$ (7.9) This is nothing but the $`A_{N2}`$ affine Toda superpotential. Thus, we have reproduced the mirror symmetry of $`𝐂\mathrm{P}^{N2}`$ model and affine Toda theory. $``$ For $`2d<N`$ where $`M`$ has positive first Chern class and the sigma model is asymptotically free, there are $`(Nd)`$ massive vacua and a vacuum at $`X_i=0`$. At $`X_i=0`$, the last term $`\mathrm{e}^{t/d}X_1\mathrm{}X_N`$ of the potential (7.6) is irrelevant compared to the first $`N`$-terms. This can be seen by computing its dimension in the superconformal field theory we will study below. (More intuitively, this is because $`t`$ flows at low energies to large negative values.) Thus, the theory at $`X_i=0`$ corresponds to the $`(𝐙_d)^{N1}`$-orbifold of the LG model with the quasi-homogeneous superpotential $$W_{IR}=X_1^d+\mathrm{}+X_N^d.$$ (7.10) In particular, we see the enhancement of the axial $`𝐙_{2(Nd)}`$ R-symmetry to $`U(1)_A`$ symmetry. For $`d=2`$ the critical point at $`X_i=0`$ is non-degenerate and does not corresponds to a non-trivial fixed point. However, because of the orbifolding, it could correspond to multiple vacua. $``$ For $`d=N`$ where $`M`$ is a Calabi-Yau manifold, there are no massive vacua but one massless vacuum at $`X_i=0`$. At $`X_i=0`$, the last term of (7.6) is equally relevant compared to the first $`N`$ terms. Thus $`t`$ remains as the parameter describing a marginal deformation of the SCFT. $``$ For $`d>N`$ where $`M`$ has negative first Chern class, there are $`(dN)`$ massive vacua and one massless vacuum at $`X_i=0`$. At $`X_i=0`$, the first $`N`$ terms of (7.6) are irrelevant compared to the last term. Thus, the IR fixed point is described by the theory with superpotential $$W_{IR}=\mathrm{e}^{t/d}X_1\mathrm{}X_N.$$ (7.11) The vacuum equation $`\mathrm{d}W_{IR}=0`$ is solved if two of $`X_i`$’s vanish. Namely, the theory is a free SCFT on $`𝐂^{N2}`$. This is expected since if $`c_1(M)<0`$ the sigma model is IR free. Thus, we have seen that the LG orbifold with the superpotential (7.6) and group $`(𝐙_d)^{N1}`$ captures the physics of the sigma model for all values of $`d`$. Also, it describes both the massive vacua and the massless vacuum that flows to a non-trivial (or trivial) IR fixed point. Below, we shall analyze the spectrum of the SCFT with the superpotential (7.10) that appears as the non-trivial fixed point for hypersurfaces with $`c_1>0`$. #### 7.1.1 The Spectrum of Chiral Primary Fields for Hypersurfaces in $`𝐂\mathrm{P}^{N1}`$ We can use the method developed in to analyze the spectrum of the chiral primary fields. We are considering the LG model with the superpotential $$W=X_1^d+\mathrm{}+X_N^d$$ (7.12) divided by the orbifold group $`(𝐙_d)^{N1}`$ acting on $`X_i`$’s as $`X_i\omega ^{\alpha _i}X_i`$ where $$\omega =\mathrm{e}^{2\pi i/d},\underset{i=1}{\overset{N}{}}\alpha _i=0(\mathrm{mod}d).$$ (7.13) Recall that we are in the non-standard convention of chirality as the LG model; the fields $`X_i`$ are twisted chiral superfields of R-charge $`q_i=\overline{q}_i=1/d`$.<sup>1</sup><sup>1</sup>1$`q`$ and $`\overline{q}`$ are the left/right R-charges related to the vector/axial R-charges by $`q_V=\overline{q}+q`$ and $`q_A=\overline{q}q`$. The model before orbifolding has only the $`ac`$ ring given by $$=𝐂[X_1,\mathrm{},X_N]/_iW=\underset{0p_id2}{}𝐂X_1^{p_1}\mathrm{}X_N^{p_N},$$ (7.14) and the central charge is $`c/3=_i(12/d)=N(d2)/d`$. ### Supersymmetric Ground States The supersymmetric (or Ramond) ground states of the $`(𝐙_d)^{N1}`$ orbifold come both from untwisted and twisted sectors. The Ramond ground states in the untwisted sector are build from the canonical ground state $`\|0_R^1`$ of R-charge $`\overline{q}=q=c/6=N(d2)/2d`$ by multiplying the $`(𝐙_d)^{N1}`$ invariant elements of (7.14). Obviously they are the powers of $`X_1\mathrm{}X_N`$. Thus, there are $`d1`$ Ramond ground states in the untwisted sector; $$(X_1\mathrm{}X_N)^p\|0_R^1,p=0,1,\mathrm{},d2.$$ (7.15) The R-charge is just the sum of $`c/6`$ and $`Np(1/d)`$; $$\overline{q}_p=q_p=\frac{N}{d}\left(p\frac{d2}{2}\right).$$ (7.16) Let us next consider the $`h`$-twisted sector for $`h=(\omega ^{\alpha _1},\mathrm{},\omega ^{\alpha _N})(𝐙_d)^{N1}`$. The $`h`$-invariant fields are $`X_i`$ with $`\omega ^{\alpha _i}=1`$. There is a (not necessarily physical) ground state $`\|0_R^h`$ in the sector on which $`g(𝐙_d)^{N1}`$ acts as the multiplication by $`detg\|_h`$ where $`g\|_h`$ is the action of $`g`$ restricted on $`h`$-invariant fields. ($`detg\|_h=1`$ if there is no $`h`$-invariant field). The Ramond ground states in the $`h`$-twisted sector are $`(𝐙_d)^{N1}`$ invariant states of the form $`X_{i_1}^{p_1}\mathrm{}X_{i_s}^{p_s}\|0_R^h`$ where $`X_{i_a}`$’s are $`h`$-invariant fields and $`0p_ad2`$. If $`\omega ^{\alpha _i}=1`$ for some $`i`$ (say $`i=1`$), then $`\omega ^{\alpha _j}1`$ for some other $`j`$ (say $`j=N`$) if $`h1`$. Then $`X_1`$ is $`h`$-invariant but $`X_N`$ is not. For $`g=(\omega ,1,\mathrm{},1,\omega ^1)`$, we have $`detg\|_h=\omega `$. But then $`g`$ acts on the state $`X_1^{p_1}\mathrm{}\|0_R^h`$ by multiplication by $`\omega ^{p_1+1}1`$. So there is no $`(𝐙_d)^{N1}`$-invariant state. Thus, Ramond ground states exist only for $`h`$ such that $`\omega ^{\alpha _i}1`$ for all $`i`$. For such an $`h`$, there is no $`h`$-invariant field and therefore there is a unique Ramond ground state $$\|0_R^h.$$ (7.17) Its R-charge can be obtained from the formula (3.2) in $$\overline{q}_h=q_h=\left(\frac{\underset{i}{}\alpha _i}{d}\frac{N}{2}\right).$$ (7.18) Since $`_i\alpha _i=0`$ mod $`d`$, we have $`_i\alpha _i=n_hd`$ where $`n_h`$ is an integer in the range $`N/dn_hN(11/d)`$. Thus, the vector R-charge $`q^V=\overline{q}_h+q_h`$ is an integer $`N2n_h`$ in the range $$\|q_h^V\|c/3<N2.$$ (7.19) Let us count the number of such $`h(𝐙_d)^{N1}`$. Let $`𝒞_N`$ be the set of such $`h`$’s. Let us consider a sequence $`(\omega ^{\alpha _1},\mathrm{},\omega ^{\alpha _{N1}})`$ such that any entry is not equal to $`1`$. If $`\alpha _{}:=\alpha _i`$ is equal to $`0`$ mod $`d`$, this determines an element of $`𝒞_{N1}`$. Otherwise this determines a unique element $`(\omega ^{\alpha _1},\mathrm{},\omega ^{\alpha _{N1}},\omega ^\alpha _{})`$ of $`𝒞_N`$. Thus we obtain the recursion relation $`(d1)^{N1}=\mathrm{}(𝒞_{N1})+\mathrm{}(𝒞_N)`$ which is solved by $$\mathrm{}(𝒞_N)=(1)^N\frac{(1d)^N(1d)}{d}.$$ (7.20) This is the total number of Ramond ground states from twisted sectors. Let us compute the Witten index $`\mathrm{Tr}(1)^F`$. We normalize $`(1)^F=1`$ on $`\|0_R^1`$. Then the ground states from untwisted sector has $`(1)^F=(1)^N`$ . Thus, the index of the SCFT is $$\mathrm{Tr}_{_{\mathrm{SCFT}}}(1)^F=(d1)+\frac{(1d)^N(1d)}{d}=\frac{(1d)^N+d^21}{d}.$$ (7.21) Together with the $`(Nd)`$ massive vacua having $`(1)^F=1`$, we obtain the total index $$\mathrm{Tr}(1)^F=(Nd)+\frac{(1d)^N+d^21}{d}=\frac{(1d)^N+Nd1}{d}.$$ (7.22) This indeed agrees with the Euler number of the hypersurface $`M`$. Note that our original sigma model possesses the unbroken $`U(1)_V`$ R-symmetry and it counts $`p+q`$ of the Hodge number $`(p,q)`$ of the harmonic form representing a vacuum. For a hypersurface in $`𝐂\mathrm{P}^{N1}`$ the off diagonal Hodge number is non-zero only for $`p+q=N2`$. Thus, the ground state $`\|0_R^h`$ with non-zero $`q_h^V=N2n_h`$ corresponds to a harmonic $`(p,q)`$-form with $`qp=q_h^V`$ and $`p+q=N2`$. In other words, the number of such $`h`$’s for a fixed $`(p,q)`$ must be equal to the Hodge number $`h^{p,q}`$. This is indeed easy to check explicitly for small values of $`N`$ and $`d`$ (e.g. using a formula for the Hodge numbers in ). For illustration, let us present the result for two cases $`N=4`$, $`d=3`$ and $`N=5`$, $`d=3`$: N=4, d=3 $`M`$ is the cubic surface in $`𝐂\mathrm{P}^3`$ which is known as $`E_6`$ del Pezzo surface. It has $`h^{0,0}=h^{2,2}=1`$, $`h^{1,1}=7`$ and the off-diagonal Hodge numbers are all zero. Thus, classically there are one ground state with $`(q,\overline{q})=(\pm 1,\pm 1)`$ and seven ground states with $`(q,\overline{q})=(0,0)`$. In the quantum theory there is a single massive vacuum and a massless vacuum. The massless vacuum corresponds to a SCFT of $`c/3=4/3`$. There are one untwisted ground state with $`(q,\overline{q})=(\pm 2/3,\pm 2/3)`$ and six twisted ground states with $`(q,\overline{q})=(0,0)`$. N=5, d=3 $`M`$ is a cubic hupersurface in $`𝐂\mathrm{P}^4`$ which has $`h^{p,p}=1`$ for $`p=0,1,2,3`$ and $`h^{2,1}=h^{1,2}=5`$. Thus, classically there are one ground state with $`(q,\overline{q})=(p3/2,p3/2)`$ for $`p=0,1,2,3`$, five with $`(q,\overline{q})=(1/2,1/2)`$ and five with $`(q,\overline{q})=(1/2,1/2)`$. In the quantum theory, there are two massive vacua and a massless vacuum. The massless vacuum corresponds to a SCFT with $`c/3=5/3`$. There are one untwisted ground state with $`(q,\overline{q})=(\pm 5/6,\pm 5/6)`$, five twisted states with $`(q,\overline{q})=(1/2,1/2)`$ and five twisted states with $`(q,\overline{q})=(1/2,1/2)`$. ### $`ac`$ Primaries Since all the vector R-charges are integral, the $`ac`$ primary states are in one to one correspondence with the Ramond ground states by spectral flow. The spectral flow simply changes the R-charges by $`\overline{q}_{\mathrm{𝑎𝑐}}=\overline{q}_R+c/6`$ and $`q_{\mathrm{𝑎𝑐}}=q_Rc/6`$. Thus, we have the following $`ac`$ primary states with R-charges: From the untwisted sector $$(X_1\mathrm{}X_N)^p\|0_{\mathrm{𝑎𝑐}}^1;\overline{q}_p=q_p=\frac{Np}{d},$$ (7.23) ($`p=0,\mathrm{},d2`$), and from the twisted sectors $$\|0_{\mathrm{𝑎𝑐}}^h;\overline{q}_h=q_h+c/3=N\left(1\frac{1}{d}\right)n_h.$$ (7.24) where $`n_h`$ is the integer in the range $`N/dn_hN(11/d)`$ defined above. ### $`cc`$ Primaries Since the axial R-charges are not necessarily integers, we must separately consider $`cc`$ primary states. The $`cc`$ primaries in the $`h`$-twisted sector for $`h=(\omega ^{\alpha _1},\mathrm{},\omega ^{\alpha _N})`$ (including $`h=1`$) can be found in the same way as the search for Ramond ground states in $`hj^1`$-twisted sector where $`j=(\omega ,\mathrm{},\omega )`$ . There is a unique $`cc`$ primary $$\|0_{\mathrm{𝑐𝑐}}^h,$$ (7.25) for each $`h`$ such that $`\beta _i:=\alpha _i1`$ obey $`_{i=1}^N\beta _i=N`$ and $`\beta _i0`$ ($`\mathrm{mod}d`$). One can choose $`\beta _i`$ in the range $`1\beta _id1`$ and then $`_i\beta _i=N+m_hd`$ where $`m_h`$ is an integer in the range $`2N/dm_hN`$. The R-charge of the $`cc`$ primary state is then $$\overline{q}=q=\frac{c}{6}\underset{i=1}{\overset{N}{}}\left(\frac{\beta _i}{d}\frac{1}{2}\right)=Nm_h.$$ (7.26) If $`N`$ is divisible by $`d`$, there are extra $`cc`$ primaries from the $`h=(\omega ,\mathrm{},\omega )`$-twisted sector; $$(X_1\mathrm{}X_N)^p\|0_{\mathrm{𝑐𝑐}}^h,p=0,\mathrm{},d2.$$ (7.27) which have R-charges $$\left(\genfrac{}{}{0pt}{}{\overline{q}}{q}\right)=\frac{c}{12}\frac{c}{12}\pm \frac{Np}{d}.$$ (7.28) The R-charges (7.26) and (7.28) are integers in the range $`0q,\overline{q}c/6`$. ### Marginal Deformations The marginal deformation of the theory preserving the $`(2,2)`$ superconformal symmetry is done by an $`ac`$ primary field of $`\overline{q}=q=1`$ or a $`cc`$ primary field of $`\overline{q}=q=1`$. From the above list, it is easy to see that there are no such $`ac`$ primaries except the special case $`N=6,d=3`$ where twenty $`n_h=3`$ twisted states does correspond to the marginal deformation. This is the case which corresponds to the conformal field theory of $`K3`$ . On the other hand, there are always $`cc`$ primaries with $`\overline{q}=q=1`$; they are the states (7.25) with $`m_h=N1`$. This corresponds to $`\alpha _i`$’s such that $`_i(d\alpha _i)=d`$ for $`0d\alpha _id2`$. The number of such $`\alpha _i`$’s is the same as the number of independent polynomial deformations of degree $`d`$ equation in $`N`$ variables. As we will see below, they correspond to the complex structure deformations of the hypersurface $`M`$. Incidentally all the cases were $`d`$ divides $`N`$ gives rise to a conformal field theory in the IR which is mirror to a Calabi-Yau manifold. Let $`N/d=k`$ be an integer. Then we obtain a conformal field theory corresponding to a Calabi-Yau with complex dimension $`(d2)k`$. In fact it is mirror to a Calabi-Yau which is an orbifold of $`k`$ copies of degree $`d`$ hypersurface in $`𝐂\mathrm{P}^{d1}`$. For example, as noted in the case $`N=9,d=3`$ is the mirror to a rigid Calabi-Yau threefold. For such cases we can now find a geometric interpretation of the mirror: The mirror of these rigid Calabi-Yau manifolds can be viewed as the infrared limit of sigma model on the corresponding hypersurface with positive $`c_1`$. #### 7.1.2 Mirror Symmetry of Orbifold Minimal Models as IR Duality An LG description of the IR fixed point in this model is actually available also in the original linear sigma model (see also ). This is the extension of the basic argument for CY/LG correspondence to the sigma model of $`c_1>0`$ hypersurfaces in $`𝐂\mathrm{P}^{N1}`$. For $`d<N`$, the Kahler parameter $`r`$ flows at low energies to large negative values $`r0`$. There we find $`(Nd)`$ massive vacua at large values of $`\sigma `$; $`\sigma ^{Nd}=\mathrm{e}^t(d)^d`$. Actually, there is also a vacuum at $`\sigma =0`$ but $`\|p\|^2=\|r\|/d`$. (Presumably the equation for $`\sigma `$ should be replaced by $`\sigma ^{N1}=\mathrm{e}^t(d)^d\sigma ^{d1}`$, which is the chiral ring relation for the sigma model on the hypersurface $`M`$ . This indeed has a solution at $`\sigma =0`$ for $`d>1`$.) Since $`p`$ has electric charge $`d`$, the $`U(1)`$ gauge symmetry is broken to $`𝐙_d`$ by the Higgs mechanism. The superfield $`P`$ is massive while other fields of unit charge $`\mathrm{\Phi }_i`$ are massless and become the relevant fields to describe the low energy theory. Because of the expectation value of $`P`$, there is a non-trivial superpotential for $`\mathrm{\Phi }_i`$’s $$W=\sqrt{\|r\|/d}\left(\mathrm{\Phi }_1^d+\mathrm{}+\mathrm{\Phi }_N^d\right).$$ (7.29) Namely, the theory flows in the IR limit to $`(Nd)`$ empty theories and the $`𝐙_d`$ orbifold of the LG model with the superpotential (7.29). This is very similar to the structure found in our effective theory (7.6) except that the group is now a single $`𝐙_d`$. Since the two LG orbifolds, one by $`𝐙_d`$ and the other by $`(𝐙_d)^{N1}`$, arise as IR fixed points of the same theory, they must agree with each other. Namely, the $`𝐙_d`$-orbifold and the $`(𝐙_d)^{N1}`$-orbifold must be mirror to each other.<sup>2</sup><sup>2</sup>2$`\mathrm{\Phi }_i`$ are chiral superfields and the model (7.29) has the standard convention of chirality. This mirror symmetry is actually the special case of the mirror symmetry between orbifolds of minimal models mentioned in section 2. It is straightforward to compute the spectrum of chiral primary fields in the $`𝐙_d`$ orbifold model. The result is of course in agreement with the one for the $`(𝐙_d)^{N1}`$ orbifold. Untwisted sector states in the $`𝐙_d`$ orbifold corresponds to twisted sector states in the $`(𝐙_d)^{N1}`$ orbifold, and vice versa. For instance, the untwisted $`cc`$ primary field $`\mathrm{\Phi }_1^{p_1}\mathrm{}\mathrm{\Phi }_N^{p_N}`$ ($`0p_id2`$) corresponds to the $`h=(\omega ^{\alpha _1},\mathrm{},\omega ^{\alpha _N})`$-twisted $`cc`$ primary state (7.25) in the $`(𝐙_d)^{N1}`$ orbifold where $`p_i=d\alpha _i`$. This is actually what is expected since $`\mathrm{\Phi }_i`$ and $`\mathrm{e}^{Y_i}=X_i^d`$ are dual to each other and a momentum (power) of one corresponds to a winding (twist) of the other. The interpretation of the marginal deformation by $`cc`$ primaries is clear in this picture. The marginal $`cc`$ deformation corresponds to deformation of the LG model (7.29) by $`\mathrm{\Phi }_1^{p_1}\mathrm{}\mathrm{\Phi }_N^{p_N}`$ with $`_ip_i=d`$ and $`0p_id2`$. This corresponds to the deformation of the superpotential by $`P\mathrm{\Phi }_1^{p_1}\mathrm{}\mathrm{\Phi }_N^{p_N}`$ in the original linear sigma model for $`M`$ which is nothing but the polynomial deformation of the defining equation of $`M`$. Actually, any deformation of the complex structure is of this form.<sup>3</sup><sup>3</sup>3This can be seen as follows (we thank R. Pandharipande for explanation). Using the long exact sequence for $`0T_MT_{𝐂\mathrm{P}^{N1}}\|_M𝒪_M(d)0`$, it is enough to show $`H^1(M,T_{𝐂\mathrm{P}^{N1}}\|_M)=0`$. From $`0𝒪𝒪(1)^NT_{𝐂\mathrm{P}^{N1}}0`$, this reduces to $`H^1(M,𝒪_M(1))=0`$ and $`H^2(M,𝒪_M)=0`$. These further reduce via $`0𝒪(d)𝒪𝒪_M0`$ to $`H^{1+i}(𝐂\mathrm{P}^{N1},𝒪(1i))=0`$ and $`H^{2+i}(𝐂\mathrm{P}^{N1},𝒪(1di))=0`$ for $`i=0,1`$. That this holds for $`d<N`$ is a standard fact (e.g. by Kodaira-Nakano vanishing theorem ). Within $`dN`$, the only case this fails is $`N=d=4`$, the famous example $`M=K3`$. Thus, the $`cc`$ deformation in the sigma model on a $`c_1>0`$ projective hypersurface is in one to one correspondence with the complex structure deformation, as in the case of Calabi-Yau sigma models. It would be interesting to better understand why this holds in this case and investigate whether this is a general fact. ### 7.2 Complete Intersections LG orbifold description of the mirror is possible also for a class of complete intersections in toric varieties. Let $`X`$ be a toric variety defined by the charge matrix $`Q_{ia}`$ ($`i=1,\mathrm{},N`$, $`a=1,\mathrm{},k`$). We consider a complete intersection $`M`$ in $`X`$ defined by the equations $`G_b=0`$ ($`b=1,\mathrm{},k`$) as many as the number of the $`U(1)`$ gauge groups to define $`X`$. Here $`G_b`$ is a polynomial of $`\mathrm{\Phi }_i`$ of a certain “degree” $`d_{ba}`$ (that is, $`G_b`$ has charge $`d_{ba}`$ for the $`a`$-th $`U(1)`$ gauge group) where we assume $`d_{ba}`$ to be an invertible matrix. The dual of the non-compact theory is described by $`N+k`$ twisted chiral fields $`Y_i`$ and $`Y_{P_b}`$ dual to $`\mathrm{\Phi }_i`$ and $`P_b`$ and $`k`$ field strengths $`\mathrm{\Sigma }_a`$. It has the twisted superpotential $$\stackrel{~}{W}=\underset{a=1}{\overset{k}{}}\mathrm{\Sigma }_a\left(\underset{i=1}{\overset{N}{}}Q_{ia}Y_i\underset{b=1}{\overset{k}{}}d_{ba}Y_{P_b}t_a\right)+\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}+\underset{b=1}{\overset{k}{}}\mathrm{e}^{Y_{P_b}}.$$ (7.30) The period integral relevant for the compact theory is given by $$\mathrm{\Pi }=\underset{a=1}{\overset{k}{}}\mathrm{d}\mathrm{\Sigma }_a\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\underset{b=1}{\overset{k}{}}\mathrm{d}Y_{P_b}\delta _1\mathrm{}\delta _k\mathrm{exp}\left(\stackrel{~}{W}\right),$$ (7.31) where $`\delta _b=_{a=1}^kd_{ba}\mathrm{\Sigma }_a`$. Then, the period is expressed as follows (a quick derivation of this is given in the next subsection in more general cases); $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\underset{b=1}{\overset{k}{}}\mathrm{d}Y_{P_b}\mathrm{e}^{Y_{P_1}}\mathrm{}\mathrm{e}^{Y_{P_k}}}`$ (7.32) $`\times {\displaystyle \underset{a=1}{\overset{k}{}}}\delta ({\displaystyle \underset{i=1}{\overset{N}{}}}Q_{ia}Y_i{\displaystyle \underset{b=1}{\overset{k}{}}}d_{ba}Y_{P_b}t_a)\times \mathrm{exp}({\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{e}^{Y_i}{\displaystyle \underset{b=1}{\overset{k}{}}}\mathrm{e}^{Y_{P_b}}).`$ Suppose there are matrices $`m_{ji}`$, $`n_{jb}`$ of non-negative integers such that $$\underset{i=1}{\overset{N}{}}m_{ji}Q_{ia}=\underset{b=1}{\overset{k}{}}n_{jb}d_{ba}.$$ (7.33) Then, the constraints for $`Y_i`$ and $`Y_{P_b}`$ can be solved as $$\mathrm{e}^{Y_i}=\underset{j=1}{\overset{N}{}}X_j^{m_{ji}},\mathrm{e}^{Y_{P_b}}=\mathrm{e}^{d_{ab}^1t_a}\underset{j=1}{\overset{N}{}}X_j^{n_{bj}}.$$ (7.34) Furthermore, if $`m_{ji}`$ is an invertible matrix and $`_bn_{jb}=1`$ for each $`j`$ (thus $`n_{jb}`$ is either $`0`$ or $`1`$), the period integral can be expressed as $$\mathrm{\Pi }=\underset{i=1}{\overset{N}{}}\mathrm{d}X_i\mathrm{exp}\left(\stackrel{~}{W}\right),$$ (7.35) for $$\stackrel{~}{W}=\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}X_j^{m_{ji}}+\underset{b=1}{\overset{k}{}}\mathrm{e}^{d_{ab}^1t_a}\underset{j=1}{\overset{N}{}}X_j^{n_{jb}}.$$ (7.36) Thus, under the condition that the charges $`Q_{ia}`$ and the degrees $`d_{ba}`$ admit such matrices $`m_{ji}`$ and $`n_{bj}`$ as the solution to (7.33), the period of the dual theory is the same as that of the LG orbifold with the twisted chiral superpotential (7.36). The orbifold group $`\stackrel{~}{\mathrm{\Gamma }}`$ is generated by $`X_j\omega _jX_j`$ where $`\omega _j`$ are phases satisfying $`_j\omega _j^{m_{ji}}=1`$ and $`_j\omega _j^{n_{jb}}=1`$. As in the example of hypersurfaces in $`𝐂\mathrm{P}^{N1}`$, one might expect that the last $`k`$ terms in (7.36) are irrelevant at low energies and the IR fixed point is described by $$\stackrel{~}{W}_{IR}=\underset{i=1}{\overset{N}{}}\underset{j=1}{\overset{N}{}}X_j^{m_{ji}}.$$ (7.37) However, as we will see, an extra condition is required for this. Intuitively, the last $`k`$ terms are irrelevant when $`d_{ab}^1t_a\mathrm{}`$ at low energies. The flow of $`t_a`$ is determined by $`\beta _a=_{i=1}^NQ_{ia}_{b=1}^kd_{ba}`$ as $`t_a(\mu )=\beta _a\mathrm{log}(\mu /\mathrm{\Lambda })`$. Then, the condition of irrelevance is $`d_{ab}^1\beta _a>0`$ or equivalently $$\underset{i,a}{}d_{ab}^1Q_{ia}>1.$$ (7.38) More precisely, let $`2q_i`$ be the axial charge of $`X_i`$ so that each of the first $`N`$ terms in (7.36) has charge $`2`$, namely $`_{j=1}^Nm_{ji}q_j=1`$. Then, from (7.33) we find $`_{i=1}^NQ_{ia}=_{b,j}q_jn_{jb}d_{ba}`$. Thus the condition (7.38) means $`_{j=1}^Nq_jn_{jb}>1`$, which is nothing but the irrelevance of $`_{j=1}^NX_j^{n_{jb}}`$. Thus, under the condition (7.38), the IR fixed point is described by the LG orbifold with superpotential (7.37). The central charge of this model is $`c/3=N2_{i,j}m_{ij}^1`$. What we have said remains true when the inequalities are relaxed to allow equalities. If an equality holds, the corresponding term in (7.36) is a marginal operator and should be kept. In the original linear sigma model, on the other hand, under the same condition we can take $`G_b`$ to be $$G_b=\underset{j=1}{\overset{N}{}}n_{jb}\underset{i=1}{\overset{N}{}}\mathrm{\Phi }_i^{m_{ji}},$$ (7.39) where the sum is over $`j`$ such that $`n_{jb}=1`$. The D-term equations for the chiral fields are expressed as $$\underset{i,a}{}d_{ab}^1Q_{ia}\|\varphi _i\|^2\|p_b\|^2=\underset{a}{}d_{ab}^1t_a.$$ (7.40) At low energies, under the condition (7.38), the right hand side flows to large negative values for all $`b`$. Under the same condition, the coefficients of $`\|\varphi _i\|^2`$ are all positive. Then, the equation (7.40) implies all $`p_b`$ are non-vanishing. The vacuum equation also requires $`_{b=1}^kp_b_iG_b=0`$. It follows from this that $`_i\mathrm{\Phi }_i^{m_{ji}}=0`$ for all $`j`$ and that $`_{j_1,b}^{}n_{j_1b}p_b_{j_2i}\mathrm{\Phi }_{j_2}^{m_{j_1j_2}}=0`$ for all $`i`$ where the sum $`^{}`$ is over such $`(j_1,b)`$ that $`m_{j_1i}=1`$. We assume that this implies $`\mathrm{\Phi }_i=0`$ for all $`i`$. (We do not attempt to prove it here. It is possible that in general an extra condition is required.) Then the gauge group $`U(1)^k`$ is broken to its subgroup $`\mathrm{\Gamma }`$ (generated by $`g_a`$ such that $`_ag_a^{d_{ba}}=1`$, $`b`$) which is a discrete subgroup since $`d_{ba}`$ is invertible, and the massless degrees of freedom are $`\mathrm{\Phi }_i`$’s only. Thus, we expect that the theory flows in the IR limit to the $`\mathrm{\Gamma }`$-orbifold of the LG model with the superpotential $$W=\underset{j,b}{}n_{jb}p_b\underset{i=1}{\overset{N}{}}\mathrm{\Phi }_i^{m_{ji}}$$ (7.41) where $`p_b`$ is the expectation value of the massive field $`p_b`$. One can relax the inequality in (7.38) to admit equality. If an equality holds, we can choose the value of $`t_a`$’s such that $`d_{ab}^1t_a`$ are all negative, and the same conclusion holds. The central charge of the model is $`c/3=N2_{i,j}m_{ij}^1`$, the same as the one for (7.37). Since the two LG orbifold models appear as the IR limit of the same theory, they must be equivalent, or mirror to each other. The equivalences of this type between LG models have been noted before . For illustration, let us present an example. We consider a complete intersection in $`𝐂\mathrm{P}^{N1}\times 𝐂\mathrm{P}^{M1}`$ ($`NM`$) defined by two equations of bi-degree $`(d_1,0)`$ and $`(1,d_2)`$ respectively. It has a non-negative first Chern class if $`Nd_1+1`$ and $`Md_2`$. The equations for the homogeneous coordinates $`S_i`$ ($`i=1,\mathrm{},N`$) and $`T_j`$ ($`j=1,\mathrm{},M`$) are $$G_1=\underset{i=1}{\overset{N}{}}S_i^{d_1},G_2=\underset{j=1}{\overset{M}{}}S_jT_j^{d_2}.$$ (7.42) Under the condition<sup>1</sup><sup>1</sup>1Note that this condition is stronger than the non-negativity of the first Chern class. $$Nd_1+M/d_2,Md_2,$$ (7.43) $`d_{ab}^1t_a`$ are (or can be chosen) large negative at low energies, and we find $`p_1`$ and $`p_2`$ are non-vanishing. Then, we can show in this case that $`_{b=1,2}p_bG_b=0`$ implies that all $`S_i`$ and $`T_j`$ must be vanishing. Thus, we find an LG orbifold description of the low energy theory where the superpotential is $$W=\underset{j=1}{\overset{M}{}}(S_j^{d_1}+S_jT_j^{d_2})+\underset{k=M+1}{\overset{N}{}}S_k^{d_1},$$ (7.44) and the orbifold group is the subgroup of $`U(1)_1\times U(1)_2`$ defined by $`g_1^{d_1}=1`$ and $`g_1g_2^{d_2}=1`$. On the other hand, the dual theory is described by the LG orbifold for $`N+M`$ twisted chiral superfields $`U_i`$ and $`V_j`$ with the twisted superpotential $$\stackrel{~}{W}=\underset{j=1}{\overset{M}{}}(U_j^{d_1}V_j+V_j^{d_2})+\underset{k=M+1}{\overset{N}{}}U_k^{d_1}+\mathrm{e}^{\frac{t_1}{d_1}\frac{t_2}{d_1d_2}}\underset{i=1}{\overset{N}{}}U_i+\mathrm{e}^{\frac{t_2}{d_2}}\underset{j=1}{\overset{M}{}}V_j,$$ (7.45) and the group acting on fields as $`V_j\gamma _jV_j,j=1,\mathrm{},M,`$ (7.46) $`U_j\eta _jU_j,j=1,\mathrm{},M,`$ (7.47) $`U_k\omega _kU_k,k=M+1,\mathrm{},N,`$ (7.48) where $`\omega _k^{d_1}=1`$, $`\gamma _j^{d_2}=1`$, $`\gamma _j\eta _j^{d_1}=1`$, and $`_k\omega _k_j\eta _j=_j\gamma _j=1`$. Under the condition (7.43), the last two terms of (7.45) are marginal or irrelevant depending on whether the equality holds or not. The two LG orbifolds must be mirror to each other. Indeed both have a central charge $`c/3=(11/d_1)(N+M2M/d_2)`$. ### 7.3 General Mirror Description As noted before, for the most general case, what we will find is that the mirror theory can be expressed as an LG theory on a non-compact Calabi-Yau manifold. To see how this works we begin with the case already discussed, i.e. hypersurfaces of degree $`d`$ in $`𝐂\mathrm{P}^{N1}`$ and provide an alternative description of the mirror. This reformulation sets the stage for the most general description which follows it. As noted before, the relevant periods (i.e. D-brane ‘masses’) are given in this case by $$\mathrm{\Pi }=dd\mathrm{\Sigma }dY_P\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\mathrm{\Sigma }\mathrm{exp}\left(\stackrel{~}{W}\right)$$ (7.49) Since $`d\mathrm{\Sigma }`$ is given by $`/Y_P`$ of the linear terms in $`\stackrel{~}{W}`$ we have $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle d\mathrm{\Sigma }\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\mathrm{d}Y_P\frac{}{Y_P}\left[\mathrm{exp}\left(\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Y_idY_Pt\right)\right)\right]\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}\mathrm{e}^{Y_P}\right)}`$ (7.50) $`=`$ $`{\displaystyle d\mathrm{\Sigma }\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\mathrm{d}Y_P\mathrm{e}^{Y_P}\mathrm{exp}\left(\mathrm{\Sigma }\left(\underset{i=1}{\overset{N}{}}Y_idY_Pt\right)\right)\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}\mathrm{e}^{Y_P}\right)}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\mathrm{d}Y_P\mathrm{e}^{Y_P}\delta \left(\underset{i=1}{\overset{N}{}}Y_idY_Pt\right)\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}\mathrm{e}^{Y_P}\right)}`$ We make the following change of variables $`\mathrm{e}^{Y_P}=\stackrel{~}{P},`$ (7.51) $`\mathrm{e}^{Y_i}=\stackrel{~}{P}U_i,\text{for}i=1,\mathrm{},d,`$ (7.52) $`\mathrm{e}^{Y_j}=U_j,\text{for}j=d=1,\mathrm{},N.`$ (7.53) Then, $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\frac{\mathrm{d}U_i}{U_i}\mathrm{d}\stackrel{~}{P}\delta \left(\mathrm{log}\left(\underset{i=1}{\overset{N}{}}U_i\right)+t\right)\mathrm{exp}\left(\stackrel{~}{P}\left(\underset{i=1}{\overset{d}{}}U_i+1\right)\underset{i=d+1}{\overset{N}{}}U_i\right)}`$ (7.54) $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\frac{\mathrm{d}U_i}{U_i}\delta \left(\mathrm{log}\left(\underset{i=1}{\overset{N}{}}U_i\right)+t\right)\delta \left(\underset{i=1}{\overset{d}{}}U_i+1\right)\mathrm{exp}\left(\underset{i=d+1}{\overset{N}{}}U_i\right)}.`$ Thus we have obtained a submanifold $`\stackrel{~}{M}^{}`$ of $`(𝐂^\times )^N`$ defined by $`{\displaystyle \underset{i=1}{\overset{N}{}}}U_i=\mathrm{e}^t,`$ (7.55) $`{\displaystyle \underset{i=1}{\overset{d}{}}}U_i+1=0.`$ (7.56) This is a non-compact manifold of dimension $`N2`$. The expression (7.54) is identical to the period of an LG model on $`\stackrel{~}{M}^{}`$ with superpotential $$W_{\stackrel{~}{M}^{}}=\underset{i=d+1}{\overset{N}{}}U_i.$$ (7.57) This model is the mirror of the sigma model on $`M`$, at least when twisted to topological field theory. The special case is the case $`d=N`$ in which $`M`$ is a compact Calabi-Yau manifold. In this case, the superpotential (7.57) is trivial and the mirror is simply the non-linear sigma model on $`\stackrel{~}{M}^{}`$. The mirror manifold $`\stackrel{~}{M}^{}`$ is actually an open subset of a Calabi-Yau manifold $`\stackrel{~}{M}`$ which is familiar to us. That is, the orbifold of the hypersurface in $`𝐂\mathrm{P}^{N1}`$ $$G(Z_1,\mathrm{},Z_N)=Z_1^N+\mathrm{}+Z_N^N+\mathrm{e}^{t/N}Z_1\mathrm{}Z_N=0,$$ (7.58) by the $`(𝐙_N)^{N2}`$ action given by $$Z_i\gamma _iZ_i,\gamma _i^N=1,\gamma _1\mathrm{}\gamma _N=1.$$ (7.59) To see this, we note that $$U_i=\mathrm{e}^{t/N}\frac{Z_i^N}{Z_1\mathrm{}Z_N}$$ (7.60) is invariant under the $`𝐂^\times \times (𝐙_N)^{N2}`$ action and solves the first equation (7.55). The second equation (7.56) becomes $`G(Z_i)=0`$. If $`Z_i`$ and $`Z_i^{}`$ yields the same $`U_i`$, it is easy to see that $`Z_i^N=Z_i^N`$ and $`Z_1\mathrm{}Z_N=Z_1^{}\mathrm{}Z_N^{}`$ modulo $`𝐂^\times `$ action. Then, this means $`Z_i=Z_i^{}`$ modulo the $`𝐂^\times \times (𝐙_N)^{N2}`$ action. Thus, the map from $`[Z_i]`$ to $`U_i`$ is one to one. Under this identification, we have $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{vol}(𝐂^\times )}\underset{i=1}{\overset{N}{}}\frac{\mathrm{d}Z_i}{Z_i}\delta \left(\frac{G(Z_1,\mathrm{},Z_N)}{Z_1\mathrm{}Z_N}\right)}`$ (7.61) $`=`$ $`{\displaystyle \underset{i=1}{\overset{N1}{}}\mathrm{d}Z_i\delta \left(G(Z_1,\mathrm{},Z_{N1},1)\right)}`$ $`=`$ $`{\displaystyle \left(\underset{i=1}{\overset{N2}{}}\mathrm{d}Z_i/\frac{G\|_{Z_N=1}}{Z_{N1}}\|_{G=0}\right)}={\displaystyle \mathrm{\Omega }},`$ which is the period of the holomorphic differential of the Calabi-Yau manifold $`\stackrel{~}{M}`$. Thus, we have shown that the A-twisted sigma model on $`M`$ is equivalent to the B-twisted sigma model on $`\stackrel{~}{M}`$. This tempts us to propose that the sigma model on $`M`$ and the LG model on $`\stackrel{~}{M}^{}`$ with the superpotential $`W_{\stackrel{~}{M}^{}}`$ are mirror to each other as (2,2) quantum field theories, not just as topological theories. This is certainly true for $`d=1`$ case with $`M=𝐂\mathrm{P}^{N2}`$ where $`\stackrel{~}{M}^{}`$ is the algebraic torus $`(𝐂^\times )^{N2}`$ and $`W_{\stackrel{~}{M}^{}}`$ is the affine Toda superpotential. However, for $`d2`$ we must partially compactify $`\stackrel{~}{M}^{}`$ for the model to be the mirror of the sigma model on $`M`$. The reason is that the superpotential $`W_{\stackrel{~}{M}^{}}`$ for $`d>1`$ has a run-away direction which yields a continuous spectrum, a property that we do not expect for a sigma model on a compact smooth manifold $`M`$. The typical case is $`d=N`$; the sigma model on a non-compact manifold has a continuous spectrum. As we have seen above, $`\stackrel{~}{M}^{}`$ can indeed be compactified to a compact CY manifold $`\stackrel{~}{M}`$ and we can claim that the sigma model on $`M`$ is mirror to the sigma model on $`\stackrel{~}{M}`$. For $`2d<N`$, we must compactify only partially since a non-trivial superpotential is not allowed on a compact complex manifold. In fact, this partial compactification can be found as a simple generalization of the $`d=N`$ case. Let us solve the first equation (7.55) for $`U_i`$’s as follows $`U_i=\mathrm{e}^{t/d}{\displaystyle \frac{Z_i^d}{Z_1\mathrm{}Z_N}},i=1,\mathrm{},d,`$ (7.62) $`U_j=Z_j^d,j=d+1,\mathrm{},N.`$ (7.63) As in the $`d=N`$ case one can see that the map from $`Z_i`$ to $`U_i`$ is one to one modulo the $`𝐂^\times \times (𝐙_d)^{N2}`$ action given by $`Z_i\lambda \gamma _iZ_i,i=1,\mathrm{},d,`$ (7.64) $`Z_j\gamma _jZ_j,j=d+1,\mathrm{},N,`$ (7.65) where $`\lambda 𝐂^\times `$ and $`\gamma _i^d=\gamma _j^d=1`$ and $`\gamma _1\mathrm{}\gamma _N=1`$. The second equation (7.56) is then expressed as $$Z_1^d+\mathrm{}+Z_d^d+\mathrm{e}^{t/d}Z_1\mathrm{}Z_dZ_{d+1}\mathrm{}Z_N=0.$$ (7.66) This is the equation for a Calabi-Yau hypersurface in $`𝐂\mathrm{P}^{d1}`$ with the $`\psi `$ parameter $`\psi =\mathrm{e}^{t/d}(Z_{d+1}\mathrm{}Z_N)`$. Thus, the manifold $`\stackrel{~}{M}^{}`$ is partially compactified to a manifold $`\stackrel{~}{M}`$ which is the $`(𝐙_d)^{N2}`$ quotient of the total space of the family of CY hypersurface in $`𝐂\mathrm{P}^{d1}`$ parametrized by $`𝐂^{Nd}`$ via $`\psi =\mathrm{e}^{t/d}(Z_{d+1}\mathrm{}Z_N)`$. Now the superpotential $`W_{\stackrel{~}{M}^{}}`$ (7.57) on $`\stackrel{~}{M}^{}`$ extends to $`\stackrel{~}{M}`$ as $$W_{\stackrel{~}{M}}=Z_{d+1}^d+\mathrm{}+Z_N^d.$$ (7.67) Repeating what we have done in the $`d=N`$ case, we can see that the period is expressed as $$\mathrm{\Pi }=\mathrm{\Omega }_{d2}\mathrm{d}Z_{d+1}\mathrm{}\mathrm{d}Z_N\mathrm{exp}\left(W_{\stackrel{~}{M}}\right),$$ (7.68) where $`\mathrm{\Omega }_{d2}`$ is the holomorphic $`(d2)`$-form of the CY hypersurface in $`𝐂\mathrm{P}^{d1}`$. Now, the superpotential $`W_{\stackrel{~}{M}}`$ on $`\stackrel{~}{M}`$ has no run-away direction; $`\stackrel{~}{M}`$ includes the limiting points of the run-away direction in $`\stackrel{~}{M}^{}`$. In particular, we expect that the theory has a discrete spectrum. Thus, we claim that the sigma model on $`M`$ is mirror to the LG model on $`\stackrel{~}{M}`$ with the superpotential (7.67). The superpotential (7.67) has $`(Nd)`$ non-degenerate critical points at $`{\displaystyle \frac{Z_i^d}{Z_1\mathrm{}Z_N}}={\displaystyle \frac{\mathrm{e}^{t/d}}{d}},i=1,\mathrm{},d,`$ (7.69) $`Z_j^d=U,j=d+1,\mathrm{},N;U^{Nd}=(d)^d\mathrm{e}^t,`$ (7.70) and a critical manifold at $$Z_{d+1}=\mathrm{}=Z_N=0,$$ (7.71) which is the CY hypersurface $`_{i=1}^dZ_i^d=0`$ of $`𝐂\mathrm{P}^{d1}`$. For $`d>2`$, this critical CY manifold has dimension $`>0`$ and also the superpotential (7.67) is degenerate there. Thus, we expect that the theory for $`d>2`$ flows in the IR limit to a non-trivial fixed point. This must be equivalent to the non-trivial fixed point studied in section 7.1. For $`d=2`$, the critical CY manifold is actually a point and the superpotential is non-degenerate. Thus, we expect that the theory has a mass gap for $`d=2`$. However, because of the orbifolding, the critical point (7.71) may correspond to multiple vacua. From the result of section 7.1.1 the actual number of vacua there is $`2`$ for even $`N`$ and $`1`$ for odd $`N`$. ### Complete Intersection in Toric Variety Let $`X`$ be the toric variety defined by the charge matrix $`Q_{ia}`$ ($`i=1,\mathrm{},N`$, $`a=1,\mathrm{},k`$). We consider the submanifold $`M`$ of $`X`$ defined by the equations $$G_\beta =0,\beta =1,\mathrm{},l,$$ (7.72) where $`G_\beta `$ are polynomials of $`\mathrm{\Phi }_i`$ of charge $`d_{\beta a}`$ for the $`a`$-th $`U(1)`$ gauge group. The sigma model on $`M`$ can be realized as the linear sigma model of gauge group $`U(1)^k`$ with chiral superfields $`\mathrm{\Phi }_i`$ of charge $`Q_{ia}`$ and $`P_\beta `$ of charge $`d_{\beta a}`$ which has a superpotential $$W=\underset{\beta =1}{\overset{l}{}}P_\beta G_\beta (\mathrm{\Phi }).$$ (7.73) The theory without this superpotential is the same as the sigma model on a non-compact toric variety $`V`$ (defined by charge matrix $`(Q_{ia},d_{\beta a})`$) and has the dual description in terms of the twisted chiral superfield $`\mathrm{\Sigma }_a`$, $`Y_i`$ and $`Y_{P_\beta }`$ where $`\mathrm{\Sigma }_a`$ is the field strength of the $`a`$-th gauge group and $`Y_i`$ and $`Y_{P_\beta }`$ are the dual variables of $`\mathrm{\Phi }_i`$ and $`P_\beta `$. The dual theory has the twisted superpotential $$\stackrel{~}{W}=\underset{a=1}{\overset{k}{}}\mathrm{\Sigma }_a\left(\underset{i=1}{\overset{N}{}}Q_{ia}Y_i\underset{\beta =1}{\overset{l}{}}d_{\beta a}Y_{P_\beta }t_a\right)+\underset{i=1}{\overset{N}{}}\mathrm{e}^{Y_i}+\underset{\beta =1}{\overset{l}{}}\mathrm{e}^{Y_{P_\beta }}.$$ (7.74) As noted before, the period integral for the compact theory is given by $$\mathrm{\Pi }=\underset{a=1}{\overset{k}{}}\mathrm{d}\mathrm{\Sigma }_a\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\underset{\beta =1}{\overset{l}{}}\mathrm{d}Y_{P_\beta }\delta _1\mathrm{}\delta _l\mathrm{exp}\left(\stackrel{~}{W}\right),$$ (7.75) where $$\delta _\beta =\underset{a=1}{\overset{k}{}}d_{\beta a}\mathrm{\Sigma }_a.$$ (7.76) We note that this can be expressed as $$\delta _\beta =\frac{}{Y_{P_\beta }}\underset{a=1}{\overset{k}{}}\mathrm{\Sigma }_a\left(\underset{i=1}{\overset{N}{}}Q_{ia}Y_i\underset{\beta =1}{\overset{l}{}}d_{\beta a}Y_{P_\beta }t_a\right).$$ (7.77) Then, via partial integration we obtain $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \underset{a=1}{\overset{k}{}}\mathrm{d}\mathrm{\Sigma }_a\underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\underset{\beta =1}{\overset{l}{}}\mathrm{d}Y_{P_\beta }\mathrm{e}^{Y_{P_1}}\mathrm{}\mathrm{e}^{Y_{P_l}}}`$ (7.78) $`\times \mathrm{exp}\left({\displaystyle \underset{a=1}{\overset{k}{}}}\mathrm{\Sigma }_a\left({\displaystyle \underset{i=1}{\overset{N}{}}}Q_{ia}Y_i{\displaystyle \underset{\beta =1}{\overset{l}{}}}d_{\beta a}Y_{P_\beta }t_a\right)\right)\mathrm{exp}\left({\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{e}^{Y_i}{\displaystyle \underset{\beta =1}{\overset{l}{}}}\mathrm{e}^{Y_{P_\beta }}\right)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\mathrm{d}Y_i\underset{\beta =1}{\overset{l}{}}\mathrm{e}^{Y_{P_\beta }}\mathrm{d}Y_{P_\beta }\underset{a=1}{\overset{k}{}}\delta \left(\underset{i=1}{\overset{N}{}}Q_{ia}Y_i\underset{\beta =1}{\overset{l}{}}d_{\beta a}Y_{P_\beta }t_a\right)}`$ $`\times \mathrm{exp}\left({\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{e}^{Y_i}{\displaystyle \underset{\beta =1}{\overset{l}{}}}\mathrm{e}^{Y_{P_\beta }}\right).`$ This is the expression for the BPS mass for the most general toric complete intersection. Now let us consider the case where one can find $`n_\beta ^i=0`$ or $`1`$ such that $$\underset{i=1}{\overset{N}{}}n_\beta ^iQ_{ia}=d_{\beta a}.$$ (7.79) If we make the following change of variables $`\mathrm{e}^{Y_{P_\beta }}=\stackrel{~}{P}_\beta ,`$ (7.80) $`\mathrm{e}^{Y_i}=U_i{\displaystyle \underset{\beta =1}{\overset{l}{}}}\stackrel{~}{P}_\beta ^{n_\beta ^i},`$ (7.81) the period integral is expressed as $`\mathrm{\Pi }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\frac{\mathrm{d}U_i}{U_i}\underset{\beta =1}{\overset{l}{}}\mathrm{d}\stackrel{~}{P}_\beta \underset{a=1}{\overset{k}{}}\delta \left(\mathrm{log}\left(\underset{i=1}{\overset{N}{}}U_i^{Q_{ia}}\right)+t_a\right)\mathrm{exp}\left(\underset{\beta =1}{\overset{l}{}}\stackrel{~}{P}_\beta \left(\underset{n_\beta ^i=1}{}U_i+1\right)\underset{n_\beta ^j=0,\beta }{}U_j\right)}`$ (7.82) $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}\frac{\mathrm{d}U_i}{U_i}\underset{a=1}{\overset{k}{}}\delta \left(\mathrm{log}\left(\underset{i=1}{\overset{N}{}}U_i^{Q_{ia}}\right)+t_a\right)\underset{\beta =1}{\overset{l}{}}\delta \left(\underset{n_\beta ^i=1}{}U_i+1\right)\times \mathrm{exp}\left(\underset{n_\beta ^j=0,\beta }{}U_j\right)}`$ Thus, we have obtained a submanifold $`\stackrel{~}{M}^{}`$ of $`(𝐂^\times )^N`$ defined by $`{\displaystyle \underset{i=1}{\overset{N}{}}}U_i^{Q_{ia}}=\mathrm{e}^{t_a},`$ (7.83) $`{\displaystyle \underset{n_\beta ^i=1}{}}U_i+1=0,`$ (7.84) where the sum in (7.84) is over $`i`$ such that $`n_\beta ^i=1`$. This is a non-compact manifold of dimension $`(Nkl)`$ which is the same as the dimension of $`M`$. The period (7.82) is identical to the period of an LG model on $`\stackrel{~}{M}^{}`$ with superpotential $$W_{\stackrel{~}{M}^{}}=\underset{n_\beta ^j=0,\beta }{}U_j,$$ (7.85) where the sum here is over $`j`$ such that $`n_\beta ^j=0`$ for all $`\beta `$. Thus, we have shown that the mirror of the sigma model on $`M`$ is given by this LG model on $`\stackrel{~}{M}^{}`$, at least when twisted to topological field theory. The expression (7.82) of the period on the mirror manifold (7.83)-(7.84) with the superpotential (7.85) is equivalent to the one that appears in . ## 8 Directions for Future Work In this paper we have seen how mirror symmetry, formulated as a duality between pairs of 2d QFT’s can be proven using rather simple physical ideas. Not only have we recovered the known formulation for mirror symmetry, but we have also generalized it to classes which were not known before. In particular, we have shown that a theory and the mirror will give rise to the same “BPS mass” for the D-branes. It is natural to ask to what extent one can derive dualities between QFT’s in higher dimensions. One idea along this line is their reduction to 2 dimensions. For example, it was shown in how certain results about 4 dimensional $`N=2`$ Yang-Mills theory compactified on a Riemann surface to 2 dimensions can be related to sigma models, and computed via mirror symmetry. Given that we have now understood how mirror symmetry can be derived, we can go back and ask to what extent all the quantum corrections of $`N=2`$ theories can be understood using results about 2d QFT’s. It would not be surprising if all the BPS data , which is all we know at the present about the $`N=2`$ theories in 4 dimensions, can be derived in this way. This in fact fits with the idea that in supersymmetric theories certain exact results can be computed rather easily. For example computation of Witten’s index can be reduced to an ordinary integral. It is thus conceivable that the computations related to exact BPS aspects for $`N=2`$ theories in 4 dimensions can be related to derivable facts about 2 dimensional quantum field theories. This program, we believe, is an interesting one to investigate. At least in one case it already seems to work: Consider pure $`N=1`$ supersymmetric Yang-Mills theory in 4 dimensions. It is known that upon compactification to 3 dimensions one obtains an effective theory captured by a superpotential of affine Toda type . The same superpotential survives in the compactification down to 2 dimensions. This can now be derived in our setup as follows: Consider compactification of the four dimensional theory on $`T^2`$, down to two dimensions. The effective theory in two dimensions will be a sigma model whose target space is the moduli space of flat connection of the corresponding group on $`T^2`$. These spaces have been studied and is found to be the weighted projective space $`𝐂\mathrm{WP}_{(1,g_1,\mathrm{},g_r)}^r`$ where $`r`$ is the rank of the gauge group and $`g_1,\mathrm{},g_r`$ are its coroot integers. The $`(2,2)`$ sigma model on this space is realized as a $`U(1)`$ gauge theory with chiral matter fields of charge $`(1,g_1,\mathrm{},g_r)`$. Using the mirror symmetry results of this paper, we obtain the mirror affine Toda potential $$\stackrel{~}{W}=\mathrm{\Lambda }\left(\mathrm{e}^{\mathrm{\Theta }_1}+\mathrm{}+\mathrm{e}^{\mathrm{\Theta }_r}+\mathrm{e}^{_{i=1}^rg_i\mathrm{\Theta }_i}\right),$$ (8.1) in agreement with . It would be interesting to extend these results to other $`N=1`$ theories in four dimensions<sup>1</sup><sup>1</sup>1For another relation between moduli of flat connections on $`T^2`$ and mirror symmetry, see .. Another aspect of our work which may find further applications is the issue of the meeting of Higgs and Coulomb branch in 2 dimensional gauge theories. This issue has been studied in . The dual of gauge theory coupled to matter that we have found in this paper seems potentially useful for studying such questions. In particular the dual of the matter fields (whose vevs mark the Higgs branch) and the scalar in the vector multiplet (whose vev marks the Coulomb branch) are in the same type of multiplet and the superpotential we have found captures certain aspects of the gauge dynamics involving these fields. Another application of our work may be to the large $`N`$ behavior of Chern-Simons theory on $`S^3`$, which has been conjectured to be dual to topological string on $`𝒪(1)+𝒪(1)`$ over $`𝐂\mathrm{P}^1`$ . In appendix A we have presented a conjectured generalization of the duality in (2,2) supersymmetric gauge theories in 2 dimensions for non-abelian groups, motivated by considering a generic point on the Coulomb branch of these theories. It would be interesting to check the validity of this conjecture. This conjecture leads to the computation of BPS structure for complete intersections defined on arbitrary flag varieties. Certain aspects of our work may also be relevant for the coupling of topological sigma model to topological gravity. In fact, it is tempting to conjecture that some matrix version of the toda theories we have found would be relevant for that reformulation. This would be interesting to study further. ## Acknowledgement We would like to thank K. Intriligator, A. Iqbal, S. Katz, A. Klemm, B. Lian, D. Morrison, H. Ooguri, R. Pandharipande, R. Plesser, N. Seiberg, R. Thomas and S.T. Yau for valuable discussions. K.H. would like to thank the Institute for Advanced Study, Princeton and Duke University where a part of this work was carried out. The research of K.H. is supported in part by NSF-DMS 9709694. The research of C.V. is supported in part by NSF grant PHY-98-02709. ## Appendix ## A Non-Abelian Gauge Theories and Complete Intersections in Grassmannians— A Conjecture In this paper we have been mainly dealing with linear sigma models realized as complete intersection in toric varieties. We have used abelian gauge theories in analyzing them. It is natural to ask if we can compute analogous quantities for the case of complete intersections in Grassmannians. In fact this class can also be realized as gauge theories as well , but in this case it will involve a $`U(N)`$ gauge theory coupled to matter. So we would need do know the analog of the dual formulation for this gauge system. Our argument for abelian gauge theory is not applicable to non-abelain gauge theories. Here we attempt to make a conjecture for what the analog dual is, which we motivate using the results already obtained for the product of $`U(1)`$’s. We also make some checks for the validity of the conjecture. We consider the $`U(N)`$ gauge theory coupled to matter in a (not necessarily irreducible) representation $``$.<sup>1</sup><sup>1</sup>1 One can generalize this conjecture to arbitrary groups. The $`U(N)`$ gauge supermultiplet contains a complex scalar in the adjoint representation of $`U(N)`$. The “generic”<sup>2</sup><sup>2</sup>2 This is why we cannot prove our conjecture: The places where some $`\mathrm{\Sigma }_i`$ and $`\mathrm{\Sigma }_j`$ are equal yields a non-abelian unbroken group and we are assuming that this does not cause any problems for our dualization. configuration which survives in the infrared is the one corresponding to the complex scalar given by a diagonal matrix. This is the generic point on the Coulomb branch of the theory. Let $`\mathrm{\Sigma }_i`$ denote the fields corresponding to these diagonal elements. These are well defined up to permutation. In other words the invariant fields involve symmetric polynomials in $`\mathrm{\Sigma }_i`$. In the generic configuration for $`\mathrm{\Sigma }_i`$ the theory becomes a $`U(1)^N/S_N`$ gauge theory. So we could then apply our results for the product of $`U(1)`$ gauge theories to obtain the dual, modulo taking into account the permutation action on the groups. There will be $`dim`$ fields $`Y^\alpha `$ obtained by dualizing the matter field in representation $`R`$. Each $`Y_\alpha `$ corresponds to a weight of $`U(N)`$ lattice, and we can associate the corresponding $`U(1)^N`$ charges $`Q_i^\alpha `$ to them. Consider the superpotential $$W=\underset{i}{}\mathrm{\Sigma }_i(Q_i^\alpha Y^\alpha t)+\underset{\alpha }{}e^{Y_\alpha }$$ (A.1) where $`t`$ denotes the FI parameter for the $`U(N)`$ gauge theory. This is almost what we propose as the dual, except that we now have to consider Weyl invariant (i.e. $`S_N`$ invariant) combination of fields. The $`S_N`$ acts on $`\mathrm{\Sigma }_i`$ by permutation, as already noted. It acts on $`Y^\alpha `$ by the permutation induced on the weights of the representation $``$ by the action of the Weyl group. Clearly the above action is invariant under the Weyl group action. We then consider the theory given by $$W_{invariant}=W//S_N$$ (A.2) where by this we mean that the fundamental fields of the theory are to be written in terms of the $`S_N`$ invariant combinations (note that this is not the same as orbifolding the theory). Concretely what this would mean in computing the periods (i.e. D-brane masses) is as follows: We will consider integrals of the form $$\underset{i}{}d\mathrm{\Sigma }_i\underset{\alpha }{}dY^\alpha \underset{i<j}{}(\mathrm{\Sigma }_i\mathrm{\Sigma }_j)e^W$$ (A.3) The insertion of $`\mathrm{\Delta }=_{i<j}(\mathrm{\Sigma }_i\mathrm{\Sigma }_j)`$ is to make the measure correspond to the symmetric measure, and is reminiscent of the $`\delta `$ insertion discussed in the context of getting hypersurfaces from non-compact toric varieties. Note that this also agrees with the dimension count for the space. Our proposed dual has infrared degrees of freedom which is too big: $`dimN`$ whereas it should have had $`dimN^2`$. Insertion of $`\mathrm{\Delta }^2`$ in the correlation functions (which is equivalent to the insertion of $`\mathrm{\Delta }`$ in the periods) changes the dimension count by $`N(N1)`$, and makes up for the discrepancy. This is very similar to how the insertion of $`\delta `$’s in the non-compact toric case was used to embed the compact cohomology ring in the non-compact one. The above periods can be easily computed from the corresponding non-compact toric varieties: The insertion of $`\mathrm{\Delta }`$ is the same as the action of $`_{i<j}(/t_i/t_j)`$ acting on the periods of the corresponding non-compact toric case (and substituting $`t_i=t_j=t`$ at the end). Similarly the case of hypersurfaces in Grassmannians would be obtained from the non-compact version of bundles over Grassmannians by inclusions of extra insertions similar to $`\delta `$. Also this conjecture naturally extends to the case of flag manifolds and complete intersections in them. The above conjecture would be interesting to verify. For the case of Grassmannian itself, this conjecture gives the correct ring structure , which was one of the motivations for the above conjecture. It would also be interesting to connect the above formulation with the corresponding LG model proposed for the Grassmannian which was further generalized in for other homogeneous spaces and complete intersections therein. ## B Supersymmetry transformation We record here the supersymmetry transformation of the vector and chiral multiplet fields (in the Wess-Zumino gauge) . This is obtained by dimensional reduction of the formulae in for $`N=1`$ supersymmetry transformation in $`3+1`$ dimensions. The reduction is in $`x^1,x^2`$ directions and the scalars in the vector multiplet is defined as $`\sigma =(v_1iv_2)`$ and $`\overline{\sigma }=(v_1+iv_2)`$. The time coordinate is still $`x^0`$ but we rename the spacial coordinate $`x^3`$ as $`x^1`$. (The normalization of vector multiplet fields used in this paper differs from the one in by factors of $`\sqrt{2}`$: if we denote the latter as $`\mathrm{\Sigma }^{}`$, $`\sigma ^{}`$, etc, the relations are $`\mathrm{\Sigma }=\sqrt{2}\mathrm{\Sigma }^{}`$, $`\sigma =\sqrt{2}\sigma ^{}`$, $`\lambda _\pm =\sqrt{2}\lambda _\pm ^{}`$, $`D=D^{}`$ and $`F_{01}=F_{01}^{}`$.) The four supersymmetry generators are combined as $$\delta =ϵ^\alpha Q_\alpha +\overline{ϵ}_\alpha \overline{Q}^\alpha =ϵ_+Q_{}ϵ_{}Q_+\overline{ϵ}_+\overline{Q}_{}+\overline{ϵ}_{}\overline{Q}_+,$$ (B.1) where $`ϵ^\pm `$ and $`\overline{ϵ}^\pm `$ are anti-commuting spinorial parameters ($`ϵ^{}=\pm ϵ_\pm `$ etc). The transformation of the vector multiplet fields is $`\delta v_\pm =\sqrt{2}i\overline{ϵ}_\pm \lambda _\pm +\sqrt{2}iϵ_\pm \overline{\lambda }_\pm ,`$ $`\delta \sigma =\sqrt{2}i\overline{ϵ}_+\lambda _{}\sqrt{2}iϵ_{}\overline{\lambda }_+,`$ $`\delta \overline{\sigma }=\sqrt{2}iϵ_+\overline{\lambda }_{}\sqrt{2}i\overline{ϵ}_{}\lambda _+,`$ $`\delta D={\displaystyle \frac{1}{\sqrt{2}}}(\overline{ϵ}_+D_{}\lambda _+\overline{ϵ}_{}D_+\lambda _{}+ϵ_+D_{}\overline{\lambda }_++ϵ_{}D_+\overline{\lambda }_{}.`$ $`.+ϵ_+[\sigma ,\overline{\lambda }_{}]+ϵ_{}[\overline{\sigma },\overline{\lambda }_+]\overline{ϵ}_{}[\sigma ,\lambda _+]\overline{ϵ}_+[\overline{\sigma },\lambda _{}]),`$ (B.2) $`\delta \lambda _+=\sqrt{2}iϵ_+(D+iF_{01}+{\displaystyle \frac{i}{2}}[\sigma ,\overline{\sigma }])+\sqrt{2}ϵ_{}D_+\overline{\sigma },`$ $`\delta \lambda _{}=\sqrt{2}iϵ_{}(DiF_{01}{\displaystyle \frac{i}{2}}[\sigma ,\overline{\sigma }])+\sqrt{2}ϵ_+D_{}\sigma ,`$ $`\delta \overline{\lambda }_+=\sqrt{2}i\overline{ϵ}_+(DiF_{01}+{\displaystyle \frac{i}{2}}[\sigma ,\overline{\sigma }])+\sqrt{2}\overline{ϵ}_{}D_+\sigma ,`$ $`\delta \overline{\lambda }_{}=\sqrt{2}i\overline{ϵ}_{}(D+iF_{01}{\displaystyle \frac{i}{2}}[\sigma ,\overline{\sigma }])+\sqrt{2}\overline{ϵ}_+D_{}\overline{\sigma },`$ where $`v_\pm =v_0\pm v_1`$, $`D_\pm =D_0\pm D_1`$. The transformation of the charged chiral multiplet fields is given by $`\delta \varphi =\sqrt{2}ϵ_+\psi _{}\sqrt{2}ϵ_{}\psi _+,`$ $`\delta \psi _+=\sqrt{2}i\overline{ϵ}_{}D_+\varphi +\sqrt{2}ϵ_+F\sqrt{2}\overline{ϵ}_+\overline{\sigma }\varphi ,`$ $`\delta \psi _{}=\sqrt{2}i\overline{ϵ}_+D_{}\varphi +\sqrt{2}ϵ_{}F+\sqrt{2}\overline{ϵ}_{}\sigma \varphi ,`$ (B.3) $`\delta F=\sqrt{2}i\overline{ϵ}_+D_{}\psi _+\sqrt{2}i\overline{ϵ}_{}D_+\psi _{}`$ $`+\sqrt{2}(\overline{ϵ}_+\overline{\sigma }\psi _{}+\overline{ϵ}_{}\sigma \psi _+)+\sqrt{2}i(\overline{ϵ}_{}\overline{\lambda }_+\overline{ϵ}_+\overline{\lambda }_{})\varphi .`$ From (B.2) it is clear that turning on an expectation value of the scalar component of the vector multiplet $`\sigma =\stackrel{~}{m}`$ and freezing the fluctuation of the entire vector multiplet (by setting the coupling zero) does not break any supersymmetry provided $`\stackrel{~}{m}`$ and its hermitian conjugate $`\overline{\stackrel{~}{m}}`$ commute with each other $$[\stackrel{~}{m},\overline{\stackrel{~}{m}}]=0.$$ (B.4) Thus, twisted mass for an abelian subgroup of the flavor symmetry preserves $`(2,2)`$ supersymmetry. The same thing can be said also for holomorphic isometry of a Kahler manifold (see the formulae (28)-(30) in ).
warning/0002/physics0002053.html	ar5iv	text	# SAMPLING PROPERTIES OF THE SPECTRUM AND COHERENCY OF SEQUENCES OF ACTION POTENTIALS ## 1 Introduction The study of spike trains is of central importance to electrophysiology. Often changes in the mean firing rate are studied but there is increasing interest in characterising the temporal structure of spike trains, and the relationships between spike trains, more completely \[Gray et al., 1989, Gerstein et al., 1985, Abeles et al., 1983\]. A natural extension to estimating the rate of neuronal firing is to estimate the autocorrelation and the cross-correlation functions<sup>1</sup><sup>1</sup>1Definitions of these quantities will be given in section 2.3.. This paper will discuss the frequency domain counterparts of these quantities. Auto- and cross-correlations correspond to spectra and cross spectra respectively. The coherency, which is the normalised cross spectrum, does not in general have a simple time domain counterpart. The frequency domain has several advantages over the time domain. Firstly often subtle structure can be detected with the frequency domain estimators which is difficult to observe with the time domain estimators. Secondly, the time domain quantities have problems which are associated with sensitivity of the estimators to weak non-stationarity and the non-local nature of the error bars \[Brody, 1998\]. These problems are greatly reduced in the frequency domain. Thirdly, reasonably accurate confidence intervals may be placed on estimates of the second order properties in the frequency domain which permits the statistical significance of features to be assessed. Fourthly, the coherency provides a normalised measure of correlations between time series, in contrast with time domain cross-correlations which are not normalisable by any simple means. This paper begins by reviewing the population spectrum and coherency for point processes and motivating their use by describing some example applications. Next direct, lag window and multitaper estimators of the spectrum and coherency are presented. The concept of degrees of freedom is introduced and used to obtain large sample error bars for the estimators. Many elements of the work discussed in the review section of this paper can be found in the references \[Percival and Walden, 1993, Cox and Lewis, 1966, Brillinger, 1978, Bartlett, 1966\]. Most of the material in these references is targeted at either spectral analysis of continuous processes or at the analysis of point processes but with less emphasis on spectral analysis. Building on this framework corrections, based on a specific model, will be given for finite sample sizes. These corrections are cast in terms of a reduction in the degrees of freedom of the estimators. For a homogeneous Poisson process the modified degrees of freedom is the harmonic sum of the the asymptotic degrees of freedom and twice the number of spikes used to construct the estimate. Modifications to this basic result are given for structured spectra and tapered data. A section is included on the treatment of point process spectra which contain lines. A statistical test for the presence of a line in a background of coloured noise is given, and the method for removal of such a line described. An example application to periodic stimulation is given. ## 2 Population measures and their interpretation ### 2.1 Counting representation of a spike train A spike train may be regarded as a point process. If the spike shapes are neglected, it is completely specified by a series of spike times $`\{t_i\}`$ and the start and end points of the recording interval $`[0,T]`$. It is convenient to introduce some notation which enables the subsequent formulae to be written in a compact form \[Brillinger, 1978\]. The counting process $`N(t)`$ is defined as the number of spikes which occur between the start of the interval $`(t=0)`$ and time $`t`$. The counting process has the property that the area beneath it grows as $`t`$ becomes larger. This is undesirable because it leads to an interval dependent peak at low frequencies in the spectrum. To avoid this a process $`\overline{N}(t)=N(t)\lambda t`$, where $`\lambda `$ is the mean rate, which has zero mean may be constructed. Note that $`d\overline{N}(t)=\overline{N}(t+dt)\overline{N}(t)`$ which is either $`1\lambda dt`$ or $`\lambda dt`$ depending on whether or not there is a spike in the interval $`dt`$. Thus $`d\overline{N}(t)/dt`$ is a series of delta functions<sup>2</sup><sup>2</sup>2A delta function is a generalized function. It has an area of one beneath it but has zero width and therefore infinite height. with the mean rate subtracted. Figure 1 illustrates the relationship between $`N(t)`$, $`\overline{N}(t)`$, and $`d\overline{N}(t)/dt`$. ### 2.2 Stationarity It will be assumed in what follows that the spike trains are second order stationary. This means that their first and second moments do not depend on the absolute time. In many electrophysiology experiments this is not the case. In awake behaving studies the animal is often trained to perform a highly structured task. Nevertheless it may still be the case that over an appropriately chosen short time window, the statistical properties are changing slowly enough for reasonable estimates of the spectrum and coherency to be obtained. As an example, neurons in primate parietal area PRR exhibit what is known as memory activity during a delayed reach task \[Snyder et al., 1997\]. The mean firing rate of these neurons varies considerably during the task but during the memory period is roughly constant. The assumption of stationarity during the memory period is equivalent to the intuitive notion that there is nothing special about 0.75s into the memory period as compared to say 0.5s. Second order stationarity implies that the mean firing rate ($`\lambda `$) is independent of time and additionally that the autocovariance depends only on the lag ($`\tau `$) and not on the absolute time. ### 2.3 Definitions Equations 1 \- 4 give the first and second order moments of a single spike train for a stationary process. The spectrum $`S(f)`$ is the Fourier transform of the autocovariance function $`(\mu (\tau )+\lambda \delta (\tau ))`$. $`{\displaystyle \frac{E\{dN(t)\}}{dt}}=\lambda `$ (1) $`{\displaystyle \frac{E\{d\overline{N}(t)\}}{dt}}=0`$ (2) $`\mu (\tau )+\lambda \delta (\tau )={\displaystyle \frac{E[d\overline{N}(t)d\overline{N}(t+\tau )]}{dtd\tau }}`$ (3) $`S(f)=\lambda +{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mu (\tau )\mathrm{exp}(2\pi if\tau )𝑑\tau `$ (4) Where $`E`$ denotes the expectation operator. The autocovariance measures how likely it is that a spike will occur at time $`t+\tau `$ given that one has occurred at time $`t`$. Usually $`\mu (\tau )`$ is estimated rather than the full autocovariance which includes a delta function at zero lag<sup>3</sup><sup>3</sup>3When estimating the autocovariance using a histogram method one generally omits the spike at the start of the interval which would always fall in the bin nearest zero.. However, in order to take the Fourier transform the full autocovariance is required. The inclusion of this delta function leads to a constant offset of the spectrum. This offset is an important difference between continuous time processes and point processes. The population coherency $`\gamma (f)`$ is defined in equations 5 \- 7. $`\mu _{ab}(\tau )={\displaystyle \frac{E[d\overline{N}_a(t)d\overline{N}_b(t+\tau )]}{dtd\tau }}`$ (5) $`S_{ab}(f)={\displaystyle _{\mathrm{}}^{\mathrm{}}}\mu _{ab}(\tau )\mathrm{exp}(2\pi if\tau )𝑑\tau +\lambda _a\delta _{ab}`$ (6) $`\gamma (f)={\displaystyle \frac{S_{12}(f)}{\sqrt{S_{11}(f)S_{22}(f)}}}`$ (7) Where indices $`1`$ and $`2`$ denote simultaneously recorded spike trains from different cells. Unlike the spectrum, which is strictly real and positive, the coherency is a complex quantity. The modulus of the coherency, which is known as the coherence<sup>4</sup><sup>4</sup>4Some authors define coherence as the modulus squared of the coherency., can only vary between zero and one. This makes coherence particularly attractive for detecting relationships between spike trains as it is insensitive to the mean spike rates. ## 3 Examples and their interpretation Before discussing the details regarding how to estimate the spectrum and coherency it will be helpful to motivate them further by considering some simple examples. ### 3.1 Example population spectra For a homogeneous Poisson process of constant rate $`\lambda `$ the autocovariance is simply $`\lambda \delta (\tau )`$ and hence the spectrum is a constant equal to the rate $`\lambda `$. At the opposite extreme consider the case where the spikes are spaced by intervals $`\mathrm{\Delta }\tau `$. This is not a stationary process but if a small amount of drift is permitted, so that over an extended period there is nothing special about a given time, it becomes stationary. The spectrum of this process contains sharp lines at integer multiples of $`f=\frac{1}{\mathrm{\Delta }\tau }`$. Due to the drift the higher harmonics will become increasingly blurred and in the high frequency limit the spectrum will tend towards a constant value of the mean rate $`\lambda `$. As a final example consider the case where $`\mu (\tau )`$ is a negative Gaussian centered on zero $`\tau `$. This form of $`\mu (\tau )`$ is consistent with the probability of firing being suppressed after firing<sup>5</sup><sup>5</sup>5This need not necessarily correspond to the biophysical refractive period but, it could arise, rather from a characteristic integration time.. The spectrum of this process will be below $`\lambda `$ at low frequencies and will go to a constant value $`\lambda `$ at high frequencies. Figure 2 illustrates these different population spectra. ### 3.2 Example population coherency The population coherency of two homogeneous Poisson processes is zero. In contrast if two spike trains are equal then the coherence is one and the phase of the coherency is zero at all frequencies. If two spike trains are identical but offset by a lag $`\mathrm{\Delta }\tau `$ then the coherence will again be one but the phase of the coherency will vary linearly with frequency with a slope proportional to $`\mathrm{\Delta }\tau `$ and given by $`\varphi (f)=2\pi f\mathrm{\Delta }\tau `$. ## 4 Estimating the spectrum Section 3 demonstrated that the population spectrum may provide insights into the nature of a spike train. In this section the question of how to estimate the spectrum from a finite section of data will be introduced. In what follows the population quantity $`\lambda `$ in the definition of $`\overline{N}(t)`$ is replaced by a sample estimate $`N(T)/T`$. ### 4.1 Direct Spectral Estimators #### 4.1.1 Definition A popular, though seriously flawed, method for estimating the spectrum is to take the modulus squared of the Fourier transform of the data $`d\overline{N}(t)`$. This estimate is known as the Periodogram and is the simplest example of a direct spectral estimator. More generally, a direct spectral estimator is the modulus squared of the Fourier transform of the data but with the data being multiplied by an envelope function $`h(t)`$, known as a taper \[Percival and Walden, 1993\]. Equations 8 \- 10 define the direct estimator. On substituting $`\overline{N}(t)`$ into equation 9 a form amenable to implementation on a computer is obtained (equation 11). In this form the Fourier transform may be computed rapidly and without the need for the binning of data. Note that equation 10 results in $`h(t)`$ scaling as $`1/\sqrt{T}`$ as the sample length is altered. This ensures proper normalization of the Fourier transformation as sample size varies. $`I^D(f)=\|J^D(f)\|^2`$ (8) $`J^D(f)={\displaystyle _0^T}h(t)e^{2\pi ift}𝑑\overline{N}(t)`$ (9) Where, $`{\displaystyle _0^T}h(t)^2𝑑t=1`$ (10) $`J^D(f)={\displaystyle \underset{j=1}{\overset{N(T)}{}}}h(t_j)e^{2\pi ift_j}{\displaystyle \frac{N(T)H(f)}{T}}`$ (11) and $`H(f)`$ is the Fourier transform of the taper. The direct estimator suffers from bias and variance problems, described below, and is of no practical relevance for a single spike train sample. #### 4.1.2 Bias It may not be immediately apparent why the above procedure is an estimate of the spectrum, especially when one is permitted to multiply the data by an arbitrary, albeit normalized, taper. The relation between $`I^D(f)`$ and the spectrum may be obtained by taking the expectation of equation 8. $`E\{I^D(f)\}=E\{{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(t)h(t^{})e^{2\pi if(tt^{})}𝑑\overline{N}(t)𝑑\overline{N}(t^{})\}`$ (12) Assuming that the integration and expectation operations may be interchanged and substituting equation 3 yields <sup>6</sup><sup>6</sup>6For the moment, we assume that the population quantity $`\lambda `$ is known. This is of course not the case in practice, and one employs the estimate $`N(T)/T`$ as stated before. The effect of this extra uncertainty is given in equation 15., $$E\{I^D(f)\}=_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}h(t)h(t^{})e^{2\pi if(tt^{})}\{\mu (tt^{})+\lambda \delta (tt^{})\}𝑑t𝑑t^{}$$ (13) Which may be rewritten in the Fourier domain as, $$E\{I^D(f)\}=_{\mathrm{}}^{\mathrm{}}S(f^{})\|H(ff^{})\|^2𝑑f^{}$$ (14) The expected value of the direct estimator is a convolution of the true spectrum and the modulus squared of the Fourier transform of the taper. The normalization condition on the taper (equation 10) is equivalent to the requirement that the kernel of the convolution has unit area underneath it. Sharp features in the true spectrum will be thus be smeared by an amount which depends on the width of the taper in the frequency domain. If the taper is well localized in the frequency domain the expected value of the direct estimate is close to the true spectrum but if the taper is poorly localized then the expected value of the direct estimator will be incorrect i.e. the direct spectral estimator is biased. There is a fundamental level beyond which the bias cannot be reduced, due to the uncertainty relation forbidding simultaneous localization of a function in the time and frequency domains below a given limit. Since the maximum width of the taper is T the minimum frequency spread is 1/T which is known as the Raleigh frequency. Figure 3 shows the smoothing kernel for a rectangular taper and a T of 0.5s. Note that this kernel has large sidelobes which is the primary motivation for using tapering. In the above argument equation 3 was used in spite of the appearance of the population quantity $`\lambda `$ rather than the sample estimate $`N(T)/T`$ for which equation 12 was defined. A more careful treatment, which includes this correction, leads to an additional term at finite sample sizes in the expectation of the direct spectral estimator at low frequencies. The full expression is given below, $$E\{I^D(f)\}=_{\mathrm{}}^{\mathrm{}}S(f^{})\|H(ff^{})\|^2𝑑f^{}\|H(f)\|^2S(0)/T$$ (15) In the case of the periodogram, where $`h(t)=1/\sqrt{T}`$, the effect is clear since in this case $`J^D(0)=0`$ and hence $`I^D(0)=0`$ for any set of spike times and any T. #### 4.1.3 Asymptotic variance In the previous section it was shown that provided the taper is sufficiently local in frequency the expected value of the direct spectral estimator will be close to the true spectrum. However, the fact that the estimate is on average close to the true spectrum belies a serious problem with direct spectral estimators, namely that the estimates have very large fluctuations about this mean. The underlying source of this problem is that one is attempting to estimate the value of a function at an infinite number of points using a finite sample of data. The problem manifests itself in the fact that direct spectral estimators are inconsistent estimators of the spectrum<sup>7</sup><sup>7</sup>7Inconsistent estimators have a finite variance even for an infinite length sample.. In fact it may be shown that, under fairly general assumptions, the estimates are distributed exponentially (or equivalently as $`S(f)\chi _2^2/2`$) for asymptotic sample sizes (i.e. $`T\mathrm{}`$) \[Brillinger, 1972\]. Figure 4 illustrates that direct spectral estimators are noisy and untrustworthy, a fact emphasised by the observation that the $`\chi _2^2`$ distribution has a standard deviation equal to its mean. In the next three subsections methods for reducing the variance of direct spectral estimators using different forms of averaging will be discussed. ### 4.2 Trial averaging If there are a number of trials ($`N_T`$) available then the variance of the direct estimator may be reduced by trial averaging. $$I^{DT}(f)=\frac{1}{N_T}\underset{n=1}{\overset{N_T}{}}I_n^D(f)$$ (16) Where $`I_n^D(f)`$ is the direct spectral estimate based on the $`n^{th}`$ trial. In the large T limit taking the average entails summing $`N_T`$ independent samples from a $`\chi _2^2`$ distribution the result of which is distributed as $`\chi _{2N_T}^2`$. The reduction in variance is inversely proportional to the number of trials corresponding to a reduction in standard deviation which is the familiar factor of $`1/\sqrt{N_T}`$. At first sight it appears one would be getting something for nothing by breaking a single section of data into $`N_T`$ segments and treating them as separate trials. This is, of course, not the case. The reason is that if the data is segmented into short length samples, there is loss of frequency resolution proportional to the inverse of segment length. Lag window and multitaper estimators use the information from these independent estimates without artificially segmenting the data. ### 4.3 Lag Window Estimates A powerful property of the frequency domain is that, unless two frequencies are very close together, direct estimates of the spectrum of a stationary process at different frequencies are nearly uncorrelated. This property arises when the covariance between frequencies falls off rapidly. If the true spectrum varies slowly over the width of the covariance then the large sample covariance of a direct spectral estimator is given by equation 17. $$cov\{I^D(f_1),I^D(f_2)\}E\{I^D(\overline{f})\}^2\left\|_{\mathrm{}}^{\mathrm{}}h(t)^2e^{2\pi i\mathrm{\Delta }ft}𝑑t\right\|^2$$ (17) Where $`\overline{f}=(f_1+f_2)/2`$ and $`\mathrm{\Delta }f=f_1f_2`$ For $`\mathrm{\Delta }f=0`$, this expression reduces to the previously mentioned result that the variance of the estimator is equal to the square of the mean. For $`\mathrm{\Delta }f>>1/T`$, $`\|_{\mathrm{}}^{\mathrm{}}h(t)^2e^{2\pi i\mathrm{\Delta }ft}𝑑t\|^20`$, since $`h(t)^2`$ is a smooth function with extent $`T`$. This implies that $`cov\{I^D(f_1),I^D(f_2)\}0`$ for $`\|f_1f_2\|>>1/T`$. The approximate independence of nearby points means that, if the true spectrum varies slowly enough, then closely spaced points will provide several independent estimates of the same underlying spectrum. This is the motivation for the lag window estimator which is simply a smoothed version of the direct spectral estimator \[Percival and Walden, 1993\]. The lag window estimator is defined in equations 18 and 19. $$I^{LW}(f)=_{\mathrm{}}^{\mathrm{}}K(ff^{})I^D(f^{})𝑑f^{}$$ (18) Where, $$_{\mathrm{}}^{\mathrm{}}K(f)𝑑f=1$$ (19) Averaging over trials may be included by using the trial averaged direct spectral estimate $`I^{DT}`$ (see equation 16) in place of $`I^D`$ in the above expression. It is assumed that $`K(f)`$ is a smoothing kernel with reasonable properties. #### 4.3.1 Bias The additional smoothing of the lag window kernel modifies the bias properties of the estimator from those expressed in equation 15. The expected value of the lag window estimator is given by, $$E\{I^{LW}(f)\}=_{\mathrm{}}^{\mathrm{}}K(ff^{})\|H(f^{}f^{\prime \prime })\|^2S(f^{\prime \prime })𝑑f^{}𝑑f^{\prime \prime }\frac{S(0)}{T}_{\mathrm{}}^{\mathrm{}}K(ff^{})\|H(f^{})\|^2𝑑f^{}$$ (20) #### 4.3.2 Asymptotic Variance The large sample variance of this estimator is readily obtained using equation 17. $$var\{I^{LW}(f)\}=\frac{\xi }{N_T}E\{I^{LW}(f)\}^2$$ (21) Where, $$\xi =_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}K(f)K(f^{})\|(ff^{})\|^2𝑑f𝑑f^{}$$ (22) and, $$(f)=_{\mathrm{}}^{\mathrm{}}h(t)^2e^{2\pi ift}𝑑t$$ (23) Equation 21 includes the reduction in variance due to trial averaging. $`1/\xi `$ can be interpreted as the effective number of independent estimates beneath the smoothing kernel, as demonstrated by the following qualitative argument. If $`\mathrm{\Delta }f`$ is the frequency width of the smoothing kernel $`K(f)`$ and $`\delta f`$ is the frequency width of the taper $`(f)`$ then since $`K(f)1/\mathrm{\Delta }f`$ it follows that $`\xi 1/(\mathrm{\Delta }f)^2_{\mathrm{\Delta }f}_{\mathrm{\Delta }f}\|(ff^{})\|^2𝑑f𝑑f^{}`$ and hence that $`\xi \delta f/\mathrm{\Delta }f.`$ ### 4.4 Multitaper Estimates While the lag window estimator is based on the idea that nearby frequencies provide independent estimates, the estimation is not very systematic, since, one should be able to explicity decorrelate nearby frequencies from the knowledge of the correlations introduced by a finite window size. This is achieved in multitaper spectral estimation. The basic idea of multitaper spectral estimation is to average the spectral estimates from several orthogonal tapers. The orthogonality of the tapers ensures that the estimates are uncorrelated for large samples (consider substituting $`h_1(t)h_2(t)`$ for $`h(t)^2`$ in equation 17). A critical question is the choice of a set of orthogonal tapers. A natural choice are the discrete prolate spheroidal sequences (dpss) or Slepian sequences, which are defined by the property that they are maximally localised in frequency. The dpss tapers maximize the spectral concentration defined as; $$\lambda =\frac{_W^W\|H(f)\|^2𝑑f}{_{\mathrm{}}^{\mathrm{}}\|H(f)\|^2𝑑f}$$ (24) Where in the time domain $`h(t)`$ is strictly confined to the interval \[0,T\]. For given values of W and T there are a finite number of tapers which have concentrations ($`\lambda `$) close to one, and therefore have well controlled bias. This number is known as the Shannon number and is $`2WT`$. This sets an upper limit on the number of independent estimates that can be obtained for a given amount of spectral smoothing. A direct multitaper estimate of the spectrum is defined in equation 25. $$I^{MT}(f)=\frac{1}{K}\underset{k=0}{\overset{K1}{}}I_k^D(f)$$ (25) The eigenspectra $`I_k^D`$ are direct spectral estimates based on tapering the data with the $`k^{th}`$ dpss function. As previously trial averaging can be included by using $`I^{DT}`$ rather than $`I^D`$. More sophisticated estimates involve adaptive (rather than constant) weighting of the data tapers \[Percival and Walden, 1993\]. Multitaper spectral estimation has been recently shown to be useful for analysing neurobiological time series, both continuous processes \[Mitra and Pesaran, 1999\] and spike trains \[Pesaran et al., 2000\]. #### 4.4.1 Bias The bias for the multitaper estimate is given by equation 15 but with $`\|H()\|^2`$ replaced by an average over tapers $`\frac{1}{K}_{k=0}^{K1}\|H_k()\|^2`$. #### 4.4.2 Asymptotic Variance The asymptotic variance of the multitaper estimator, including trial averaging, is given by equation 26. $$var\{I^{MT}(f)\}=\frac{1}{N_TK}E\{I^{MT}(f)\}^2$$ (26) ### 4.5 Degrees of freedom At this point it is useful to introduce the concept of the degrees of freedom ($`\nu _0`$) of an estimate. The degrees of freedom is twice the number of independent estimates of the spectrum. Degrees of freedom is a useful concept as it permits the expressions for the variance of the different estimators to be written in a common format. $$var\{I^X(f)\}=\frac{2E\{I^X(f)\}^2}{\nu _0}$$ (27) Where, \| $`X`$ \| $`D`$ \| $`DT`$ \| $`LW`$ \| $`MT`$ \| \| --- \| --- \| --- \| --- \| --- \| \| $`\nu _0`$ \| $`2`$ \| $`2N_T`$ \| $`2N_T/\xi `$ \| $`2N_TK`$ \| Degrees of freedom is also a useful framework in which to cast both finite size corrections and the confidence limits for the spectra and coherence. The variance of estimators of the spectrum can be estimated using internal methods such as the bootstrap or jackknife \[Efron and Tibshirani, 1993\],\[Thomson and Chave, 1991\]. Jackknife estimates can be constructed over trials or over tapers. If $`\nu _0`$ is large ($`>20`$), then the theoretical and Jackknife variance are in close agreement. If distributional assumptions can be validly made about the point process, theoretical error bars have an important advantage over internal estimates since they enable the understanding of different factors which enter into the variance in order to guide experimental design and data analysis. However Jackknife estimates are less sensitive to failures in distributional assumptions, and this provides them with statistical robustness. It is conventional to display spectra on a log scale. This is because taking the log of the spectrum stabilizes the variance and leads a distribution which is approximately Gaussian. ### 4.6 Confidence intervals The expected values of the estimators and also their variance have been discussed for several different spectral estimators but it is desirable to put confidence intervals on the spectral estimates rather than standard deviations. As previously mentioned in section 4.2 the averaging of direct spectral estimates from different trials yields, in the large sample limit, estimates which are distributed as $`\chi _{2N_T}^2`$. In general for the other estimates a well known approximation \[Percival and Walden, 1993\] is to assume that the estimate is distributed as $`\chi _{\nu _0}^2`$. Confidence intervals can then by placed on estimates on the basis of this $`\chi _{\nu _0}^2`$ distribution. The confidence interval applies for the population spectrum $`S(f)`$ and is obtained from the following argument. $$P\left[q_1\chi _{\nu _0}^2q_2\right]=12p$$ (28) Where $`P`$ indicates probability, $`q_1`$ is such that $`P[\chi _{\nu _0}^2q_1]=p`$ and $`q_2`$ is such that $`P[\chi _{\nu _0}^2q_2]=p`$. It follows that, $$P\left[q_1\nu _0I^X(f)/S(f)q_2\right]=12p$$ (29) Hence an approximate $`100\%\times (12p)`$ confidence interval for $`S(f)`$ is given by, $$P\left[\nu _0I^X(f)/q_2S(f)\nu _0I^X(f)/q_1\right]$$ (30) For large $`\nu _0`$ ($`>`$ 20) these confidence intervals do not differ substantially from those based on a Gaussian ($`\pm `$2 standard deviations) but at small $`\nu _0`$ the difference can be substantial as for these values the $`\chi _{\nu _0}^2`$ distribution has long tails. ### 4.7 High Frequency limit The population spectrum goes to a constant value equal to the rate $`\lambda `$ in the high frequency limit. In practice spectra calculated from a finite sample will go to a value close to $`\lambda `$ but fluctuations in the number of spikes in the interval will lead to an error in this estimate. For a given sample the spectrum will go the value given by equation 31. $$I(f\mathrm{})=\frac{1}{N_TK}\underset{k=0}{\overset{K1}{}}\underset{n=1}{\overset{N_T}{}}\underset{j=1}{\overset{N_n(T)}{}}h_k(t_j^n)^2$$ (31) Where $`t_j^n`$ is the $`j^{th}`$ spike in the $`n^{th}`$ trial and $`N_n(T)`$ is the total number of spikes in the $`n^{th}`$ trial. In the case of direct and lag window estimators the averaging over tapers need not be performed. This expression yields a value which is typically very close to the sample estimate of the mean rate<sup>8</sup><sup>8</sup>8It is exactly the sample estimate of the mean rate for a rectangular taper.. It is significant departures from this high frequency limit which are of interest when interpreting the spectrum as these indicate enhancement or suppression relative to a homogeneous Poisson process. ### 4.8 Choice of estimator, taper and lag window The preceding section discussed the large sample statistical properties of direct, lag window and multitaper estimates of the spectrum. The choice of which estimator to use remains a contentious one \[Percival and Walden, 1993\]. The multitaper method is the most systematic of the estimators but the lag window estimators should perform almost as well for those spike train spectra which have reasonably small dynamic ranges<sup>9</sup><sup>9</sup>9Dynamic range is a measure of the variation in the spectrum as a function of frequency and is defined as $`10log_{10}(\frac{max_fS(f)}{min_fS(f)})`$.. However, it is possible to construct spike trains with widely different time scales, which can possess a large dynamic range. In addition, the multitaper technique leads to a simple jackknife procedure by leaving out one data taper in turn. A further important property of the multitaper estimator is that it gives more weight to events at the edges of the time interval and thus ameliorates the arbitrary downweighting of the edges of the data introduced by single tapers. If using the lag window estimator there are many choices available for both the taper and the lag window. The choice of taper is generally not critical provided that the taper goes smoothly to zero at the start and end of the interval. A rectangular taper has particularly large sidelobes in the frequency domain which can lead to significant bias. The choice of lag window is also usually not critical and typically a Gaussian kernel will be satisfactory. ## 5 Estimating the Coherency Sample coherency, which may be estimated using equation 32, may be evaluated using any of the previously mentioned spectral estimators. The principle difference is that the direct estimator, in terms of which the other estimators are expressed, is given by equations 33 and 34 rather than 8 and 9. $`C(f)`$ $`=`$ $`I_{12}^X/\sqrt{I_{11}^XI_{22}^X}`$ (32) $`I_{12}^D(f)`$ $`=`$ $`J_1^D(f)J_2^D(f)^{}`$ (33) $`J_a^D(f)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h(t)e^{2\pi ift}𝑑\overline{N}_a(t)`$ (34) Where the $`\overline{N}_1(t)`$ and $`\overline{N}_2(t)`$ are simultaneously recorded spike trains from two different cells and $`X`$ denotes the type of spectral estimator. Possible choices of estimator $`X`$ include; $`D`$ direct, $`DT`$ trial averaged direct, $`LW`$ lag window or $`MT`$ multitaper. Lag window and multitaper coherency estimates may be constructed by substituting $`I_{12}^D()`$ in place of $`I^D()`$ in equations 18 and 25. The estimates are biased over a frequency range equal to the width of the smoothing although the exact form for the bias is difficult to evaluate. ### 5.1 Confidence limits for the Coherence The treatment of error bars is somewhat different between the spectrum and the coherency, since the coherency is a complex quantity. Usually one is interested is in establishing whether there is significant coherence in a given frequency band. In order to do this the sample coherence should be tested against the null hypothesis of zero population coherence. The distribution of the sample coherence under this null hypothesis is given below. $$P(\|C\|)=(\nu _02)\|C\|(1\|C\|^2)^{(\nu _0/22)}0\|C\|1$$ (35) A derivation of this result is given in \[Hannan, 1970\]. In outline the method is to rewrite the coherence in such a way that it is equivalent to a multiple correlation coefficient \[Anderson, 1984\]. The distribution of a multiple correlation coeffient is then a known result from multivariate statistics. In the case of coherence estimates based on lag window estimators the appropriate $`\nu _0`$ may be used although this is only approximately valid because this method of derivation assumes integer $`\nu _0/2`$. It is straightforward to calculate a confidence level based on this distribution. The coherence will only exceed $`\sqrt{1p^{1/(\nu _0/21)}}`$ in $`p\times 100\%`$ of experiments. In addition it is notable that the quantity $`(\nu _0/21)\|C\|^2/(1\|C\|^2)`$ is distributed as $`F_{2,\nu _02)}`$ under this null hypothesis. It is useful to apply a transformation to the coherence before plotting it which aids in the assessment of significance. The variable $`q=\sqrt{(\nu _02)log(1\|C\|^2)}`$ has a Raleigh distribution which has density $`p(q)=qe^{q^2/2}`$. This density function does not depend on $`\nu _0`$ and furthermore has a tail which closely resembles a Gaussian. For certain values of a fitting parameter<sup>10</sup><sup>10</sup>10A reasonable choice for $`\beta `$ is 23/20. $`\beta `$, a further linear transformation $`r=\beta (q\beta )`$ leads to a distribution which closely resembles a standard normal Gaussian for $`r>2`$. This means that for $`r>2`$ one can interpret $`r`$ as the number of standard deviations by which the coherence exceeds that expected under the null hypothesis. ### 5.2 Confidence Limits for the Phase of the Coherency If there is no population coherency then the phase of the sample coherency is distributed uniformly. If, however, there is population coherency then the distribution of the sample phase is approximately Gaussian provided that the tails of the Gaussian do not extend beyond a width $`2\pi `$. An approximate 95% confidence interval for the phase \[Rosenberg et al., 1989, Brillinger, 1974\] is given below. $$\widehat{\varphi }(f)\pm 2\sqrt{\frac{2}{\nu _0}\left(\frac{1}{\|C(f)\|^2}1\right)}$$ (36) Where $`\widehat{\varphi }(f)`$, the sample estimate of the coherency phase, is evaluated using $`tan^1\{Im(C)/Re(C)\}`$. ## 6 Finite Size Effects In the preceding sections error bars were given for estimators of the spectrum and the coherence. However these error bars were based on large sample sizes (they apply asymptotically as $`T\mathrm{}`$). Neurophysiological data are not collected in this regime and, particularly in awake behaving studies where data is often sparse, corrections arising at small T are potentially important. In order to estimate the size of these corrections a particular model for the point process is required. The model studied was chosen primarily for its analytical tractability while still maintaining a non-trivial spectrum. The model and the final results will be presented here but the details of the analysis are reserved until appendix A. The model is a doubly stochastic inhomogeneous Poisson process with a Gaussian rate function. A specific realization of a spike train is generated from the model in the following manner. Firstly a population spectrum $`S_G(f)`$ is specified. From this a realization of a zero mean Gaussian process $`\lambda _G(f)`$ is generated. A constant $`\lambda `$, the mean rate, is then added to this realization. This function is then considered to be the rate function for an inhomogeneous Poisson process. A realization of this inhomogeneous Poisson process is then generated. A schematic of the model is shown in figure 5. Technically this is not a valid process because the rate function $`\lambda (t)`$ may be negative. However, if the area underneath the spectrum is small enough then the fluctuations about the mean rate seldom cross zero and corrections due to this effect are negligible. In addition large violations of this area constraint have been tested by Monte Carlo simulation and the results still apply to a good approximation. An important feature of this model is that the population spectrum of the spike trains is simply the spectrum of the inhomogeneous Poisson process rate function plus an offset equal to the mean rate<sup>11</sup><sup>11</sup>11This result does not depend on the Gaussian assumption.. The spectrum of the rate function is a positive real quantity and therefore within this model the population spectrum cannot be less than the mean rate at any frequency. Intuitively, the reason for this is that the process must be more variable than a homogeneous Poisson process at all frequencies. To make the nature of the result clear a simplified version is given in equation 37. This version is for the particular case of a homogeneous Poisson process (which has a flat population spectrum) and a rectangular taper<sup>12</sup><sup>12</sup>12The expression also holds approximately for the multitaper estimate provided all tapers up to the Shannon limit are used.. $$var\{I^X(f)\}=\lambda ^2\left[\frac{2}{\nu _0}+\frac{1}{N_TT\lambda }\right]$$ (37) Where $`\lambda `$ is the mean rate. A sample based estimate of $`N_TT\lambda `$ is the total number of spikes over all trials. It is to be noted that finite size effects reduce the degrees of freedom. This result implies that there is a point beyond which additional smoothing does not decrease the variance further and this point is approximately when $`\nu _0`$ is equal to twice the total number of spikes. The full result is given in equations 38 \- 43. $`var\{I^X(f)\}={\displaystyle \frac{2E\{I^X(f)\}^2}{\nu (f)}}`$ (38) $`{\displaystyle \frac{1}{\nu (f)}}={\displaystyle \frac{1}{\nu _0}}+{\displaystyle \frac{C_h^X\mathrm{\Phi }(f)}{2TN_TE\{I^X(f)\}^2}}`$ (39) Where, $$C_h^X=\{\begin{array}{c}_0^1f(t)^4𝑑t\mathrm{If}\mathrm{X}=\mathrm{LW},\mathrm{D}\mathrm{or}\mathrm{DT}\hfill \\ \frac{1}{K^2}\underset{k,k^{}}{}_0^1f_k(t)^2f_k^{}(t)^2𝑑t\mathrm{If}\mathrm{X}=\mathrm{MT}\hfill \end{array}$$ (40) $$f(t/T)=\sqrt{T}h(t)$$ (41) and, $$\mathrm{\Phi }(f)=\lambda _{hf}+4[E\{I^X(f)\}\lambda _{hf}]+2[E\{I^X(0)\}\lambda _{hf}]+[E\{I^X(2f)\}\lambda _{hf}]$$ (42) $$\lambda _{hf}=E\{I^X(f\mathrm{})\}$$ (43) $`C_h^X`$ is a constant of order unity which depends on the taper. When a taper is used to control bias some of the spikes are effectively disregarded and this has an effect on the size of the correction. The function $`f(t)`$ has the same form as the taper $`h(t)`$ but is defined for the interval $`[0,1]`$. $`C_h^X`$ is the integral of the fourth power of $`f`$ and obtains its minimum value of one for a rectangular taper. It is typically between 1 and 2 for other tapers. In the multitaper case cross terms between tapers are included. Equation 42 describes how the finite size correction depends on the structure of the spectrum. $`\mathrm{\Phi }(f)`$ is the sum of four terms. The first term is the only term which is present for a flat spectrum. The second term is a correction which depends on the spectrum at the frequency being considered. The next two terms depend on the spectrum at zero frequency and the spectrum at twice the frequency being considered. The latter three terms all depend on the difference between the spectrum at some frequency and the high frequency limit. Equation 42 applies provided that the spike train is well described by the model. However, this is not necessarily the case and a suppression of the spectrum, which cannot be described by the model, often occurs at low frequencies<sup>13</sup><sup>13</sup>13Note that any spike train spectra displaying significant suppression below the mean firing rate can immediately rule out the inhomogeneous Poisson process model.. In the event that there is a significant suppression of the spectrum $`\mathrm{\Phi }(f)`$ may become small or even negative. To avoid this a modified form for $`\mathrm{\Phi }(f)`$ may be used which prevents this. $`\mathrm{\Phi }(f)=\lambda _{hf}+4\mathrm{m}\mathrm{a}\mathrm{x}([\mathrm{E}\{\mathrm{I}^\mathrm{X}(\mathrm{f})\}\lambda _{\mathrm{hf}}],0)`$ $`+`$ $`2\mathrm{m}\mathrm{a}\mathrm{x}([\mathrm{E}\{\mathrm{I}^\mathrm{X}(0)\}\lambda _{\mathrm{hf}}],0)\mathrm{}`$ (44) $`+`$ $`\mathrm{max}([\mathrm{E}\{\mathrm{I}^\mathrm{X}(2\mathrm{f})\}\lambda _{\mathrm{hf}}],0)`$ The above modification to the result is somewhat ad hoc so Monte Carlo simulations of spike trains with enforced refractory periods have been performed to test its validity. These simulations demonstrated that, although the correction derived using 44 was significantly different from that obtained from the Monte Carlo simulations in the region of the suppression, equation 44 provided a pessimistic estimate in all cases studied. This increases confidence that applying finite size corrections using equation 44 will provide reasonable error bars for small samples. Equation 39 gives the finite size correction in terms of a reduction in $`\nu _0`$. The new $`\nu (f)`$ may be used to put confidence intervals on the results, as described in section 4.6, although the accuracy of the $`\chi _\nu ^2`$ assumption will be reduced. In the case of the coherence an indication of the correction to the confidence level can be obtained by using the smaller of the two $`\nu (f)`$ from the spike train spectra to calculate the confidence level using equation 35. In all cases if the effect being observed only achieves significance by an amount which is of the same order as the finite size correction then it is recommended that more data be collected. ## 7 Experimental Design Often it is useful to know in advance how many trials or how long a time interval one needs in order to resolve features of a certain size in the spectrum or the coherence. To do this one needs to estimate the asymptotic degrees of freedom $`\nu _0`$. This depends on the size of feature to be resolved $`\alpha `$, the significance level for which confidence intervals will be calculated $`p`$ and the fraction of experiments which will achieve significance $`𝒫`$. In addition the reduction in the degrees of freedom due to finite size effect depends on the total number of spikes $`N_s`$ and also $`C_h`$ (see section 6). An estimate of $`v_0`$ may be obtained in two stages. Firstly $`\alpha `$,$`p`$ and $`𝒫`$ are specified and used to calculate a degrees of freedom $`\nu `$. Secondly the asymptotic degrees of freedom $`\nu _0`$ is estimated using $`\nu `$, $`N_s`$ and $`C_h`$. The feature size $`\alpha =(S\lambda )/\lambda `$ is the minimum size of feature which the experimenter is content to resolve. For example, a value of 0.5 indicates that where the population spectrum exceeds $`1.5\lambda `$ the feature will be resolved. The significance level should be set to the same value that will be used for calculating the confidence interval for the spectrum, typically be $`0.05`$. For a given $`p`$ there is some probability $`𝒫`$ that an experiment will achieve significance. To calculate $`\nu `$ one begins with a guess $`\nu _g`$. Then $`q_1`$ is chosen such that $`P\left[\chi _{\nu _g}^2q_1\right]=p/2`$. On the basis of this one then evaluates $`𝒫_\mathcal{0}=1\mathrm{\Phi }\left[q_1/(1+\alpha )\right]`$ where $`\mathrm{\Phi }`$ is the cumulative $`\chi _{\nu _g}^2`$ distribution<sup>14</sup><sup>14</sup>14These formulae apply for $`\alpha >0`$. If $`\alpha <0`$ then $`P\left[\chi _{\nu _g}^2q_1\right]=p/2`$ and $`𝒫_\mathcal{0}=\mathrm{\Phi }\left[q_1/(1+\alpha )\right]`$ should be used.. If $`𝒫_\mathcal{0}`$ is equal to the specified fraction $`𝒫`$ then $`\nu =\nu _g`$ otherwise a different $`\nu _g`$ is chosen. This procedure is readily implemented as a minimization of $`(𝒫𝒫_\mathcal{0}(\nu _g))^2`$ on a computer. Having obtained $`\nu `$ one can estimate $`\nu _0`$ using the following expression. $$\frac{1}{\nu _0}=\frac{1}{\nu }\frac{C_h\left[1+4\alpha \right]}{2N_s\left[1+\alpha \right]^2}$$ (45) Where the $`4\alpha `$ is omitted from the numerator if $`\alpha <0`$. Figure 6 illustrates example design curves generated using this method. These curves show the asymptotic degrees of freedom as a function of feature size for different total numbers of spikes. The existence of a region bounded by vertical asymptotes implies that as long as the total number of measured spikes is finite, modulations in the spectrum below a certain level cannot be detected no matter how much the spectrum is smoothed. These curves may be used to design experiments capable of resolving spectral features of a certain size. In the case of the coherence one calculates how many degrees of freedom are required for the confidence line to lie at a certain level as described in section 5.1. ## 8 Line Spectra One of the assumptions underlying the estimation of spectra is that the population spectrum varies slowly over the smoothing width ($`W`$ for multitaper estimators). While this is often the case there are situations in which the spectrum contains very sharp features which are better approximated by lines than by a continuous spectrum. This corresponds to periodic modulations of the underlying rate, such as when a periodic stimulus train is presented. In such situations it is useful to be able to test for the presence of a line in a background of colored noise (i.e. in a locally smooth but otherwise arbitrary continuous population spectrum). Such a test has been previously developed, in the context of multitaper estimation, for continuous processes \[Thomson, 1982\] and in the following section the analogous development for point processes is presented. ### 8.1 F-test for point processes A line in the spectrum has an exactly defined frequency and consequently the process $`N(t)`$ has a non-zero first moment. The natural model in the case of a single line is given by equation 46. $$E\{dN(t)\}/dt=\lambda _0+\lambda _1cos(2\pi f_1t+\varphi )$$ (46) A zero mean process ($`\overline{N}`$) may be constructed by subtraction of an estimate of $`\lambda _0t`$. Provided that the product of the line frequency($`f_1`$) and the sample duration(T) is much greater than one the sample quantity $`N(T)/T`$ is an approximately unbiased estimate of $`\lambda _0`$. The resultant zero mean process $`\overline{N}`$ has a Fourier transform which has a non-zero expectation. $`J_k(f)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}h_k(t)e^{2\pi ift}𝑑\overline{N}(t)`$ (47) $`E\{J_k(f)\}`$ $`=`$ $`c_1H_k(ff_1)+c_1^{}H_k(f+f_1)`$ (48) Where, $$c_1=\lambda _1e^{i\varphi }/2$$ (49) In the case where $`f>0`$ and $`f_1>W`$, $$E\{J_k(f)\}c_1H_k(ff_1)$$ (50) The estimates of $`J_k(f_1)`$ from different tapers provide a set of uncorrelated estimates of $`c_1H_k(0)`$. It is hence possible to estimate the value of $`c_1`$ by complex regression. $$\widehat{c}_1=\frac{_kJ_k(f_1)H_k(0)}{_k\|H_k(0)\|^2}$$ (51) Under the null hypothesis that there is no line in the spectrum ($`c_1=0`$) it may readily be shown that $`E\{\widehat{c}_1\}=0`$ and $`var\{\widehat{c}_1\}=S(f_1)/_k\|H_k(0)\|^2`$. The residual spectrum<sup>15</sup><sup>15</sup>15It is also possible to estimate a residual coherency. In order to do this one uses a residual cross-spectrum $`\widehat{S}_{xy}(f)=\frac{1}{K}_k(J_k^x(f)\widehat{c}_1^xH_k(ff_1))^{}(J_k^y(f)\widehat{c}_1^yH_k(ff_1))`$, together with the residual spectra to evaluate the usual expression for coherency., which has the line removed, may be estimated using equation 52. $$\widehat{S}(f)=\frac{1}{K}\underset{k}{}\|J_k(f)\widehat{c}_1H_k(ff_1)\|^2$$ (52) In the large sample limit the distributions of both $`\widehat{c}_1`$ and $`\widehat{S}(f_1)`$ are known \[Percival and Walden, 1993\] and may be used to derive relation 53. $$\frac{\|\widehat{c}_1\|^2_k\|H_k(0)\|^2(K1)}{_k\|J_k(f_1)\widehat{c}_1H_k(0)\|^2}F_{2,2(K1)}$$ (53) Where $``$ denotes ‘is distributed as’. The null hypothesis may be tested using this relation and, if rejected, the line can be removed using equation 52 to estimate the residual spectrum. It is worth noting that although relation 53 was derived for large samples the test is remarkably robust as the sample size is decreased. Numerical tests indicate that the tail of the F distribution is well reproduced even in situations where there are as low as 5 spikes in total. ### 8.2 Periodic Stimulation A common paradigm in neurobiology where line spectra are particularly important is that of periodic stimulation. When a neuron is driven by a periodic stimulation of frequency $`f_1`$ the spectrum may contain lines at any of the harmonics $`nf_1`$. Provided that $`f_1>2W`$ the analysis of section 8.1 applies with each harmonic being separately tested for significance. The first moment of the process, which has period $`1/f_1`$, is given by equation 54 and may be estimated using $`\widehat{c}_n`$. $$\lambda (t)=\lambda _0+\underset{n}{}\lambda _ncos(2\pi nf_1t+\varphi _n)$$ (54) Where $`\lambda _n=2\|c_n\|`$, $`\varphi _n=\mathrm{tan}^1\{Im(c_n)/Re(c_n)\}`$, the sum is taken over all the significant coefficients. This rate function $`\lambda (t)`$ is the average response to a single stimulus or impulse response. The coefficients $`c_n`$ are the Fourier series representation of $`\lambda (t)`$. ### 8.3 Error Bars It is possible to put confidence intervals on both the modulus and the phase of the coefficients $`\widehat{c}_n`$. For large samples($`>10`$ spikes) the real and imaginary parts of $`\widehat{c}_n`$ are distributed as independent Gaussians each with standard deviation $`\sigma _n=\sqrt{S(nf_1)/(2_k\|H_k(0)\|^2)}`$. For $`c_n/\sigma _n>3`$ the distribution of $`\|\widehat{c}_n\|`$ is well approximated by a Gaussian centered on $`\|c_n\|`$ and with standard deviation $`\sigma _n`$. In addition the estimated phase angle ($`\widehat{\varphi }_n`$) is also almost Gaussian with mean $`\varphi _n`$ and standard deviation $`\sigma _n/\|c_n\|`$. Approximate error bars or confidence intervals may be obtained using a sample based estimate of $`\sigma _n`$, $`\widehat{\sigma }_n=\sqrt{\widehat{S}(nf_1)/(2_k\|H_k(0)\|^2)}`$. Estimating error bars for the impulse response function is more involved due to their non-local nature (if one of the Fourier coefficients is varied the impulse response function changes everywhere). It is therefore of interest to estimate a global confidence interval, defined as any interval such that the probability of the function crossing the interval anywhere is some predefined probability. A method for estimating a global confidence band is detailed in \[Sun and Loader, 1994\] and outlined here. First a basis vector $`\mathrm{\Phi }(t)`$ is constructed. $$\mathrm{\Phi }(t)=\left[\begin{array}{c}\widehat{\sigma }_1cos(2\pi f_1t)\\ \mathrm{}\\ \widehat{\sigma }_Ncos(2\pi f_Nt)\\ \widehat{\sigma }_1sin(2\pi f_1t)\\ \mathrm{}\\ \widehat{\sigma }_Nsin(2\pi f_Nt)\end{array}\right]$$ (55) Where $`N`$ is the total number of harmonics. The elements of this vector have unit variance and a standard approximation may be applied. $$P(sup\|\lambda (t)E\{\lambda (t)\}\|>c\mathrm{\Phi }(t))2(1N(c))+(k/\pi )e^{c^2/2}$$ (56) Where $`sup`$ is the maximum value of its operand, $`\mathrm{\Phi }(t)`$ denotes the length of vector $`\mathrm{\Phi }(t)`$, $`N(c)`$ is the cumulative standard normal distribution and $`k`$ is a constant. $`k`$ may be evaluated by constructing the $`2\times N`$ matrix $`X(t)=[\mathrm{\Phi }(t)d\mathrm{\Phi }(t)/dt]`$, forming its $`QR`$ decomposition \[Press et al., 1992\] and then evaluating $`k=_0^T\|R_{22}(t)/R_{11}(t)\|𝑑t`$. Confidence intervals for the residual spectrum are calculated in the usual manner (using $`\chi _\nu ^2`$) although at the line frequencies the interval is slightly broadened due to the loss of a degree of freedom incurred by estimation of $`c_n`$. Section 11 contains an example application of the methods described in this section. ## 9 Example Spectra Figure 7 is a spectrum calculated from data collected from a single cell recorded from area PRR in the parietal cortex of an awake behaving monkey during a delayed memory reach task \[Snyder et al., 1997\]. The spectrum is calculated over an interval of 0.5s during which the firing rate is reasonably stationary and is averaged over 5 trials. The spectrum shows two features which achieve significance. There is enhancement of the spectrum in the frequency band 20-40 Hz indicating the presence of an underlying broad band oscillatory mode in the neuronal firing rate. In addition there is suppression of the spectrum at low frequencies. As discussed previously a suppression of this sort is consistent with an effective refactory period during which the neuron is less likely to fire. Care must be taken at low frequencies since at frequencies comparable to the smoothing width the spectrum is particularly sensitive to any non-stationarity in the data. ## 10 Example Coherency To illustrate the estimation of coherency simulated spike trains were generated from a coupled doubly stochastic Poisson process. For a given trial a pair of rate functions were drawn from a Gaussian process. The realizations share a coherent mode which is linearly mixed into the rates of both cells. These coupled rate functions are then used to independently draw a realization of an inhomogeneous Poisson process for each cell. Using this method 15 trials of duration 0.5s were generated. The coherent mode was set such that the population coherence was a Gaussian of height 0.35 and standard deviation 5 Hz centered on 20 Hz. The phase of this mode was set to $`180^o`$. Figure 8 indicates that this coherent mode is reasonably estimated. ## 11 Example Periodic Stimulation An example of an analysis of a periodic stimulus paradigm is shown in figure 9. The data is a single cell recording collected from the barrel cortex of an awake behaving rat during periodic whisker stimulation at 5.5 Hz \[Sachdev et al., 1999\]. There is a single trial of duration 50s. The estimated impulse response function $`\widehat{\lambda }(t)`$ is seen to have two distinct sharp peaks outside of which the response does not differ significantly from zero. The moduli of the Fourier coefficients are significant out to $`n=25`$. This automatically sets the smoothing of $`\widehat{\lambda }(t)`$ as structure on a time scale of less than $`1/(25\times 5.5)=7`$ ms does not achieve significance. Note that the coefficients are enhanced at multiples of 6 (i.e.$`33`$ Hz) which comes from having two peaks in the time domain $`\lambda (t)`$ which are separated by $`30`$ ms. The phase of the coefficients closely follows a straight line but there is a small periodic deviation from this line which is again at index multiples of 6. The gradient of the straight line depends on the time delay of the response. The residual spectrum was calculated by first evaluating a multitaper estimate from which the significant harmonics were removed. This spectrum had a bandwidth of 1.5 Hz chosen to avoid overlap of the harmonics leading to the multitaper estimate being undersmoothed. A further smoothing was performed using a lag window<sup>16</sup><sup>16</sup>16The previous theory developed for lag window estimators applies to this hybrid estimator with $`\|H()\|^2`$ replaced by $`\frac{1}{K}_{k=0}^{K1}\|H_k()\|^2`$ in equation 20 and $`\|()\|^2`$ replaced by $`\frac{1}{K}_{k=0}^{K1}\|_k()\|^2`$ in equation 22. and the resultant spectrum, displays a slight but significant suppression relative to a Poisson process out to almost 200 Hz. Such a spectrum is characteristic of a short time scale refractive period. The residual spectrum is particularly useful because rate non-stationarity has been removed. ## 12 Summary It is the belief of the authors that spectral analysis is a fruitful and under exploited analysis technique for spike trains. In this paper an attempt has been made to collect the machinery necessary for performing spectral analysis on spike train data into a single document. Starting from the population definitions the statistical properties of estimators of the spectrum and coherency have been reviewed. Estimation methods for both continuous spectra and spectra which contain lines have been included. In addition new corrections to asymptotic error bars have been presented which increase confidence in applying spectral techniques in practical situations where data is often sparse. Tables 1 to 5 summarize the important formulae. Matlab software implementing the methods discussed in this paper is available from xxx.lanl.gov/archive/neuro-sys. ## Appendix A Derivation of Finite Size Correction The following is an outline derivation of the finite size corrections described in section 6. Firstly the characteristic functionals \[Bartlett, 1966\] for the processes $`\overline{N}`$ and the inhomogeneous Poisson process rate $`\lambda (t)`$ are related. $$C_{\overline{N}}(\theta (t))=E\{exp(i_0^T\theta (t)𝑑\overline{N})\}=E_\lambda \{exp(_0^T\lambda (t)b(\theta (t))𝑑t)\}$$ (57) $$b(\theta (t))=exp[i\theta (t)\frac{i}{T}_0^T\theta (t^{})𝑑t^{}]$$ (58) Under the Gaussian process assumption for $`\lambda (t)`$ this integral may be done. $$C_{\overline{N}}=exp[\frac{1}{2}_0^T_0^Tb(t)\mathrm{\Lambda }(t,t^{})b(t^{})𝑑t𝑑t^{}+\overline{\lambda }_0^Tb(t)𝑑t]$$ (59) $$\mathrm{\Lambda }(t,t^{})=E_\lambda \{(\lambda (t)\overline{\lambda })(\lambda (t^{})\overline{\lambda })\}$$ (60) Note that $`\overline{\lambda }`$ denotes the mean rate. Taking the log of the characteristic functionals now yields the following relation between the resultant cumulant functionals. $$K_{\overline{N}}=lnE\{exp(i_0^T\theta (t)𝑑\overline{N})\}=\frac{1}{2}_0^T_0^Tb(t)\mathrm{\Lambda }(t,t^{})b(t^{})𝑑t𝑑t^{}+\overline{\lambda }_0^Tb(t)𝑑t$$ (61) Next $`\theta (t)`$ is chosen appropriately and substituted into $`K_{\overline{N}}`$. The form for $`\theta (t)`$ which is required to obtain the covariance of multitaper estimators is; $$i\theta (t)=\theta _1h_k(t)e^{2\pi if_1t}+\theta _2h_k(t)e^{2\pi if_1t}+\theta _3h_k^{}(t)e^{2\pi if_2t}+\theta _4h_k^{}(t)e^{2\pi if_2t}$$ (62) Substituting into the cumulant functional for $`\overline{N}`$ yields; $$K_{\overline{N}}=lnE\{exp(\theta _1J_k^D(f_1)+\theta _2J_k^D(f_1)+\theta _3J_k^{}^D(f_2)+\theta _4J_k^{}^D(f_2))\}$$ (63) Where $`J_k^D`$ is the Fourier transform of the data tapered by a function indexed by $`k`$. Application of the cumulant expansion theorem \[Ma, 1985\] then leads to; $$K_{\overline{N}}=E\{exp(\theta _1J_k^D(f_1)+\theta _2J_k^D(f_1)+\theta _3J_k^{}^D(f_2)+\theta _4J_k^{}^D(f_2))1\}_C$$ (64) This may then be differentiated and set to zero. $$K_{lmno}=\frac{K_{\overline{N}}}{\theta _1^l\theta _2^m\theta _3^n\theta _4^o}\|_{\theta _1=\theta _2=\theta _3=\theta _4=0}=E\{J_k^{Dl}(f_1)J_k^{Dm}(f_1)J_k^{}^{Dn}(f_2)J_k^{}^{Do}(f_2)\}_C$$ (65) Moments of the estimators may be expressed in terms of these cumulant derivatives. The expressions are simplified by the fact that all cumulant derivatives which have indices summing to an odd number are zero because $`\overline{N}`$ is a zero mean process. $$E\{I_k^D(f)\}=K_{1100}$$ (66) $$var\{I^{MT}(f)\}=\frac{1}{K^2}\underset{k=0}{\overset{K1}{}}\underset{k^{}=0}{\overset{K1}{}}cov\{I_k^D(f),I_k^{}^D(f)\}$$ (67) $$cov\{I_k^D(f),I_k^{}^D(f)\}=K_{1010}K_{0101}+K_{1111}+K_{1001}K_{0110}$$ (68) The problem has now been reduced to that of calculating these derivatives within the model. This is done by substituting the expression for $`\theta (t)`$ into the RHS of equation 61. Considerable algebra then leads to the following exact result. $$K_{lmno}=K_{lmno}^A+K_{lmno}^B$$ (69) Where, $`K_{lmno}^A={\displaystyle \frac{1}{2}}{\displaystyle \underset{l_i,m_i,n_i,o_i}{}}{\displaystyle \frac{l!m!n!o!}{\mathrm{\Pi }l_i!\mathrm{\Pi }m_i!\mathrm{\Pi }n_i!\mathrm{\Pi }o_i!}}\left[{\displaystyle \frac{H_1(f_1)}{T}}\right]^{l_2+l_4}\left[{\displaystyle \frac{H_1(f_1)^{}}{T}}\right]^{m_2+m_4}\mathrm{}`$ $`\left[{\displaystyle \frac{H_1(f_2)}{T}}\right]^{n_2+n_4}\left[{\displaystyle \frac{H_1(f_2)^{}}{T}}\right]^{o_2+o_4}I_{l_3,m_3,n_3,o_3}^{l_1,m_1,n_1,o_1}`$ (70) Where $`_il_i=l`$ and cases where $`l_1+l_2=l`$ or $`l_3+l_4=l`$ are excluded (and also for $`n,m,o`$). $`I_{l_3,m_3,n_3,o_3}^{l_1,m_1,n_1,o_1}={\displaystyle _{\mathrm{}}^{\mathrm{}}}S_\lambda (f)H_{l_1+m_1+n_1+o_1}[f_1(l_1m_1)+f_2(n_1o_1)f]\mathrm{}`$ $`H_{l_3+m_3+n_3+o_3}^{}[f_1(l_3m_3)+f_2(n_3o_3)f]df`$ (71) Where $`S_\lambda (f)`$ is the spectrum of the Gaussian process and $`H_l`$ is; $`H_l(f)={\displaystyle _{\mathrm{}}^{\mathrm{}}}h^l(t)exp(2\pi ift)𝑑t`$ (72) $`H_0(f)=Texp(i\pi fT)sinc(\pi fT)`$ (73) $`K_{lmno}^B=\overline{\lambda }{\displaystyle \underset{p=0}{\overset{l}{}}}{\displaystyle \underset{q=0}{\overset{m}{}}}{\displaystyle \underset{r=0}{\overset{n}{}}}{\displaystyle \underset{s=0}{\overset{o}{}}}[\begin{array}{c}l\hfill \\ p\hfill \end{array}][\begin{array}{c}m\hfill \\ q\hfill \end{array}][\begin{array}{c}n\hfill \\ r\hfill \end{array}][\begin{array}{c}o\hfill \\ s\hfill \end{array}]H_{p+q+r+s}[f_1(pq)+f_2(rs)]\mathrm{}`$ (82) $`\left[{\displaystyle \frac{H_1(f_1)}{T}}\right]^{(lp)}\left[{\displaystyle \frac{H_1(f_1)^{}}{T}}\right]^{(mq)}\left[{\displaystyle \frac{H_1(f_2)}{T}}\right]^{(nr)}\left[{\displaystyle \frac{H_1(f_2)^{}}{T}}\right]^{(os)}`$ (83) The preceding result is somewhat cumbersome but readily evaluated computationally for a given spectrum. The expression simplifies greatly when only frequencies above the smoothing width are considered and many of the terms may be neglected. Restricting attention to the second order properties there are only a few remaining dominant terms. Terms from $`K_{1001}`$ lead to the previously discussed asymptotic results but there are corrections which arise from the term $`K_{1111}`$. Assuming that the population spectrum varies slowly over the width of the tapers leads to the result given by equations 38 \- 43. The validity of this assumption has been tested computationally and was found to be very accurate even for spectra with sharp peaks. ## Acknowledgment The authors thank C. Buneo, and R. Sachdev for providing example datasets, Clive Lauder for help with the calculation of global error bars and D. R. Brillinger and D. J. Thomson for comments which substantially improved the manuscript. M. Jarvis is grateful to R. A. Andersen both for his continued support of theoretical work in his lab and also for his careful reading of the manuscript. M. Jarvis acknowledges the generous support of the Sloan foundation for theoretical neuroscience.
warning/0002/hep-ph0002302.html	ar5iv	text	# I Introduction ## I Introduction Some time ago, after QCD started to become established as the theory of hadronic interactions, a number of authors looked into the possibility of measuring $`T`$-odd effects in current-induced interaction The phrase $`T`$-odd observables refers to observables that change sign under simultaneous reflection of particle momenta and spins and does not refer to truly $`T`$-violating observables. that would result from QCD rescattering effects. The rescattering effects were calculated from the absorptive parts of the relevant next-to-leading order QCD one-loop contributions. The authors of considered $`T`$-odd effects in the decay of a $`J^{PC}=1^{}`$ quarkonium state into three gluonic jets. $`T`$-odd effects in $`e^+e^{}`$-annihilation into three partonic jets were considered in excepting quark loop contributions. First, it came as a surprise that, for mass zero quarks, there are no leading order $`𝒪(\alpha _s^2)`$ $`T`$-odd effects in this reaction . This was understood more systematically later on in from the observation that the absorptive parts are necessarily proportional to the Born term and are thus unobservable. Measurable $`T`$-odd effects in these reactions are generated by quark mass effects which were investigated in . The non-observability of $`𝒪(\alpha _s^2)`$ $`T`$-odd effects in $`e^+e^{}`$-interactions for the massless case does not carry over to the crossed channels of deep inelastic scattering (DIS) and the Drell-Yan (DY) process. $`T`$-odd effects in DIS were explored in and in the DY process in . Since the early proposals to measure $`T`$-odd effects in current induced interactions in the early eighties experimental facilities and techniques have considerably been improved. Luminosities of lepton-hadron and hadron-hadron colliders have dramatically increased providing for much higher event rates than was possible in the earlier experiments. The energy range of the collliders has been extended such that high momentum transfers can now be routinely probed. For example, at HERA one is starting to probe weak interaction $`Z`$-exchange effects in neutral current events at very high momentum transfers. This opens the door for the investigation of $`T`$\- and $`P`$-odd effects in neutral current DIS. Powerful jet finding and flavour tagging algorithms have been developed that allow one to define asymmetry measures related to $`T`$-odd effects in DIS that use parton jet observables instead of the semi-inclusive particle observables used in the calculation of . Finally, there have been dramatic improvements in the availability of polarized beams which again can be utilized to define new $`T`$-odd observables . It is therefore timely to take a fresh look at the subject of $`T`$-odd observables in current- induced reactions generated by QCD rescattering effects or, in a different language, by the absorptive parts of the corresponding one-loop contributions. In this paper we point out that the $`𝒪(\alpha _s^2)`$ absorptive parts of the relevant one-loop contributions in DIS and in the DY process can be obtained through crossing from the well-known one-loop contributions to $`e^+e^{}q\overline{q}g`$ annihilation calculated in . This is theoretically appealing and provides an independent check of the results presented before in DIS and DY . We also fill out some small odds and ends on the subject of $`T`$-odd observables in these reactions which had not been covered in the earlier publications. Our paper is structured as follows. In Sec. 2 we derive crossing rules that allow one to cross from the $`e^+e^{}`$-annihilation channel to the DIS and DY channels. To obtain the absorptive parts in the respective channels it is necessary to discuss the analyticity structure of the one-loop contributions in the complex plane of the relevant kinematical variables. The absorptive parts originate from logarithmic and dilogarithmic functions in the one-loop amplitude when their arguments take values on cuts in the analytic plane. We identify the range of values of the kinematical variables in the three processes and show how to analytically continue the one-loop functions from the $`e^+e^{}`$-channel to the DIS and DY channels. Sec. 3 is devoted to a detailed discussion of $`T`$-odd effects in DIS. We first provide a complete list of the nonsingular one-loop contributions in $`e^+e^{}`$-annihilation. Using the crossing rules laid down in Sec. 2 we analytically continue the one-loop amplitudes to DIS. The absorptive one-loop amplitudes are then folded with the Born term contributions. There are three $`T`$-odd hadronic structure functions, $`H_5`$, $`H_8`$, and $`H_9`$, whose functional form is given for the quark, antiquark and gluon initiated cases. We then define helicity structure functions which appear as angular coefficients in the angular decay distribution of the DIS process when the hadronic tensor is folded with the leptonic tensor. This allows us to compare our results with the results of . We find complete agreement with the $`T`$-odd results presented in these papers. In Sec. 4 we do the crossing and the analytical continuation to the DY process. Again we find agreement of our results for the $`T`$-odd structures with those given in after corrections for a typographical error reported in . In Sec. 5 we give our summary and provide an outlook on possible further applications of our results to spin-dependent $`T`$-odd observables. In an Appendix we present results on the dispersive and absorptive contributions of the triangle anomaly graph to the DIS process. ## II General principles of analyticity and crossing In this chapter we will develop the framework necessary to obtain the $`𝒪(\alpha _s^2)`$ one-loop corrections in the DIS and DY channels from the known results in the $`e^+e^{}`$-annihilation channel . In particular, we will derive crossing rules that allow one to determine the whole set of invariant hadronic structure functions $`H_i`$ (i=1,…,9), including the absorptive $`T`$-odd structure functions for the DIS and the DY processes. In the absence of polarisation the definition of $`T`$-odd observables in current-induced interactions involves the analysis of parton processes with at least three partons necessary to form triple momentum products. In order to fix our notation let us write down the momentum configuration for $`e^+e^{}`$ annihilation into a quark, antiquark and a gluon: $$l^{}(k)+l^+(k^{})\stackrel{\gamma ,Z}{}q(p_1)+\overline{q}(p_2)+g(p_3)$$ (1) where $`l^{}`$ and $`l^+`$ are massless leptons, $`q`$ and $`\overline{q}`$ are massless quarks and antiquarks, respectively, and $`g`$ is a bremsstrahlung gluon. The momentum of the time-like gauge boson is determined from four-momentum conservation and is given by $`q=p_1+p_2+p_3`$. The leading order contributions to the $`T`$-odd observables come from the interference of the absorptive parts of the one-loop amplitudes and the Born term amplitude. In Fig. 1 we show the $`𝒪(\alpha _s^2)`$ one-loop diagrams that contribute to $`e^+e^{}q\overline{q}g`$ (with the leptonic part omitted). They divide into the eleven contributions without quark loops and the two diagrams with a quark loop. In the main part of this paper we will be mostly concerned with the first eleven non-quark loop contributions (a) - (k). A discussion of the so-called triangle-anomaly quark loop contributions (l) and (m) is deferred to the Appendix. At the three-parton level DIS proceeds through the following three subprocesses: $`l(k)+q(p)`$ $`\stackrel{\gamma ,Z,W^\pm }{}`$ $`l^{}(k^{})+q(p^{})+g(p_3),`$ (2) $`l(k)+\overline{q}(p)`$ $`\stackrel{\gamma ,Z,W^\pm }{}`$ $`l^{}(k^{})+\overline{q}(p^{})+g(p_3),`$ (3) $`l(k)+g(p)`$ $`\stackrel{\gamma ,Z,W^\pm }{}`$ $`l^{}(k^{})+q(p^{})+\overline{q}(p_3).`$ (4) They are referred to as the quark-, antiquark- and gluon-initiated DIS processes, respectively. The momentum of the space-like gauge boson is now given by $`q=p+p^{}+p_3`$. In the DY process the next-to-leading order contributions come from the following subprocesses: $`q(p_a)+\overline{q}(p_b)`$ $`\stackrel{\gamma ,Z,W^\pm }{}`$ $`l(k)+l^{}(k^{})+g(p_3),`$ (5) $`q(p_a)+g(p_b)`$ $`\stackrel{\gamma ,Z,W^\pm }{}`$ $`l(k)+l^{}(k^{})+q(p_3),`$ (6) $`\overline{q}(p_a)+g(p_b)`$ $`\stackrel{\gamma ,Z,W^\pm }{}`$ $`l(k)+l^{}(k^{})+\overline{q}(p_3),`$ (7) where (5) is the annihilation subprocess and (6) and (7) are the so-called quark- and antiquark-initiated ”Compton” subprocesses, respectively. The momentum of the time-like gauge boson is given by $`q=p_a+p_bp_3`$. Differing from $`e^+e^{}`$-annihilation there are also charged gauge boson contributions to DIS and the DY process. In what follows we need to discuss only the hadronic part of the three-parton processes listed in (2) and (5-7). The contraction with the leptonic part will lead to angular factors and some $`y`$-dependence in the DIS case. The contraction with the leptonic tensor will be discussed in the subsequent sections when we compare our results with the calculation of Hagiwara et al. . The relevant one-loop contributions to DIS and the DY process can be obtained from the one-loop contibutions to $`e^+e^{}`$-annihilation calculated in through crossing, i.e. through the exchange of incoming and outgoing particle momenta in the one-loop diagrams in Fig.1. For the real parts of the one-loop contributions and for the Born term contribution crossing can be implemented in a straightforward manner. Crossing is more subtle for the imaginary parts of the one-loop amplitudes and needs a careful discussion of the analyticity properties of the one-loop amplitude. The crossing of external lines in Feynman diagrams implies a sign change of the four momentum associated with that line. Thereby, the values of the kinematic variables associated with the respective momentum undergoes a discontinuous change. Massless one-loop amplitudes contain log and dilog functions which depend on these kinematic variables and which may be indefinite in certain ranges of their domains of definition, i.e. they may be multivalued. One can choose among the possible values by defining the value of the function at a given point. Starting from this point one determines the value of the multivalued function on the cuts by analytic continuation. The kinematic variables are taken to be complex in this procedure. In order to obtain a smooth continuation one makes use of the imaginary parts of the one-loop amplitudes given by the ($`i\tau `$)-form of the relevant propagators. In this way one avoids possible ambiguities. As an example we take the natural logarithm. The logarithm is taken as a complex-valued analytic function with a cut on the negative real axis and a branch point located at zero. Next consider the natural logarithm of an arbitrary positive real number $`x`$. To obtain its value at ($`x`$) we use $`\mathrm{ln}(x)=\mathrm{ln}(\|x\|)+i(2k+1)\pi ,\mathrm{with}k𝒵.`$ The integer number $`k`$ will be determined from the phase angle of the complex number ($`x`$). If one excludes multiple rotations in the complex plane then one remains with only two possibilities: $`k=0,1`$. For all complex numbers of the form ($`x\pm i\tau `$) with an infinitesimally small positive $`\tau `$ one would have the following identity: $$\underset{\tau 0}{lim}\mathrm{ln}(x\pm i\tau )=\mathrm{ln}(\|x\|)\pm i\pi .$$ (8) The results of the calculation of the one-loop contributions to $`e^+e^{}`$-annihilation listed in contain no explicit $`i\tau `$ prescription. This is adequate for the $`e^+e^{}`$-reaction since in this case the results are given for regions in the complex plane away from the singularities. The results are valid only in this restricted region and need to be analytically continued to the other regions in the complex plane accessible in DIS and in the DY case. It is, however, possible to restore the omitted $`i\tau `$ prescriptions in in a straightforward way. For the relevant kinematic variables the infinitesimal imaginary parts are provided by the Feynman rules if one takes the full propagators in their original ($`i\tau `$)-form. In this case one has the following terms from solving the loop integrals $`(s_{ij}(p_i+p_j)^2=2p_ip_j)`$: $$\frac{s_{ij}+i\tau }{q^2+i\tau },1\frac{s_{ij}+i\tau }{q^2+i\tau },\mathrm{and}(q^2i\tau )^\epsilon ,$$ (9) where $`\epsilon =2d/2`$, and $`d`$ is the continuous space-time dimension. Here one should notice that the form of energy-momentum conservation for the s-channel annihilation ensures relative plus signs between the three scalar invariants $`s_{ij}=2p_ip_j`$ and $`i\tau `$, as well as between $`q^2`$ and $`i\tau `$ in the denominators of the respective Feynman integrals. Thus, more generally, one has the following rules for s-channel annihilation: $$q^2q^2+i\tau ,s_{ij}s_{ij}+i\tau .$$ (10) With these rules the results of are valid in any kinematical region. For the kinematical variables $`y_{ij}=s_{ij}/q^2`$ used in one finds the following replacements: $`y_{ij}`$ $``$ $`y_{ij}+{\displaystyle \frac{i\tau }{q^2}}(1y_{ij})+𝒪(\tau ^2),`$ (11) $`1y_{ij}`$ $``$ $`1y_{ij}{\displaystyle \frac{i\tau }{q^2}}(1y_{ij})+𝒪(\tau ^2),`$ (12) $`(q^2)^\epsilon `$ $``$ $`1\epsilon \mathrm{ln}(q^2i\tau )+𝒪(\epsilon ^2).`$ (13) For every contributing subprocess in DIS and DY one has to perform a detailed investigation of the range of values of the $`y_{ij}`$’s after crossing and then one can analytically continue the log and dilog functions and thereby remove the ambiguity which occurs when one changes the sign of their arguments. The analytic continuation of the logarithm function is given in (8). For the dilogarithms, when $`x>1`$, we use the identity $$\mathrm{Li}_2(x)=\mathrm{Li}_2(1x)+\zeta (2)\mathrm{ln}(x)\mathrm{ln}(1x),$$ (14) and treat the complex logarithm $`\mathrm{ln}(1x)`$ according to (8). At this point we introduce the usual hadronic DIS variables $`x`$ and $`z`$ $$x=\frac{q^2}{2qp},z=\frac{pp^{}}{qp},$$ (15) and proceed with the crossing procedure as described above. The crossing from the $`e^+e^{}`$-channel to the quark-initiated subprocesses in DIS is given by the following change of the momenta $$p_1p^{}\mathrm{and}p_2p,$$ (16) according to the momentum definitions in (1) and (2). The change of the kinematical variables $`y_{ij}`$’s and the new ranges of their values are given by $`y_{12}{\displaystyle \frac{2p_1p_2}{q^2}}y_{12}^q{\displaystyle \frac{2pp^{}}{q^2}}`$ $`=`$ $`{\displaystyle \frac{z}{x}}[1,\mathrm{}],`$ (17) $`y_{13}{\displaystyle \frac{2p_1p_3}{q^2}}y_{13}^q{\displaystyle \frac{2p^{}p_3}{q^2}}`$ $`=`$ $`1{\displaystyle \frac{1}{x}}[\mathrm{},0],`$ (18) $`y_{23}{\displaystyle \frac{2p_2p_3}{q^2}}y_{23}^q{\displaystyle \frac{2pp_3}{q^2}}`$ $`=`$ $`{\displaystyle \frac{1z}{x}}[1,\mathrm{}].`$ (19) One should note that $`q^2<0`$ in the crossed DIS channel. The corresponding crossing to the antiquark-initiated subprocess involves the momentum changes $`p_1p`$ and $`p_2p^{}`$. However, we need not explicitly discuss crossing for this case since one can use $`CP`$-invariance in the final result for the quark-initiated case to obtain the corresponding antiquark-initiated results. Similarly, the crossing from the $`e^+e^{}`$-annihilation to the gluon-initiated subprocess in DIS is effected by $$p_1p^{}\mathrm{and}p_3p.$$ (20) The resulting $`y_{ij}^g`$’s are: $`y_{12}y_{12}^g{\displaystyle \frac{2p^{}p_2}{q^2}}`$ $`=`$ $`1{\displaystyle \frac{1}{x}}[\mathrm{},0],`$ (21) $`y_{13}y_{13}^g{\displaystyle \frac{2pp^{}}{q^2}}`$ $`=`$ $`{\displaystyle \frac{z}{x}}[1,\mathrm{}],`$ (22) $`y_{23}y_{23}^g{\displaystyle \frac{2pp_2}{q^2}}`$ $`=`$ $`{\displaystyle \frac{1z}{x}}[1,\mathrm{}].`$ (23) The one-loop results in are presented in terms of the variables $`x_k=2p_kq/q^2`$ $`(k=1,2,3)`$. The $`x_k`$ are related to the $`y_{ij}`$ via $`x_k=1y_{ij},(ki,j).`$ Next we turn to the Drell-Yan process. For the annihilation subprocess (5) crossing implies the following change of momenta $$p_1p_a\mathrm{and}p_2p_b.$$ (24) For the quark-initiated ”Compton” scattering one has to change $$p_1p_a\mathrm{and}p_3p_b,$$ (25) with $`qq`$ in both cases. Again we omit explicit reference to the antiquark-initiated ”Compton” scattering case because its structure follows from the quark-initiated ”Compton” scattering case through $`CP`$ invariance. The corresponding kinematical variables in the annihilation subprocesses are now $`y_{12}y_{12}^a{\displaystyle \frac{2p_ap_b}{q^2}}`$ $`=`$ $`{\displaystyle \frac{1}{x_a}}+{\displaystyle \frac{1}{x_b}}1[1,\mathrm{}],`$ (26) $`y_{13}y_{13}^a{\displaystyle \frac{2p_ap_3}{q^2}}`$ $`=`$ $`1{\displaystyle \frac{1}{x_b}}[\mathrm{},0],`$ (27) $`y_{23}y_{23}^a{\displaystyle \frac{2p_bp_3}{q^2}}`$ $`=`$ $`1{\displaystyle \frac{1}{x_a}}[\mathrm{},0].`$ (28) The crossed $`y_{ij}`$’s for the quark-initiated ”Compton” subprocesses are $`y_{12}y_{12}^{C_q}{\displaystyle \frac{2p_ap_3}{q^2}}`$ $`=`$ $`1{\displaystyle \frac{1}{x_b}}[\mathrm{},0],`$ (29) $`y_{13}y_{13}^{C_q}{\displaystyle \frac{2p_ap_b}{q^2}}`$ $`=`$ $`{\displaystyle \frac{1}{x_a}}+{\displaystyle \frac{1}{x_b}}1[1,\mathrm{}],`$ (30) $`y_{23}y_{23}^{C_q}{\displaystyle \frac{2p_bp_3}{q^2}}`$ $`=`$ $`1{\displaystyle \frac{1}{x_a}}[\mathrm{},0].`$ (31) In the above equation we have introduced the DY variables $`x_a`$ and $`x_b`$ defined by $$x_a=\frac{M^2}{2qp_a},x_b=\frac{M^2}{2qp_b},$$ (32) where $`M^2`$ is the mass of the exchanged gauge bosons $`W^\pm `$ and $`Z`$, or, for the electromagnetic interaction, one has $`M^2=q^2`$. Note the $`23`$ symmetry between the annihilation and the quark-initiated Compton subprocess in terms of the subprocess variables $`x_a`$ and $`x_b`$. The momentum changes involved in crossing from the $`e^+e^{}`$-channel to DIS and the DY channel are summarized in Fig. 2 using diagram 1(f) as an illustrative example. As we have already mentioned before, the results of the calculation in $`e^+e^{}`$-annihilation channel contain log and dilog functions, which, after applying the rules from (17-21) and from (26-29) have to be analytically continued to the new kinematically allowed regions. Using the cut structure defined in (8) and (14), one obtains their absorptive parts. The results of crossing and analytic continuation of logarithmic and dilogarithmic functions appearing in our calculation of $`T`$-odd amplitudes in DIS and DY channels are summarized in Table I. ## III Crossing Results for Hadronic Structure Functions in DIS We now proceed to derive the $`T`$-odd hadronic structure functions in DIS by applying the analyticity and crossing rules derived in the proceeding section starting with the $`e^+e^{}`$ results in . We will be mostly interested in the imaginary (absorptive) parts but will also briefly comment on the crossing properties of the real parts of the one-loop amplitudes and the Born term amplitudes. The results of the crossing procedure will then be compared to the corresponding results in DIS in the three quark, antiquark and gluon-initiated cases . As was mentioned in the introduction, the $`T`$-odd structure functions in the $`e^+e^{}`$ case vanish identically at the one-loop level for the set of graphs (1a)-(1k) with no quark loopsThere are contributions to the $`T`$-odd structure functions coming from the quark loop graphs (l) and (m) in Fig. 1 even for massless external quarks due to an incomplete cancellation of the $`b`$\- and $`t`$-quark in the quark loop. These contributions have been shown to be very small . The reason is that the absorptive parts of the $`e^+e^{}`$-annihilation massless one-loop amplitudes are proportional to the singular terms of its dispersive part which in turn have Born term structure . The absorptive parts are thus not measurable at this order of perturbation theory. However, the kinematics is different in the crossed processes and this proportionality no longer holds leading to nonvanishing $`T`$-odd effects in the crossed channels. We begin our discussion by recapitulating the $`e^+e^{}`$-annihilation one-loop results given in . They are needed as a starting point for the crossing procedure. The transition amplitude $`T_{\mu \beta }^V`$ for the vector current transition $`V_\mu q\overline{q}g`$ can be expanded in $`d`$-dimensions along the seven independent covariants $`T_{\mu \beta }^V`$ $`=`$ $`N_iC_{\mu \beta }^i(i=1,\mathrm{},7)`$ (33) $`=`$ $`q^4N_1\widehat{p}_{+\mu }\stackrel{~}{p}_\beta \mathit{}_3+q^4N_2\widehat{p}_\mu \stackrel{~}{p}_\beta \mathit{}_3`$ (37) $`+q^2N_3\widehat{\gamma }_\mu \stackrel{~}{p}_\beta +q^2N_4\widehat{p}_{+\mu }\stackrel{~}{\gamma }_\beta `$ $`+q^2N_5\widehat{p}_\mu \stackrel{~}{\gamma }_\beta +q^2N_6\widehat{\stackrel{~}{g}}_{\mu \beta }\mathit{}_3`$ $`+N_7\left(\gamma _\mu {\displaystyle \frac{\mathit{}_2+\mathit{}_3}{s_{23}}}\gamma _\beta \gamma _\beta {\displaystyle \frac{\mathit{}_1+\mathit{}_3}{s_{13}}}\gamma _\mu \right)`$ where $`q=p_1+p_2+p_3`$ and $`p_\pm =p_1\pm p_2`$. The covariants $`C_{\mu \beta }^i`$ are defined through Eq. (33) and the symbols “$``$” and “$``$” denote the gauge invariant completions $`\widehat{p}_\mu `$ $`=`$ $`p_\mu {\displaystyle \frac{pq}{q^2}}q_\mu ,`$ (38) $`\widehat{\gamma }_\mu `$ $`=`$ $`\gamma _\mu {\displaystyle \frac{\mathit{}}{q^2}}q_\mu ,`$ (39) $`\stackrel{~}{p}_\beta `$ $`=`$ $`p_\beta {\displaystyle \frac{pp_3}{p_3q}}q_\beta ,`$ (40) $`\stackrel{~}{\gamma }_\beta `$ $`=`$ $`\gamma _\beta {\displaystyle \frac{\mathit{}_3}{p_3q}}q_\beta ,`$ (41) $`\widehat{\stackrel{~}{g}}_{\mu \beta }`$ $`=`$ $`g_{\mu \beta }{\displaystyle \frac{q_\beta p_{3\mu }}{p_3q}}.`$ (42) In writing down the covariant expansion it is understood that the covariants are taken between the relevant spin wave functions, i.e. $`\overline{u}_1T_{\mu \beta }^Vv_2\epsilon ^\beta `$. We emphasize that the $`T`$-odd structure functions resulting from the one-loop amplitudes are infrared (IR) and collinear (M) finite in $`d=4`$ dimensions. This follows indirectly from the Lee-Nauenberg theorem in that there are no corresponding tree graph contributions that could cancel the divergencies if there were any. We could in principle therefore keep $`d=4`$ in our calculation. In this case one has overcounted the number of covariants in the above expansion. There are in fact only six independent covariants in four dimensions as can be verified by counting the number of independent helicity amplitudes. For the sake of completeness we list the linear relation between the seven covariants $`C_{\mu \beta }^i`$ in $`d=4`$ dimensions (taken between spin wave functions) which can be obtained fromWe take this opportunity to correct Eq. (A.7) in . Eq. (A.7) in should read: $`(1x_3)k_{\beta \mu }^7=k_{\beta \mu }^1\frac{1}{2}x_3(k_{\beta \mu }^3k_{\beta \mu }^5)\frac{1}{2}(2x_1x_3)k_{\beta \mu }^4+(1x_3)k_{\beta \mu }^6`$. . One has $`(1x_1)(1x_2)(1x_3)C_{\mu \beta }^7`$ $`=`$ $`x_3C_{\mu \beta }^1{\displaystyle \frac{1}{2}}x_3(x_1+x_2)C_{\mu \beta }^3{\displaystyle \frac{1}{2}}(x_1^2x_2^2)C_{\mu \beta }^4`$ (44) $`+{\displaystyle \frac{1}{2}}x_3^2C_{\mu \beta }^5(x_1x_2)(1x_3)C_{\mu \beta }^6.`$ As concerns the present application it is nevertheless technically advantageous to work with the (overcounted) set of the above seven covariants. Note that the seventh covariant $`C_{\mu \beta }^7`$ in (33) has been chosen to have Born term structure. This will be important to keep in mind in what follows. As we are dealing with massless quarks the case of axial vector current transition $`A_\mu q\overline{q}g`$ can be easily dealt with. We simply have to multiply the vector current amplitude by $`\gamma _5`$ from the left. One has $$T_{\mu \beta }^A=\gamma _5T_{\mu \beta }^V.$$ (45) What is needed are explicit forms of the one-loop amplitudes in the $`e^+e^{}`$ channel. The IR and M singular one-loop contributions have Born term structure and will not be needed in the following. As concerns the nonsingular one-loop contributions, we decided to reproduce the two relevant tables, Table II and Table III, from because the results of may not be readily available to everyone (in those early days there was no electronic publishing). There are two colour structures in the one-loop amplitudes referred to as the QCD and QED structure. The corresponding amplitudes are denoted by $`\stackrel{~}{N}_i`$ (Table II) and $`\widehat{N}_i`$ (Table III), respectively, $`N_l^{(j)}=g^3\{\stackrel{~}{N}_lif^{ijk}t_kt_i+\widehat{N}_lt_it_jt_i\}.`$ As a next step one folds the above one-loop amplitude with the Born term amplitude proportional to $`C_{\mu \beta }^7`$ and sums over the spins and colors of the final partons. The result can be expanded in terms of nine gauge invariant covariants which define the nine invariant structure functions $`H_i(i=1,\mathrm{},9)`$ of the $`e^+e^{}`$-annihilation process. $`H_{\mu \nu }`$ $`=`$ $`H_1\left(g_{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right)`$ (54) $`+H_2q^2\widehat{p}_{1\mu }\widehat{p}_{1\nu }`$ $`+H_3q^2\widehat{p}_{2\mu }\widehat{p}_{2\nu }`$ $`+H_4q^2(\widehat{p}_{1\mu }\widehat{p}_{2\nu }+\widehat{p}_{1\nu }\widehat{p}_{2\mu })`$ $`+H_5q^2(\widehat{p}_{1\mu }\widehat{p}_{2\nu }\widehat{p}_{1\nu }\widehat{p}_{2\mu })`$ $`+H_6q^2iϵ_{\mu \nu \alpha \beta }q^\alpha p_1^\beta `$ $`+H_7q^2iϵ_{\mu \nu \alpha \beta }q^\alpha p_2^\beta `$ $`+H_8q^4(\widehat{p}_{1\mu }F_\nu +\widehat{p}_{1\nu }F_\mu )`$ $`+H_9q^4(\widehat{p}_{2\mu }F_\nu +\widehat{p}_{2\nu }F_\mu ),`$ where $`F_\mu =iϵ_{\mu \alpha \beta \gamma }p_1^\alpha p_2^\beta q^\gamma `$. Throughout this paper we use the convention $`ϵ^{0123}=1`$. From the hermiticity property $`H_{\mu \nu }=H_{\nu \mu }^{}`$ one concludes that $`H_1,H_2,H_3,H_4,H_6`$ and $`H_7`$ are purely real and $`H_5,H_8`$ and $`H_9`$ are purely imaginary functions. Also, the first five structure functions $`H_1H_5`$ are parity conserving and $`H_6H_9`$ are parity violating. As was mentioned before the $`T`$-odd invariant structures $`H_5,H_8`$ and $`H_9`$ vanish in $`e^+e^{}`$-annihilation at this order of perturbation theory for massless quarks. In the case of DIS the corresponding expansion into a set of nine covariants reads $`H_{\mu \nu }`$ $`=`$ $`H_1\left(g_{\mu \nu }{\displaystyle \frac{q_\mu q_\nu }{q^2}}\right)`$ (63) $`+H_2q^2\widehat{p}_\mu \widehat{p}_\nu `$ $`+H_3q^2\widehat{p}_\mu ^{}\widehat{p}_\nu ^{}`$ $`+H_4q^2(\widehat{p}_\mu \widehat{p}_\nu ^{}+\widehat{p}_\nu \widehat{p}_\mu ^{})`$ $`+H_5q^2(\widehat{p}_\mu \widehat{p}_\nu ^{}\widehat{p}_\nu \widehat{p}_\mu ^{})`$ $`+H_6q^2iϵ_{\mu \nu \alpha \beta }q^\alpha p^\beta `$ $`+H_7q^2iϵ_{\mu \nu \alpha \beta }q^\alpha p^\beta `$ $`+H_8q^4(\widehat{p}_\mu \stackrel{~}{F}_\nu +\widehat{p}_\nu \stackrel{~}{F}_\mu )`$ $`+H_9q^4(\widehat{p}_\mu ^{}\stackrel{~}{F}_\nu +\widehat{p}_\nu ^{}\stackrel{~}{F}_\mu ),`$ where now $`\stackrel{~}{F}_\mu =iϵ_{\mu \alpha \beta \gamma }p^\alpha p^\beta q^\gamma `$. It is clear that the actual form of the parton level tensor $`H_{\mu \nu }`$ and the hard scattering structure functions $`H_i`$ depend on which of the three quark-, antiquark- and gluon-initiated cases are being discussed. If necessary, we shall affix additional superfixes $`q,\overline{q}`$ and $`g`$, respectively, to the parton level structure functions in order to differentiate between the three cases. For the Born term contribution and for the real parts of the one-loop contributions the crossing can be done directly on the structure function expansion (54) where one may have a reshuffling of covariants in (63) depending on whether one is crossing to the quark-, antiquark- or the gluon-initiated caseThe relevant tables for the real structure functions $`H_1,H_2,H_3,H_4,H_6`$ and $`H_7`$ in $`e^+e^{}`$-annihilation can be found in .. In order to determine the absorptive contributions to the $`T`$-odd structure functions $`H_5,H_8`$ and $`H_9`$ one has to go back to the $`e^+e^{}`$ one-loop amplitude expressions in Table II and Table III, do the crossing, extract the imaginary parts and finally fold them with the Born term amplitude. It is clear that $`Im(N_7)`$ does not contribute to the $`T`$-odd structure functions in this sequence of steps after folding with the Born term amplitude since it is proportional to the Born term itself. Our final expressions for all structure functions include spin and color averaging factors for partons in the initial states. As indicated in Fig. 2 the crossing of a fermion line brings in an overall minus sign. With all this we are now in a position to write down the $`T`$-odd hard scattering structure functions $`H_5^q,H_8^q`$ and $`H_9^q`$ referring to the quark-initiated case of DIS. One uses the rules (17) and the results of Tables I, II and III, together with the change of momenta defined in (16) and the expansion (63). Note that the final antifermion line transforms to an initial fermion line which implies an overall minus sign (see also Fig. 2). Averaging over the spin and color of the initial quark and summing over the spins and colors of the final quarks and gluons one obtains the following results for the $`T`$-odd structure functions in the quark-initiated case of DIS: $`i\pi H_5^q`$ $`=`$ $`g^4x^2[{\displaystyle \frac{C_FN_C}{2}}(1{\displaystyle \frac{2xz}{(1x)(1z)}})C_F(C_F{\displaystyle \frac{N_C}{2}})(1+{\displaystyle \frac{2z}{(1x)(1z)}}`$ (65) $`+\mathrm{ln}(z){\displaystyle \frac{2z}{(1x)(1z)^2}})],`$ $`i\pi H_8^q`$ $`=`$ $`g^4{\displaystyle \frac{x^3(1xz)}{(1x)(1z)}}[{\displaystyle \frac{C_FN_C}{2}}+C_F(C_F{\displaystyle \frac{N_C}{2}})({\displaystyle \frac{3z}{1z}}+\mathrm{ln}(z){\displaystyle \frac{2}{(1z)^2}})],`$ (66) $`i\pi H_9^q`$ $`=`$ $`g^4{\displaystyle \frac{x^3}{(1x)(1z)}}[{\displaystyle \frac{C_FN_C}{2}}3+C_F(C_F{\displaystyle \frac{N_C}{2}})({\displaystyle \frac{13z}{1z}}+\mathrm{ln}(z){\displaystyle \frac{2(12z)}{(1z)^2}})],`$ (67) where $`x`$ and $`z`$ are defined in (15) and the color factors are given by $`C_F=4/3`$ and $`N_C=3`$. The antiquark-initiated absorptive structure functions $`H_i^{\overline{q}}`$ ($`i=5,8,9`$) can be obtained from the CP-invariance of the reaction. They read $$H_5^{\overline{q}}=H_5^q,H_8^{\overline{q}}=H_8^q,H_9^{\overline{q}}=H_9^q.$$ (68) In the gluon-initiated case one uses a corresponding crossing and colour averaging procedure and obtains $`i\pi H_5^g`$ $`=`$ $`g^4{\displaystyle \frac{x^2}{z(1z)}}[x2xz\mathrm{ln}(z){\displaystyle \frac{1x}{1z}}+\mathrm{ln}(1z){\displaystyle \frac{1x}{z}}](C_F{\displaystyle \frac{N_C}{2}}),`$ (69) $`i\pi H_8^g`$ $`=`$ $`g^4{\displaystyle \frac{x^3}{z(1z)}}[2(1x){\displaystyle \frac{x+z1}{z(1z)}}\mathrm{ln}(z){\displaystyle \frac{x+z1}{(1z)^2}}\mathrm{ln}(1z){\displaystyle \frac{xz1}{z^2}}](C_F{\displaystyle \frac{N_C}{2}}),`$ (70) $`i\pi H_9^g`$ $`=`$ $`g^4{\displaystyle \frac{x^3}{z(1z)}}[{\displaystyle \frac{12z}{z(1z)}}\mathrm{ln}(z){\displaystyle \frac{1}{(1z)^2}}+\mathrm{ln}(1z){\displaystyle \frac{1}{z^2}}](C_F{\displaystyle \frac{N_C}{2}}).`$ (71) As it can be noticed from (69) the QCD like contributions resulting from the 3-gluon coupling vanish dynamically for the gluon-initiated case (see also ). We are now in the position to compare our results with those in Ref. . In order to facilitate the comparison we contract the hadron tensor with the lepton tensor $`L_{\mu \nu }`$. In the case of a pure vector coupling one has $$L^{\mu \nu }=2[k^\mu k^\nu +k^\nu k^\mu +\frac{q^2}{2}g_{\mu \nu }iϵ^{\mu \nu \rho \sigma }k_\rho k_\sigma ^{}]$$ (72) where the upper(lower) sign holds for a purely left-(right-) handed initial lepton. In the case of a left chiral ($`VA`$) axial-vector charged current without lepton polarization the form of (72) stays the same except that one has only the upper (’minus’) sign for the epsilon term. In this paper we are dealing with parton momenta only whereas in the angular decay distributions are written down in terms of hadron momenta. We can nevertheless compare the angular structures of the two approaches because in the parton model the parton momenta are assumed to be collinear with the hadron momenta from which they originate or into which they fragment. This allows us to meaningfully compare the angular structure of the two approaches without having to set up the whole parton model framework with its integrations over distribution and fragmentation functions. Calculating the hard scattering cross section we will have to take into account an overall phase space factor $$PS=\frac{1}{2^7\pi ^5}\frac{1}{(q^2)}$$ (73) that multiplies the structure functions $`H_i`$, together with the appropriate coupling constants. By contracting our parton level tensor $`H_{\mu \nu }`$ with the leptonic tensor we obtain a result which has exactly the same form as the one in that stood for process variables. We find $$q^2L^{\mu \nu }H_{\mu \nu }=\frac{2}{y^2}\{A+B\mathrm{cos}\varphi +C\mathrm{cos}2\varphi +D\mathrm{sin}\varphi +E\mathrm{sin}2\varphi \}$$ (74) with $`A`$ $`=`$ $`2(1y)F_L+[1+(1y)^2]F_T\pm y(2y)F_3,`$ (75) $`B`$ $`=`$ $`\sqrt{1y}[(2y)F_4\pm yF_5],`$ (76) $`C`$ $`=`$ $`(1y)F_6,`$ (77) $`D`$ $`=`$ $`\sqrt{1y}[\pm yF_7+(2y)F_8],`$ (78) $`E`$ $`=`$ $`(1y)F_9,`$ (79) where the $`T`$-odd coefficients $`D`$ and $`E`$ multiply the $`T`$-odd angular correlation factors $`\mathrm{sin}\varphi `$ and $`\mathrm{sin}2\varphi `$, respectively. The upper(lower) sign is for the left-(right-) handed lepton beam. The scaling variable $`y`$ is defined by $`y={\displaystyle \frac{qp}{kp}}`$and is the same for parton and hadron kinematics. The same statement holds true for the azimuthal angle $`\varphi `$ between the lepton scattering plane and the hadron production plane (see Fig.3). As in , we neglect the nonperturbative effects from the intrinsic transverse momentum spread in the nucleon. Intrinsic transverse momentum effects have very little influence on the $`T`$-odd asymmetries . We therefore assume that the azimuthal angle of the produced hadron is identical to the azimuthal angle of the parton from which it originates. From the definition of three-momenta drawn in Fig.3 one has $`ϵ^{\alpha \beta \gamma \delta }k_\alpha p_\beta p_{}^{}{}_{\gamma }{}^{}q_\delta ={\displaystyle \frac{\kappa }{2x}}q^4{\displaystyle \frac{\sqrt{1y}}{y}}\mathrm{sin}\varphi ,`$where $`\kappa =p_T^{}/(q^2)^{1/2}`$. For the quark-initiated case one obtains the following expressions for the angular coefficients $`F_i,i=1,\mathrm{},9`$, in terms of the corresponding nine structure functions: $`F_L`$ $`=`$ $`H_1^q+{\displaystyle \frac{1}{4x^2}}H_2^q+{\displaystyle \frac{1}{4}}\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)^2H_3^q+{\displaystyle \frac{1}{2x}}\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)H_4^q,`$ (80) $`F_T`$ $`=`$ $`H_1^q+{\displaystyle \frac{\kappa ^2}{2}}H_3^q,`$ (81) $`F_3`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)H_7^q+{\displaystyle \frac{1}{x}}H_6^q\right],`$ (82) $`F_4`$ $`=`$ $`\kappa \left[\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)H_3^q+{\displaystyle \frac{1}{x}}H_4^q\right],`$ (83) $`F_5`$ $`=`$ $`2\kappa H_7^q`$ (84) $`F_6`$ $`=`$ $`\kappa ^2H_3^q,`$ (85) $`F_7`$ $`=`$ $`i{\displaystyle \frac{\kappa }{x}}H_5^q,`$ (86) $`F_8`$ $`=`$ $`i{\displaystyle \frac{\kappa }{2x}}\left[{\displaystyle \frac{1}{x}}H_8^q+\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)H_9^q\right],`$ (87) $`F_9`$ $`=`$ $`i{\displaystyle \frac{\kappa ^2}{x}}H_9^q.`$ (88) The relations (80) agree with the corresponding relations in when $`H_6^q`$ and $`H_7^q`$ are expressed by the $`H_9^{HHK},H_{10}^{HHK}`$ as defined in by means of Schouten’s identity. In the quark-initiated case they read: $`H_6^q`$ $`=`$ $`{\displaystyle \frac{1}{4x}}\left[\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)H_9^{HHK}+x\left(z{\displaystyle \frac{\kappa ^2}{z}}\right)^2H_{10}^{HHK}\right],`$ (89) $`H_7^q`$ $`=`$ $`{\displaystyle \frac{1}{4x}}\left[{\displaystyle \frac{1}{x}}H_9^{HHK}+\left(z+{\displaystyle \frac{\kappa ^2}{z}}\right)H_{10}^{HHK}\right].`$ (90) In the gluon-initiated case the relations between the angular coefficients and the structure functions $`H_i^g`$ are identical to those in the quark-initiated case. The $`\kappa `$-factor in (80)-(89) can be expressed as $`\kappa =\sqrt{{\displaystyle \frac{1x}{x}}z(1z)}.`$ Note that all the above results have been obtained for vector and axial vector currents with unit strength, e.g. no numerical factors in the vertices are taken into account. For the charged current case one has the factor $`{\displaystyle \frac{g_w^4}{16}}{\displaystyle \frac{1}{M_w^4}}U_{ij}^2`$where $`U_{ij}`$ denotes the Kobayashi-Maskawa matrix element and $`g_w=e/\mathrm{sin}\theta _w`$. It is clear that the $`P`$-odd structure functions $`F_8`$ and $`F_9`$ vanish for purely electromagnetic interaction. In the dispersive structure functions $`F_L,F_T,F_3,F_4,F_6`$ and $`F_7`$ have been calculated from the $`𝒪(\alpha _s)`$ Born term contributions. The absorptive $`𝒪(\alpha _s^2)`$ structure functions $`F_5,F_8`$ and $`F_9`$ have been calculated in directly in the DIS channel by determining the appropriate cut contributions of the one-loop diagrams in the DIS channel. The results on $`F_5,F_8`$ and $`F_9`$ in agree with our results derived from crossing. In a full $`𝒪(\alpha _s^2)`$ calculation of lepton-hadron correlations in (2+1) jet production in DIS one would also need the one-loop contributions to the set of dispersive structure functions . As remarked on earlier these can easily be obtained from the $`e^+e^{}`$ one-loop expressions listed in and in this paper through crossing. For the sake of completeness we write down the result for the hard scattering cross section in the case of the $`Z`$-boson exchange: $$\frac{k_0^{}p_0^{}d\widehat{\sigma }}{d^3k^{}d^3p^{}}=\frac{𝒫}{2s}\frac{2}{y^2}\{𝒜+\mathrm{cos}\varphi +𝒞\mathrm{cos}2\varphi +𝒟\mathrm{sin}\varphi +\mathrm{sin}2\varphi \}$$ (91) with $`𝒜`$ $`=`$ $`2(1y)_L+[1+(1y)^2]_T\pm y(2y)_3,`$ (92) $``$ $`=`$ $`\sqrt{1y}[(2y)_4\pm y_5],`$ (93) $`𝒞`$ $`=`$ $`(1y)_6,`$ (94) $`𝒟`$ $`=`$ $`\sqrt{1y}[\pm y_7+(2y)_8],`$ (95) $``$ $`=`$ $`(1y)_9`$ (96) and $`𝒫`$ and $`s`$ are defined by $`𝒫={\displaystyle \frac{1}{2^7\pi ^5}}{\displaystyle \frac{z}{q^4}}\delta (\kappa ^2{\displaystyle \frac{1x}{x}}z(1z)),s=(k+p)^2.`$ For the parity conserving angular coefficients $`_i^{PC},i=L,T,4,6,7`$, which originate from the vector-vector and axial-axial couplings, one has $$_i^{PC}=F_i\left\{e^4e_l^2e_q^2+e^2e_le_qg_z^2\frac{C_V^l\pm C_A^l}{2}C_V^q\text{Re}D+g_z^4(C_V^l\pm C_A^l)^2\frac{(C_V^q)^2+(C_A^q)^2}{16}D^2\right\}.$$ (97) For the parity violating angular coefficients $`_i^{PV},i=3,5,8,9`$, which originate from the vector-axial vector interference contributions, we obtain $$_i^{PV}=F_i\left\{e^2e_le_qg_z^2\frac{C_V^l\pm C_A^l}{2}C_A^q\text{Re}D+g_z^4\frac{(C_V^l\pm C_A^l)^2}{8}C_V^qC_A^qD^2\right\},$$ (98) where $`e_l,e_q`$ are lepton and quark charges, respectively. The upper (lower) sign is for left- (right-) handed initial leptons, and $`D`$ $`=`$ $`q^2(q^2m_Z^2+im_Z\mathrm{\Gamma }_Z)^1,`$ (99) $`g_z`$ $`=`$ $`{\displaystyle \frac{e}{\mathrm{sin}\theta _w\mathrm{cos}\theta _w}},`$ (100) $`C_V^l`$ $`=`$ $`{\displaystyle \frac{1}{2}}+2\mathrm{sin}^2\theta _w,C_A^l={\displaystyle \frac{1}{2}}\text{for leptons with charge -1},`$ (101) $`C_V^q`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{4}{3}}\mathrm{sin}^2\theta _w,C_A^q={\displaystyle \frac{1}{2}}\text{for up quarks},`$ (102) $`C_V^q`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{2}{3}}\mathrm{sin}^2\theta _w,C_A^q={\displaystyle \frac{1}{2}}\text{for down quarks}.`$ (103) The structure functions $`F_i`$ have to be taken from (80) for the quark- and gluon-initiated cases. We stress that in (91) terms proportional to Im$`D`$ have been dropped as they are of a higher order in the electroweak coupling. ## IV Crossing Results for Hadronic Structure Functions in DY In this section we present our results of crossing from the $`e^+e^{}`$-annihilation channel to the DY process. The changes in momenta resulting from crossing are illustrated on the right-hand side of Fig. 2. The calculation proceeds in analogy to that in DIS by following the rules defined in Sec. II. The decomposition of the $`T`$-odd parts of the hadron tensor $`H_{\mu \nu }`$ for the annihilation subprocess and both Compton scattering subprocesses is given by (in this section we concentrate on the absorptive contributions) $`H_{\mu \nu }`$ $`=`$ $`H_5q^2(\widehat{p}_{a\mu }\widehat{p}_{b\nu }\widehat{p}_{a\nu }\widehat{p}_{b\mu })`$ (106) $`+H_8q^4(\widehat{p}_{a\mu }F_\nu +\widehat{p}_{a\nu }F_\mu )`$ $`+H_9q^4(\widehat{p}_{b\mu }F_\nu +\widehat{p}_{b\nu }F_\mu ),`$ with $`\widehat{p}_{a\mu }`$ $`=`$ $`p_{a\mu }{\displaystyle \frac{p_aq}{q^2}}q_\mu ,`$ (107) $`\widehat{p}_{b\mu }`$ $`=`$ $`p_{b\mu }{\displaystyle \frac{p_bq}{q^2}}q_\mu ,`$ (108) $`F_\mu `$ $`=`$ $`iϵ_{\mu \alpha \beta \gamma }p_a^\alpha p_b^\beta q^\gamma .`$ (109) The results for our spin and color averaged T-odd structure functions $`H_5,H_8,H_9`$ are given below. For the annihilation subprocess we obtain: $`i\pi H_5^a`$ $`=`$ $`g^4C_F{\displaystyle \frac{(1+c)(x_a^2x_b^2)}{4(1x_a)(1x_b)}}`$ (112) $`+g^4{\displaystyle \frac{C_F}{N_C}}(C_F{\displaystyle \frac{N_C}{2}})({\displaystyle \frac{cx_a}{(1x_a)^2(1x_b)}}\mathrm{ln}{\displaystyle \frac{c}{x_a}}+{\displaystyle \frac{x_a((x_a+x_b)(1x_ax_b)+x_a^2+x_b^2)}{2(1x_a)(1x_b)}})`$ $`\{x_ax_b\},`$ $`i\pi H_8^a`$ $`=`$ $`g^4C_F{\displaystyle \frac{x_a(x_a^23x_b^2)}{4(1x_a)(1x_b)}}`$ (115) $`+g^4{\displaystyle \frac{C_F}{N_C}}(C_F{\displaystyle \frac{N_C}{2}})x_a({\displaystyle \frac{x_a(x_a2c)}{(1x_a)^3(1x_b)}}\mathrm{ln}{\displaystyle \frac{c}{x_a}}{\displaystyle \frac{x_b^2}{(1x_a)(1x_b)^3}}\mathrm{ln}{\displaystyle \frac{c}{x_b}}`$ $`+{\displaystyle \frac{c(x_a^23x_b^2+2x_a+2x_b)3x_a^2+x_b^2}{2(1x_a)^2(1x_b)^2}}),`$ $`i\pi H_9^a`$ $`=`$ $`i\pi H_8^a(x_ax_b),`$ (116) where $`c=x_a+x_bx_ax_b`$. For the quark-initiated Compton scattering subprocess we find: $`i\pi H_5^{C_q}`$ $`=`$ $`{\displaystyle \frac{g^4}{2}}x_a^2{\displaystyle \frac{2x_ax_bc(1+x_a)}{4c(1x_a)}}`$ (119) $`+{\displaystyle \frac{g^4}{2N_C}}(C_F{\displaystyle \frac{N_C}{2}})x_a(x_a^2(1x_b)[{\displaystyle \frac{x_b}{c^2(1x_a)}}\mathrm{ln}(1c)+{\displaystyle \frac{\mathrm{ln}(c/x_a)}{c(1x_a)^2}}]`$ $`+{\displaystyle \frac{cx_a(1+x_a)2x_ax_b(1+x_b)4x_b^3(1x_a)}{2c(1x_a)}}),`$ $`i\pi H_8^{C_q}`$ $`=`$ $`{\displaystyle \frac{g^4}{2}}x_a^3{\displaystyle \frac{x_a2x_b(1+x_a)}{4c(1x_a)}}`$ (122) $`+{\displaystyle \frac{g^4}{2N_C}}(C_F{\displaystyle \frac{N_C}{2}})x_a^3(x_a[{\displaystyle \frac{x_b^2}{c^3(1x_a)}}\mathrm{ln}(1c)+{\displaystyle \frac{\mathrm{ln}(x_a/c)}{c(1x_a)^3}}]`$ $`+{\displaystyle \frac{cx_a(x_a+4x_b2x_ax_b+1)2(x_a^2+x_ax_b^2x_b^2)}{2c^2(1x_a)^2}}),`$ $`i\pi H_9^{C_q}`$ $`=`$ $`{\displaystyle \frac{g^4}{2}}{\displaystyle \frac{3x_a^3x_b}{4c(1x_a)}}+{\displaystyle \frac{g^4}{2N_C}}(C_F{\displaystyle \frac{N_C}{2}})x_a^2x_b(x_a[{\displaystyle \frac{x_b(2cx_b)}{c^3(1x_a)}}\mathrm{ln}(1c)+{\displaystyle \frac{\mathrm{ln}(x_a/c)}{c(1x_a)^3}}]`$ (124) $`+{\displaystyle \frac{cx_a((12x_b)^23x_a)+4cx_b(1x_b)+2x_a(x_a+x_b)}{2c^2(1x_a)^2}}).`$ Note that we are using the same two colour structures as those in our $`e^+e^{}`$-annihilation expressions. For the antiquark-initiated Compton subprocess one has $$H_5^{C_{\overline{q}}}=H_5^{C_q},H_8^{C_{\overline{q}}}=H_8^{C_q},H_9^{C_{\overline{q}}}=H_9^{C_q}.$$ (125) We consider only DY processes that proceed through $`W`$-exchange. This is sufficient to compare our results with the results in Ref. . Following , we consider the hard scattering cross section $`d\widehat{\sigma }/dQ_T^2d\mathrm{cos}\widehat{\theta }d\mathrm{cos}\theta d\varphi `$. For the $`T`$-odd part of the cross section we use the expansion $$\frac{d\widehat{\sigma }^{Todd}}{dQ_T^2d\mathrm{cos}\widehat{\theta }d\mathrm{cos}\theta d\varphi }=\mathrm{sin}\theta \mathrm{sin}\varphi F_7+\mathrm{sin}2\theta \mathrm{sin}\varphi F_8+\mathrm{sin}^2\theta \mathrm{sin}2\varphi F_9.$$ (126) The definition of angles is the same as in . The $`W`$+jet production is decribed by the transverse momentum $`\stackrel{}{Q}_T`$ of the jet and the scattering angle $`\widehat{\theta }`$ in the $`W`$+jet center of mass frame, $`\theta `$ and $`\varphi `$ are the polar and the azimuthal angle of the lepton emerging from the decay $`Wl\nu `$ in the Collins-Soper frame as shown in Fig. 4. A simple exercise in particle and parton kinematics shows that the variables $`\theta ,\varphi `$ and $`\stackrel{}{Q}_T`$ are identical in the hadron and parton processes. The angular coefficients can be projected from the hadron tensor by means of the Collins-Soper frame projectors $`P_7^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{i}{\sqrt{2}q^2}}{\displaystyle \frac{x_ax_b}{\sqrt{c(1c)}}}(p_a^\mu p_b^\nu p_a^\nu p_b^\mu ),`$ (127) $`P_8^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}q^4}}{\displaystyle \frac{x_a^2x_b^2}{c\sqrt{1c}}}[{\displaystyle \frac{1}{x_b}}(p_a^\mu ϵ^{\nu \alpha \beta \gamma }p_{a\alpha }p_{b\beta }q_\gamma +\mu \nu ){\displaystyle \frac{1}{x_a}}(p_b^\mu ϵ^{\nu \alpha \beta \gamma }p_{a\alpha }p_{b\beta }q_\gamma +\mu \nu )],`$ (128) $`P_9^{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{\sqrt{c}}{q^4}}{\displaystyle \frac{x_a^2x_b^2}{c(1c)}}[{\displaystyle \frac{1}{x_b}}(p_a^\mu ϵ^{\nu \alpha \beta \gamma }p_{a\alpha }p_{b\beta }q_\gamma +\mu \nu )+{\displaystyle \frac{1}{x_a}}(p_b^\mu ϵ^{\nu \alpha \beta \gamma }p_{a\alpha }p_{b\beta }q_\gamma +\mu \nu )],`$ (129) which we use to extract the required components $`F_i`$ from the parton level tensor $`H_{\mu \nu }`$. Then, taking into account the necessary numerical factors, we have the following relations: $`K^1f_7={\displaystyle \frac{3}{\sqrt{2}}}P_7^{\mu \nu }H_{\mu \nu }`$ $`=`$ $`i\pi {\displaystyle \frac{3}{4}}{\displaystyle \frac{\sqrt{c(1c)}}{x_ax_b}}H_5,`$ (130) $`K^1f_8={\displaystyle \frac{3}{2\sqrt{2}}}P_8^{\mu \nu }H_{\mu \nu }`$ $`=`$ $`i\pi {\displaystyle \frac{3}{16}}{\displaystyle \frac{c\sqrt{1c}}{x_a^2x_b^2}}(x_bH_8x_aH_9),`$ (131) $`K^1f_9={\displaystyle \frac{3}{4}}P_9^{\mu \nu }H_{\mu \nu }`$ $`=`$ $`i\pi {\displaystyle \frac{3}{16}}{\displaystyle \frac{\sqrt{c}(1c)}{x_a^2x_b^2}}(x_bH_8+x_aH_9),`$ (132) where $`K=8`$ for the quark- and antiquark-initiated Compton subprocesses and $`K=3`$ for the annihilation subprocess. The functions $`f_7,f_8,f_9`$ are related to the functions $`F_7`$, $`F_8`$ and $`F_9`$ appearing in eq. (126) and are defined in eqs. (6) and (7) of . Substituting our expressions from (112) and (119) to (130) we find complete agreement with the corresponding results of Eqs. (8) and (9) of ref. , except for the sign before the logarithmic term of $`f_{A9}`$. This typo graphical error was corrected in the footnote of on p. 179. ## V Summary and Outlook Starting from the known one-loop results for the quark-quark-gluon gauge boson four-point function in the $`e^+e^{}`$-channel we have used analiticity and crossing to derive the absorptive parts of the same four-point function in the DIS channel and in the DY process. Whereas the imaginary parts of the one-loop four-point function generate a nonmeasurable phase in $`e^+e^{}`$-annihilation one obtains measurable phase effects in DIS and in the DY process leading to the nonvanishing of $`T`$-odd observables which we have derived. We have compared our results with the results of previous calculations where the absorptive parts in DIS and in the DY process were calculated directly in the respective channels. In this paper we have mainly considered $`T`$-odd observables built from triple products of the three-momenta of the respective processes. One can also consider $`T`$-odd observables built from triple products involving also spins. We have discussed such a triple product observable involving the spin of the initial lepton in DIS. As shown in Sec. III the relevant $`T`$-odd observable is fed from the parity conserving $`T`$-odd structure function $`H_5`$ given in this paper. The three absorptive invariant structure functions $`H_5`$, $`H_8`$ and $`H_9`$ presented in this paper suffice to calculate any $`𝒪(\alpha _s^2)`$ contribution to $`T`$-odd observables as long as the spins of partons or hadrons are summed over. If the $`T`$-odd triple product involves the spins of hadrons participating in the process the relevant new spin-dependent $`𝒪(\alpha _s^2)`$ contributions can be easily calculated from the absorptive parts of the $`𝒪(\alpha _s^{3/2})`$ one-loop amplitudes given in this paper. When one folds these with the respective Born term expressions one has to keep the relevant spins unsummed. We hope to return to the subject of spin-dependent $`T`$-odd contributions to DIS and the DY process in the future. ###### Acknowledgements. We would like to thank K. Hagiwara and K. Hikasa for providing us with details of their calculations. We also thank T. Brodkorb, L. Brücher and J. Franzkowski for participating in the early stages of this work. Z.M. thanks A. Davydychev and H.S. Do for discussions. B.M. and Z.M. would like to thank the Institut für Physik, Universität Mainz for hospitality and the BMBF, Germany, under contract 06MZ865, for support. This work was partially supported by the Ministry of Science and Technology of the Republic of Croatia under Contract No. 00980102. ## A The Triangle Anomaly in DIS In this Appendix we present the results for the axial-vector part of the quark loop-induced $`Zgg`$ coupling contribution to DIS. To the best of our knowledge these results have not been presented before. The relevant diagrams and assignments of momenta are shown in Fig. 5. We use the results for the $`Zgg`$ vertex of and apply the above crossing procedure to obtain new results. Following , with momentum assignments $`Z(q)g(k_1)+g(k_2)`$ we have the general decomposition for the transition amplitude: $`T_{\mu \alpha \beta }=f_1(k_2^2ϵ_{\mu \alpha \beta \rho }k_1^\rho +k_{2\alpha }ϵ_{\mu \beta \rho \sigma }k_{2\rho }k_{1\sigma })+f_2(k_1^2ϵ_{\mu \alpha \beta \rho }k_2^\rho +k_{1\beta }ϵ_{\mu \alpha \rho \sigma }k_{2\rho }k_{1\sigma })+`$ (A1) $`f_3(k_1+k_2)_\mu ϵ_{\alpha \beta \rho \sigma }k_{2\rho }k_{1\sigma }+f_4(k_1k_2)_\mu ϵ_{\alpha \beta \rho \sigma }k_{2\rho }k_{1\sigma }),`$ (A2) where the functions $`f_i`$ depend on the Lorentz scalars $`k_1^2,k_2^2`$ and $`q^2=(k_1+k_2)^2`$. As the gluon with momentum $`k_2`$ is on shell, $`f_1`$ does not contribute. Due to the property $`q_\mu L^{\mu \nu }=0`$ in the massless lepton limit $`f_3`$ does not contribute. The amplitude $`f_4`$ vanishes identically at the one loop level due to charge conjugation invariance. Technically speaking, the two contributions from the clockwise and counter-clockwise quark flows in the quark loop cancel in $`f_4`$. Thus, only the term proportional to $`f_2`$ remains. It has to be folded with the Born term amplitude. As it is well known, gauge invariance for the triangle graph with respect to the $`Z`$-boson momenta is broken. This implies that the squared matrix amplitude can be represented in the form of the decomposition (54) plus additional terms proportional to $`q_\mu `$ or $`q_\nu `$ which, however, do not contribute after folding with the leptonic tensor. Performing the necessary crossing and averaging over initial spins and colors, as defined in Sec. 3, one obtains the anomaly contribution to the DIS structure functions $`H_i^q`$. For the quark-initiated case we have: $`H_1^q`$ $`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{x+z}{z}}q^2\mathrm{Re}(f_2^q),`$ (A3) $`H_2^q`$ $`=`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{x(1+xz)}{z(1z)}}q^2\mathrm{Re}(f_2^q),`$ (A4) $`H_3^q`$ $`=`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{x(12x)}{z(1x)}}q^2\mathrm{Re}(f_2^q),`$ (A5) $`H_4^q`$ $`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{(1x)(2z)z(1z)}{\kappa ^2}}q^2\mathrm{Re}(f_2^q),`$ (A6) $`H_6^q`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{(1z)(1x^2+2xz+(1z)^2)+x(zx)}{x\kappa ^2}}q^2\mathrm{Re}(f_2^q),`$ (A7) $`H_7^q`$ $`=`$ $`{\displaystyle \frac{2}{3}}{\displaystyle \frac{(1z)(2(1x)(12x)+xz)+z^2(1x)}{x\kappa ^2}}q^2\mathrm{Re}(f_2^q).`$ (A8) For the antiquark-initiated case results are the same except for $`H_6^{\overline{q}}=H_6^q`$ and $`H_7^{\overline{q}}=H_7^q`$. Note that there are no absorptive parts in the (anti)quark-initiated case because $`f_2^q`$ does not have an imaginary part as we shell see later on. For the gluon-initiated case we similarly obtain: $`H_1^g`$ $`=`$ $`{\displaystyle \frac{12x}{2(1x)}}q^2\mathrm{Re}(f_2^g),`$ (A9) $`H_2^g`$ $`=`$ $`{\displaystyle \frac{z(1+xz)}{\kappa ^2}}q^2\mathrm{Re}(f_2^g),`$ (A10) $`H_3^g`$ $`=`$ $`{\displaystyle \frac{1}{\kappa ^2}}q^2\mathrm{Re}(f_2^g),`$ (A11) $`H_4^g`$ $`=`$ $`{\displaystyle \frac{1x+2xz}{2\kappa ^2}}q^2\mathrm{Re}(f_2^g),`$ (A12) $`H_5^g`$ $`=`$ $`i{\displaystyle \frac{x(12z)}{2z(1z)}}q^2\mathrm{Im}(f_2^g),`$ (A13) $`H_6^g`$ $`=`$ $`{\displaystyle \frac{(1z)(2x(1z)^2(zx)^2)xz(1x)}{4x\kappa ^2}}q^2\mathrm{Re}(f_2^g),`$ (A14) $`H_7^g`$ $`=`$ $`{\displaystyle \frac{1x2z(1z)(1+2x)}{4x\kappa ^2}}q^2\mathrm{Re}(f_2^g),`$ (A15) $`H_8^g`$ $`=`$ $`i{\displaystyle \frac{xz}{\kappa ^2}}q^2\mathrm{Im}(f_2^g),`$ (A16) $`H_9^g`$ $`=`$ $`i{\displaystyle \frac{x(12z)}{\kappa ^2}}q^2\mathrm{Im}(f_2^g).`$ (A17) The functions $`f_2^q`$ and $`f_2^g`$ can be easily written down using eqs. (2.5)-(2.9) of for the DIS region where $`q^2<0`$. For the (anti)quark-initiated case one also has $`k_1^2<0`$. Then we arrive at the following expression for $`f_2^q`$: $`f(w,r)f_2^q={\displaystyle \frac{1}{2\pi ^2q^2}}{\displaystyle \frac{x^2}{(xz)^2}}[{\displaystyle \frac{1}{r}}\{\mathrm{ln}^2(\sqrt{1rw}+\sqrt{rw})\mathrm{ln}^2(\sqrt{1r}+\sqrt{r})\}`$ (A18) $`2\{\sqrt{{\displaystyle \frac{1rw}{rw}}}\mathrm{ln}(\sqrt{1rw}+\sqrt{rw})\sqrt{{\displaystyle \frac{1r}{r}}}\mathrm{ln}(\sqrt{1r}+\sqrt{r})\}+\mathrm{ln}w],`$ (A19) with $`w=k_1^2/q^2`$ and $`r=q^2/4m^2`$ and where all quark masses are set to zero except for the top quark mass denoted by $`m`$. This approximation is valid at high enough energies when only the mass of the top quark is important. Then the contributions from the first two generations of ”light” quarks cancel out between the up($`u,c`$)- and down($`d,s`$)-quarks and only $`b`$\- and $`t`$-quarks contribute to the above functions. One can see that $`f_2^q`$ is a real function, and therefore the $`T`$-odd structure functions in the (anti)quark-initiated case of DIS do not receive any anomaly contribution. However, for the gluon initiated case we have $`k_1^2>0`$, and, clearly, there will be nonvanishing imaginary contribution from the triangle diagram. We have to separately consider the two regions below and above top threshold. Below top threshold with $`q^2<0,\mathrm{\hspace{0.17em}\hspace{0.17em}0}<k_1^2<4m^2(0<rw<1)`$ we get: $`\mathrm{Re}f_2^g`$ $`=`$ $`{\displaystyle \frac{x^2}{2\pi ^2q^2}}[{\displaystyle \frac{1}{r}}\{(\mathrm{sin}^1\sqrt{rw})^2+\mathrm{ln}^2(\sqrt{1r}+\sqrt{r})\}+`$ (A21) $`2\{\sqrt{{\displaystyle \frac{1rw}{rw}}}\mathrm{sin}^1\sqrt{rw}\sqrt{{\displaystyle \frac{1r}{r}}}\mathrm{ln}(\sqrt{1r}+\sqrt{r})\}\mathrm{ln}w],`$ $`\mathrm{Im}f_2^g`$ $`=`$ $`\pi `$ (A22) Since one is below top quark threshold there is only an imaginary part coming from the $`b`$-quark. The contributions of the ($`u,d`$) and ($`c,s`$) quarks cancel pairwise in the real and in the imaginary parts. The $`b`$-quark contribution in the real part is proportional to the last term in (A21). Above top threshold with $`q^2<0,k_1^24m^2(rw1)`$ we have a $`k_1^2`$-dependent imaginary part: $`\mathrm{Re}f_2^g`$ $`=`$ $`{\displaystyle \frac{x^2}{2\pi ^2q^2}}[{\displaystyle \frac{1}{r}}\{{\displaystyle \frac{\pi ^2}{4}}\mathrm{ln}^2(\sqrt{rw}+\sqrt{rw1})+\mathrm{ln}^2(\sqrt{1r}+\sqrt{r})\}+`$ (A24) $`2\{\sqrt{{\displaystyle \frac{rw1}{rw}}}\mathrm{ln}(\sqrt{rw}+\sqrt{rw1})\sqrt{{\displaystyle \frac{1r}{r}}}\mathrm{ln}(\sqrt{1r}+\sqrt{r})\}\mathrm{ln}w],`$ $`\mathrm{Im}f_2^g`$ $`=`$ $`\pi \left[1+{\displaystyle \frac{1}{r}}\mathrm{ln}(\sqrt{rw}+\sqrt{rw1})\sqrt{{\displaystyle \frac{rw1}{rw}}}\right].`$ (A25) The trivial imaginary term in (A22) comes entirely from the ”light” $`b`$-quark contribution. It is also present in (A25), which now also has an imaginary part due to the $`t`$-quark because now one is above top threshold. In the limit of zero $`t`$-quark mass the contribution from the triangle graph should vanish. In this case (A21) and (A22) would be obviously absent and, as one can easily see, the real and imaginary contributions in (A24) and (A25) vanish in that limit, serving as a partial check of the correctness of the above expressions.
warning/0002/math0002073.html	ar5iv	text	# Noncommutative Pieri operators on posets ## 1. Introduction The algebra $`𝒬\text{sym}`$ of quasi-symmetric functions was introduced by Gessel as a source of generating functions for $`P`$-partitions . Since then, quasi-symmetric functions have played an important role as generating functions in combinatorics . The relation of $`𝒬\text{sym}`$ to the more familiar algebra of symmetric functions was clarified by Gelfand et. al. who defined the graded Hopf algebra $`NC`$ of noncommutative symmetric functions and identified $`𝒬\text{sym}`$ as its Hopf dual. Joni and Rota made the fundamental observation that many discrete structures give rise to natural Hopf algebras whose coproducts encode the disassembly of these structures (see also ). A seminal link between these theories was shown by Ehrenborg , whose flag $`f`$-vector quasi-symmetric function of a graded poset gave a Hopf morphism from a Hopf algebra of graded posets to $`𝒬\text{sym}`$. This theory was augmented in where it was shown that the quasi-symmetric function associated to an edge-labelled poset similarly gives a Hopf morphism. That quasi-symmetric function generalised a quasi-symmetric function encoding the structure of the cohomology of a flag manifold as a module over the ring of symmetric functions . We extend and unify these results by means of a simple construction. Given a graded representation of $`NC`$ on the $``$-linear span $`P`$ of a graded poset $`P`$, the matrix coefficients of such an action are linear maps on $`NC`$ and hence quasi-symmetric functions. In Section 2 we show how this situation gives rise to a Hopf morphism as before. In Section 3, we extend this construction to an arbitrary oriented multigraph $`G`$. Sections 4, 6, and 7 give examples of this construction, including rank selection in posets, flag $`f`$-vectors of polytopes, $`P`$-partitions, Stanley symmetric functions, and the multiplication of Schubert classes in the cohomology of flag manifolds. In Section 5, we discuss how properties of the combinatorial structure of $`G`$ may be understood through the resulting quasi-symmetric function. This analysis allows us to relate work of Bayer, Billera, and Liu on Eulerian posets with work of Stembridge on enriched P-partitions. More precisely, we show that the quotient of $`NC`$ by the ideal of the generalised Dehn-Somerville relations is dual to the Hopf subalgebra of peak functions in $`𝒬\text{sym}`$. We also solve the conjecture of , showing that the shifted quasi-symmetric functions form a Hopf algebra. These functions were introduced by Billey and Haiman to define Schubert polynomials for all types. In Section 7, we show how a natural generating function for enumerating peaks in a labelled poset is the quasi-symmetric function for an enriched structure on the poset. Special cases of this combinatorics of peaks include Stembridge’s theory of enriched $`P`$-partitions , the Pieri-type formula for type $`B`$ and $`C`$ Schubert polynomials in , and Stanley symmetric functions of types $`B`$, $`C`$, and $`D`$. These examples linking the diverse areas of Schubert calculus, combinatorics of polytopes and $`P`$-partitions illustrate how this theory transfers techniques and ideas between disparate areas of combinatorics. We thank Sarah Witherspoon who contributed to the appendix on Hopf algebras, and Geanina Tudose for her assistance with fusion coefficients. ## 2. Pieri operators on posets Many interesting families of combinatorial constants can be understood as an enumeration of paths in a ranked partially ordered set (poset) which satisfy certain conditions. One example of this is the Littlewood-Richardson rule in the theory of symmetric functions . This rule describes the multiplicity $`c_{\mu ,\nu }^\lambda `$ of a Schur function $`S_\lambda `$ in the product $`S_\mu S_\nu `$ of two others. The constants $`c_{\mu ,\nu }^\lambda `$ can be seen as an enumeration of all paths in Young’s lattice from $`\mu `$ to $`\lambda `$ satisfying some conditions imposed by $`\nu `$. We note that the constants $`c_{\mu ,\nu }^\lambda `$ are invariant under certain isomorphisms of intervals in Young’s lattice, namely $`c_{\mu ,\nu }^\lambda =c_{\tau ,\nu }^\pi `$ whenever $`\lambda /\mu =\pi /\tau `$. The skew Schur functions $`S_{\lambda /\mu }`$ are generating functions of these constants as we have $$S_{\lambda /\mu }=\underset{\nu }{}c_{\mu ,\nu }^\lambda S_\nu .$$ We generalise the principles of this example, introducing families of algebraic operators to select paths in a given poset. Here, analogues of the Littlewood-Richardson constants count paths in the poset satisfying some conditions imposed by the family of operators. These enumerative combinatorial invariants of the poset are encoded by generating functions which generalise the skew Schur functions. We show that the association of such a generating function to a poset induces a Hopf morphism to $`𝒬\text{sym}`$. Let $`(P,<)`$ be a graded poset with rank function $`rk:P^+`$ and let $`P`$ be the free graded $``$-module generated by the elements of $`P`$. For an integer $`k>0`$, a (right) Pieri operator on $`P`$ is a linear map $`\overline{h}_k:PP`$ which respects the poset structure. By this we mean that for all $`xP`$, the support of $`x.\overline{h}_kP`$ consists only of elements $`yP`$ such that $`x<y`$ and $`rk(y)rk(x)=k`$. We note that such an operator $`\overline{h}_k`$ is of degree $`k`$ on $`P`$. Gelfand et. al. define the Hopf algebra $`NC`$ of noncommutative symmetric functions to be the free associative algebra $`h_1,h_2,\mathrm{}`$ with a generator $`h_k`$ in each positive degree $`k`$ and coproduct $`\mathrm{\Delta }h_k=_{i=0}^kh_ih_{ki}`$, where $`h_0=1`$. It follows that given a family of Pieri operators $`\{\overline{h}_k\}_{k>0}`$ on a poset $`P`$, the map $`h_k\overline{h}_k`$ turns $`P`$ into a graded (right) $`NC`$-module. Conversely, any graded right action of $`NC`$ on $`P`$ which respects the poset structure of $`P`$ gives a family of Pieri operators on $`P`$. When the context is clear, we may identify the generator $`h_k`$ with the operator $`\overline{h}_k`$. Given such a representation of $`NC`$ on $`P`$ and $`x,yP`$, the association of $`\mathrm{\Psi }NC`$ to the coefficient of $`y`$ in $`x.\mathrm{\Psi }`$ is a linear map on $`NC`$. These matrix coefficients are elements of the Hopf dual of $`NC`$ which is the Hopf algebra $`𝒬\text{sym}`$ of quasi-symmetric functions . These coefficients vanish unless $`xy`$. This gives a collection of quasi-symmetric functions $`K_{[x,y]}`$ associated to every interval $`[x,y]`$ of $`P`$. Let $`P`$ be the free $``$-module with basis given by Cartesian products of intervals $`[x,y]`$ of $`P`$, modulo identifying all singleton intervals $`[x,x]`$ with the unit $`1`$ and empty intervals with zero. Then $`P`$ is a graded $``$-algebra whose product is the cartesian product of intervals, and whose grading is induced by the rank of an interval of $`P`$. It has a natural coalgebra structure induced by $$\mathrm{\Delta }A=\underset{xA}{}[\widehat{0}_A,x][x,\widehat{1}_A],$$ where $`A=[\widehat{0}_A,\widehat{1}_A]`$ is an interval of $`P`$ with minimal element $`\widehat{0}_A`$ and maximal element $`\widehat{1}_A`$. Projection onto $``$ of the degree 0 component of $`P`$ is the counit. It follows that $`P`$ is a bialgebra. It is graded, therefore by Proposition A.2, there is a unique antipode and $`P`$ is a Hopf algebra. ###### Theorem 2.1. For any graded poset $`P`$, $`P`$ is a Hopf algebra. Suppose we have a family of Pieri operators on a poset $`P`$. Since $`NC`$ is a Hopf algebra, the action of $`NCNC`$ on $`PP=(P\times P)`$ pulls back along the the coproduct $`\mathrm{\Delta }`$ to give an action of $`NC`$ on $`(P\times P)`$. We iterate this and use coassociativity to get an action of $`NC`$ on the $``$-linear span of $`P^k`$. Since a product of intervals of $`P`$ is an interval in such an iterated product of $`P`$ with itself, we may extend the definition of $`K`$ to the generators of $`P`$ and then by linearity to $`P`$ itself, obtaining a $``$-linear homogeneous map $`K:P𝒬\text{sym}`$. Let $`,`$ be the bilinear form on $`P`$ induced by the Kronecker delta function on the elements of $`P`$. ###### Theorem 2.2. The map $`K:P𝒬\text{sym}`$ is a morphism of Hopf algebras. ###### Proof. We show that $`K`$ respects product and coproduct, which suffices. For $`x,yP`$ and $`\mathrm{\Psi }NC`$, we have $`K_{[x,y]}(\mathrm{\Psi })=x.\mathrm{\Psi },y`$. Thus for $`xP`$ and $`\mathrm{\Psi }NC`$, $$x.\mathrm{\Psi }=\underset{y}{}x.\mathrm{\Psi },yy=\underset{y}{}K_{[x,y]}(\mathrm{\Psi })y.$$ Let $`A=[\widehat{0}_A,\widehat{1}_A]`$ and $`B=[\widehat{0}_B,\widehat{1}_B]`$ be intervals of $`P^{k_1}`$ and $`P^{k_2}`$ respectively. For $`\mathrm{\Psi }NC`$, using Sweedler notation for the coproduct $$\mathrm{\Delta }\mathrm{\Psi }=\mathrm{\Psi }_a\mathrm{\Psi }_b,$$ and the duality between the product of $`𝒬\text{sym}`$ and the coproduct of $`NC`$, we obtain $`K_{A\times B}(\mathrm{\Psi })`$ $`=`$ $`(\widehat{0}_A\widehat{0}_B).\mathrm{\Psi },\widehat{1}_A\widehat{1}_B`$ $`=`$ $`{\displaystyle \widehat{0}_A}.\mathrm{\Psi }_a\widehat{0}_B.\mathrm{\Psi }_b,\widehat{1}_A\widehat{1}_B`$ $`=`$ $`{\displaystyle }\widehat{0}_A.\mathrm{\Psi }_a,\widehat{1}_A\widehat{0}_B.\mathrm{\Psi }_b,\widehat{1}_B`$ $`=`$ $`{\displaystyle K_A(\mathrm{\Psi }_a)K_B(\mathrm{\Psi }_b)}=(K_AK_B)(\mathrm{\Delta }\mathrm{\Psi })=(K_AK_B)(\mathrm{\Psi }).`$ Let $`A=[\widehat{0}_A,\widehat{1}_A]`$ be an interval of $`P^k`$ and $`\mathrm{\Psi },\mathrm{\Phi }NC`$. Using the duality between the coproduct of $`𝒬\text{sym}`$ and the product of $`NC`$, we have $`(\mathrm{\Delta }K_A)(\mathrm{\Psi }\mathrm{\Phi })`$ $`=`$ $`K_A(\mathrm{\Psi }\mathrm{\Phi })=(\widehat{0}_A.\mathrm{\Psi }).\mathrm{\Phi },\widehat{1}_A`$ $`=`$ $`{\displaystyle \underset{y}{}}(K_{[\widehat{0}_A,y]}(\mathrm{\Psi })y).\mathrm{\Phi },\widehat{1}_A`$ $`=`$ $`{\displaystyle \underset{x,y}{}}K_{[\widehat{0}_A,y]}(\mathrm{\Psi })K_{[y,x]}(\mathrm{\Phi })x,\widehat{1}_A`$ $`=`$ $`{\displaystyle \underset{yA}{}}K_{[\widehat{0}_A,y]}(\mathrm{\Psi })K_{[y,\widehat{1}_A]}(\mathrm{\Phi })=K_{\mathrm{\Delta }A}(\mathrm{\Psi }\mathrm{\Phi }).\text{}`$ The map $`K`$ is a generating function for the enumerative combinatorial invariants associated to the $`NC`$-structure of $`P`$. Let $`\{a_\alpha \}`$ be a graded basis of $`NC`$ and let $`\{b_\alpha \}`$ be the corresponding dual basis in $`𝒬\text{sym}`$. Then (2.1) $$K_{[x,y]}=\underset{\alpha }{}x.\overline{a}_\alpha ,yb_\alpha .$$ We interpret the coefficient of $`b_\alpha `$ in $`K_{[x,y]}`$ as the number of paths from $`x`$ to $`y`$ satisfying some condition imposed by $`a_\alpha `$. We reformulate Equation 2.1 in terms of the Cauchy element of Gelfand et. al. . This element relates each graded basis of the Hopf algebra $`NC=_{n0}NC_n`$ to its corresponding dual basis in the Hopf algebra $`𝒬\text{sym}=_{n0}𝒬\text{sym}_n`$. More precisely, let $`\alpha =(\alpha _1,\alpha _2,\mathrm{},\alpha _{\mathrm{}})`$ with $`\mathrm{}0`$ be a sequence of positive integers. Such a sequence is a composition of $`n`$, denoted $`\alpha n`$, if $`n=_{i=1}^{\mathrm{}}\alpha _i`$. By convention the empty sequence for $`\mathrm{}=0`$ is the unique composition of $`0`$. The complete $`NC`$-functions $`\{S^\alpha =h_{\alpha _1}h_{\alpha _2}\mathrm{}h_\alpha _{\mathrm{}}\}_{\alpha n}`$ and the ribbon $`NC`$-functions $`\{R_\alpha \}_{\alpha n}`$ form two bases of $`NC_n`$. Similarly the monomial quasi-symmetric functions $`\{M_\alpha \}_{\alpha n}`$ and the complete quasi-symmetric functions $`\{F_\alpha \}_{\alpha n}`$ form two bases of $`𝒬\text{sym}_n`$. In the graded completion of $`_{n0}NC_n𝒬\text{sym}_n`$ we have the following Cauchy element: (2.2) $$𝒞:=\underset{\alpha }{}a_\alpha b_\alpha =\underset{\alpha }{}R_\alpha F_\alpha =\underset{\alpha }{}S^\alpha M_\alpha ,$$ where $`\{a_\alpha \}`$ and $`\{b_\alpha \}`$ is any pair of dual graded bases. The right action of $`NC`$ on $`P`$ extends linearly to an action of the graded completion of $`_{n0}NC_n𝒬\text{sym}_n`$ on the completion of $`𝒬\text{sym}_{}P`$. The following theorem is simply a reformulation of Equation 2.1. ###### Theorem 2.3. For any family of Pieri operators on a poset $`P`$ and $`x,yP`$, $$K_{[x,y]}=x.𝒞,y.$$ Expanding the quasi-symmetric function $`K_{[x,y]}`$ in a basis of $`𝒬\text{sym}`$ gives a family of enumerative combinatorial invariants for the given action of $`NC`$ on the poset $`P`$. In this way, the functions $`K_{[x,y]}`$ are seen to be analogues of the skew Schur functions presented at the beginning of this section. ## 3. Pieri operators on graphs We extend this simple construction on graded posets to (locally finite) oriented multigraphs. Let $`G=(V,E)`$ be a multigraph where $`V`$ is the set of vertices and $`E`$ is a function $`V\times V^+`$ such that $$\underset{y^{}V}{}E(x,y^{})\text{and}\underset{x^{}V}{}E(x^{},y)$$ are both finite for all $`x,yV`$. The value $`E(x,y)`$ identifies the number of arrows from $`x`$ to $`y`$ in $`G`$. The function $`E`$ is the incidence matrix of the graph $`G`$, and $`E^r`$ is the matrix product of $`r`$ copies of $`E`$. Given $`x,yV`$, let $`[x,y]`$ be the set of all paths from $`x`$ to $`y`$. Consider a graded version of this set, $$[x,y]=\underset{r0}{}[x,y]^{(r)},$$ where the interval $`[x,y]^{(r)}`$ is the set of all paths of length $`r`$ from $`x`$ to $`y`$. Note that $`\|[x,y]^{(r)}\|=E^r(x,y)`$ is finite. Let $`G`$ denote the free $``$-module generated by $`V`$. Here, a Pieri operator is a linear map $`\overline{h}_k:GG`$ where for all $`xV`$, the support of $`x.\overline{h}_kG`$ consists of elements $`yV`$ such that $`E^k(x,y)>0`$. As before, a family $`\{\overline{h}_k\}_{k>0}`$ of Pieri operators induces on $`G`$ the structure of an $`NC`$-module. We thus obtain a collection of linear maps $`K_{[x,y]^{(r)}}:NC`$ given by $`\mathrm{\Psi }x.\mathrm{\Psi }^{(r)},y`$ where $`\mathrm{\Psi }^{(r)}`$ is the $`r`$th-homogeneous component of $`\mathrm{\Psi }`$, and thus quasi-symmetric functions $`K_{[x,y]^{(r)}}𝒬\text{sym}_r`$. We define a Hopf algebra $`G`$ associated to $`G`$. Define the product of two intervals by $`[x,y]^{(r)}\times [u,v]^{(s)}:=[(x,u),(y,v)]^{(r+s)}`$, an interval in $`G\times G`$. Let $`G`$ be the free $``$-module with basis given by products of intervals $`[x,y]^{(r)}`$ in $`G`$, modulo identifying all intervals $`[x,x]^{(0)}`$ with the unit $`1`$ and all empty intervals $`[x,y]^{(r)}`$ with zero. If we let $`r`$ be the degree of an element $`[x,y]^{(r)}`$, then $`G`$ is a graded $``$-algebra with product $`\times `$. The algebra $`G`$ has a natural coalgebra structure induced by $$\mathrm{\Delta }[x,y]^{(r)}=\underset{s=0}{\overset{r}{}}\underset{zV}{}[x,z]^{(s)}[z,y]^{(rs)}.$$ The counit is again the projection onto $``$ of the degree $`0`$ component. Since the bialgebra $`G`$ is graded, we have the following theorem. ###### Theorem 3.1. For any oriented multigraph $`G`$, $`G`$ is a Hopf algebra. Suppose we have a family of Pieri operators on a graph $`G`$. As in Section 2, we have an action of $`NC`$ on the $``$-linear span of $`G^k`$, for any positive integer $`k`$. Since the generators of $`G`$ are sets of the form $`[w,z]^{(t)}`$ in $`G^k`$, we may define the quasi-symmetric function $`K`$ on each of these generators of $`G`$, and then extend by linearity to $`G`$ itself, obtaining a $``$-linear graded map $`K:G𝒬\text{sym}`$. We leave to the reader the straightforward extension of Theorem 2.2. ###### Theorem 3.2. The map $`K:G𝒬\text{sym}`$ is a morphism of Hopf algebras. To extend Theorem 2.3, we decompose the Cauchy element $`𝒞`$ into its homogeneous components: $`𝒞=_{r0}𝒞_r`$. For example, we can use $`𝒞_r=_{\alpha r}S^\alpha M_\alpha `$. The following is immediate. ###### Theorem 3.3. For any family of Pieri operators on a graph $`G`$ and $`x,yG`$, $$K_{[x,y]^{(r)}}=x.𝒞_r,y.$$ ###### Remark 3.4. These constructions generalise those of Section 2. Given a ranked poset $`P`$ we associate to it the incidence graph $`G_P=(P,E)`$ where $`E(x,y)=1`$ if $`y`$ covers $`x`$ and $`E(x,y)=0`$ otherwise. The intervals $`[x,y]^{(r)}`$ are empty unless $`r=rk(y)rk(x)`$ in which case $`[x,y]^{(r)}`$ is equal to the saturated chains in $`[x,y]`$. In such a case we omit the superscript $`(r)`$ and arrive again at the results of Section 2. ## 4. Three simple examples We give three simple examples to illustrate our theory. ###### Example 4.1. Simple path enumeration. Given a graph $`G`$, define the Pieri operator $`\overline{h}_k:GG`$ by $$x.\overline{h}_k=\underset{yG}{}E^k(x,y)y.$$ This action of $`NC`$ satisfies $`\overline{h}_a\overline{h}_b=\overline{h}_{a+b}`$ for all $`a,b^+`$. From this we deduce that $`K_{[x,y]^{(r)}}=E^r(x,y)_{\alpha r}M_\alpha `$. Thus $`K`$ simply enumerates all paths of length $`r`$ from $`x`$ to $`y`$. When $`G`$ is a graded poset $`P`$, $`E^r(x,y)=0`$ unless $`r=rk(y)rk(x)`$ and $`xy`$ in $`P`$. In this case, $`E^r(x,y)`$ counts the saturated chains in the interval $`[x,y]`$. ###### Example 4.2. Skew Schur functions. Let $`(P,<)`$ be Young’s lattice of partitions. For $`\mu P`$, define $`\mu .\overline{h}_k`$ to be the sum of all partitions $`\lambda `$ such that $`\lambda /\mu `$ is a horizontal strip and $`\|\lambda \|\|\mu \|=k`$. This family of Pieri operators lifts the action of the algebra $`\mathrm{\Lambda }`$ of symmetric functions on itself. It follows that $`K_{[\mu ,\lambda ]}`$ is the skew Schur function $`S_{\lambda /\mu }`$. ###### Example 4.3. Rank selection Pieri operators and flag $`f`$-vectors. Given any ranked poset $`P`$, consider the Pieri operator obtained by setting $`x.\overline{h}_k`$ equal to the sum of all $`y>x`$ such that $`rk(y)rk(x)=k`$. In this case, $`x.\overline{S}^\alpha ,y`$ counts all chains in the rank-selected poset obtained from $`[x,y]`$ with ranks given by $`\alpha `$, and $`K`$ is Ehrenborg’s flag $`f`$-vector quasi-symmetric generating function . ## 5. Structure from Hopf subalgebras Suppose, for a family of posets, we have a class of enumerative combinatorial invariants which possesses some additional structure. In many situations, the associated families of Pieri operators satisfy some relations, and the resulting actions of $`NC`$ are carried by a Hopf quotient of $`NC`$. Equivalently, the images of $`K`$ lie in the dual of this quotient, a Hopf subalgebra of $`𝒬\text{sym}`$. More precisely, let an action of $`NC`$ on $`G`$ be given by a homomorphism $`\varphi `$ from $`NC`$ to the linear endomorphism ring End$`(G)`$. Let $``$ be an ideal generated by some relations satisfied by the Pieri operators. When $``$ is a Hopf ideal, so that we have $`\mathrm{\Delta }()NC+NC`$, we have the commuting diagram $``$ as the functions $`K_{[w,z]^{(t)}}`$ are characters of representations $`\varphi ^k`$ on $`G^k`$. In particular, Equation 2.1 has a more specialised form. Given a basis $`\{c_\lambda \}`$ of $`NC/`$ and $`\{d_\lambda \}`$ its dual basis inside $`𝒬\text{sym}`$, we have (5.1) $$K_{[x,y]^{(r)}}=\underset{\lambda }{}x.c_\lambda ,yd_\lambda $$ where the sum is only over the index set of the given basis for $`NC/`$. Here the numbers $`x.c_\lambda ,y`$ are special cases of the enumerative invariants in Equation 2.1. We illustrate these principles in a series of examples which introduce certain classes of Pieri operators defined by quotients of $`NC`$. ###### Example 5.1. Simple path enumeration. In Example 4.1, the ideal $``$ is generated by $`h_{a+b}h_ah_b`$ for all $`a,b>0`$. This is not a Hopf ideal since $`\mathrm{\Delta }(h_2h_1h_1)`$ $`=`$ $`h_21+h_1h_1+1h_2(h_11+1h_1)^2`$ $`=`$ $`(h_2h_1h_1)1+1(h_2h_1h_1)h_1h_1,`$ which is not contained in $`NC+NC`$. ###### Example 5.2. Symmetric Pieri operators. A family of Pieri operators is symmetric if $`\overline{h}_a\overline{h}_b=\overline{h}_b\overline{h}_a`$ for all $`a,b>0`$. In this case, $`NC/[h_1,h_2,\mathrm{}]`$, which is the self-dual Hopf algebra $`\mathrm{\Lambda }`$ of symmetric functions (see ), and thus $``$ is a Hopf ideal. Symmetric Pieri operators satisfy $`x.S^\alpha =x.S^\beta `$ whenever $`\alpha `$ and $`\beta `$ determine the same partition, and hence by Equation 2.1 we can write $`K_{[x,y]}`$ in the form $$\underset{\lambda r}{}A^\lambda \underset{\lambda (\alpha )=\lambda }{}M_\alpha $$ where $`r`$ is the rank of the interval $`[x,y]`$, $`\lambda (\alpha )`$ is the partition determined by $`\alpha `$, and $`A^\lambda `$ is some constant. By definition, $`_{\lambda (\alpha )=\lambda }M_\alpha `$ is the symmetric function $`m_\lambda `$, and so we see again that the image of $`K`$ lies in $`\mathrm{\Lambda }`$. Symmetric Pieri operators can be found in Example 4.2 and in Sections 6 and 7. It is interesting to use other known dual bases of $`\mathrm{\Lambda }`$ in Equation 5.1, in particular, its self-dual basis $`\{S_\lambda \}`$ of Schur functions. ###### Example 5.3. Flag $`f`$-vectors of Eulerian posets. Consider Example 4.3 when the given ranked poset $`P`$ is Eulerian. The flag $`f`$-vectors of Eulerian posets satisfy the linear generalised Dehn-Sommerville or Bayer-Billera relations . Billera and Liu \[9, Proposition 3.3\] show that the ideal of relations satisfied by such Pieri operators is generated by the (even) Euler relations (5.2) $$\underset{i+j=2n}{}(1)^i\overline{h}_i\overline{h}_j=2\overline{h}_{2n}+\underset{i=1}{\overset{2n1}{}}(1)^i\overline{h}_i\overline{h}_{2ni}=0,$$ where $`n`$ is a positive integer. As in , let $``$ be the ideal of $`NC`$ generated by $$X_{2n}:=\underset{i+j=2n}{}(1)^ih_ih_j=2h_{2n}+\underset{i=1}{\overset{2n1}{}}(1)^ih_ih_{2ni}.$$ Then we have the following algebra isomorphism, $$NC/y_1,y_3,y_5,\mathrm{}.$$ where $`y_i`$ has degree $`i`$. We identify the dual $`\left(NC/\right)^{}`$ inside $`𝒬\text{sym}`$ to be the peak Hopf algebra $`\mathrm{\Pi }`$ introduced by Stembridge in his study of enriched P-partitions. This shows that $``$ is a Hopf ideal over $``$. ###### Theorem 5.4. $`(NC/)^{}\mathrm{\Pi }`$. Let us clarify some notation for the proof of Theorem 5.4. Given compositions $`\alpha ,\beta m`$, write $`\beta \alpha `$ if $`\beta `$ is a refinement of $`\alpha `$ and let $`\beta ^{}`$ be the refinement of $`\beta `$ obtained by replacing all components $`\beta _i>1`$ of $`\beta `$ for $`i>1`$ with $`[1,\beta _i1]`$. Given a composition $`\alpha m`$ with $`\alpha _1>1`$ if $`m>1`$, the Billey-Haiman shifted quasi-symmetric functions are shown to have the formula (5.3) $$\theta _\alpha =\underset{\genfrac{}{}{0pt}{}{\beta m}{\beta ^{}\alpha }}{}2^{k(\beta )}M_\beta ,$$ where $`k(\beta )`$ is the number of components of $`\beta `$. If $`\alpha `$ is a composition with all components greater than 1, except perhaps the last, then we call $`\alpha `$ a peak composition and $`\theta _\alpha `$ a peak function. In Stembridge shows that the linear span $`\mathrm{\Pi }`$ of the peak functions is a subalgebra of $`𝒬\text{sym}`$. In fact, $`\mathrm{\Pi }`$ is a Hopf subalgebra of $`𝒬\text{sym}`$ . Recall that in the identification of $`𝒬\text{sym}`$ as the graded linear dual of $`NC`$, the families $`\{M_\alpha \}`$ and $`\{S^\alpha \}`$ are dual bases. That is, $`M_\alpha (S^\beta )=1`$ if $`\alpha =\beta `$ and $`0`$ otherwise. Given any two compositions $`\eta =(\eta _1,\eta _2,\mathrm{})`$ and $`ϵ=(ϵ_1,ϵ_2,\mathrm{})`$, let $`\eta ϵ`$ be the concatenation $`(\eta _1,\eta _2,\mathrm{},ϵ_1,ϵ_2,\mathrm{})`$. ###### Lemma 5.5. The peak algebra $`\mathrm{\Pi }`$ annihilates the ideal $``$. ###### Proof. We show that a peak function $`\theta _\alpha `$ annihilates any function of the form $`S^\beta X_{2n}S^\gamma `$. Since $`M_\eta (S^ϵ)=0`$ unless $`\eta =ϵ`$, it follows that we need only study those summands $`2^{k(\delta )}M_\delta `$ in $`\theta _\alpha `$ such that either $`\delta =\beta 2n\gamma `$ or else $`\delta =\beta i(2ni)\gamma `$. Now if $`2^{k(\beta 2n\gamma )}M_{\beta 2n\gamma }`$ is a summand of $`\theta _\alpha `$, then it follows that all summands of the form $`2^{k(\beta i(2ni)\gamma )}M_{\beta i(2ni)\gamma }`$ will also belong to $`\theta _\alpha `$. By the Euler relations 5.2, it follows immediately that $`S^\beta X_{2n}S^\gamma `$ is annihilated by $`\theta _\alpha `$. We are left to consider the case where $`2^{k(\beta i(2ni)\gamma )}M_{\beta i(2ni)\gamma }`$ is a summand of $`\theta _\alpha `$ but not $`2^{k(\beta 2n\gamma )}M_{\beta 2n\gamma }`$. Observe that from the definition 5.3 we must have $`n>1`$ since if $`2^{k(\beta 11\gamma )}M_{\beta 11\gamma }`$ is a summand of $`\theta _\alpha `$ then $`2^{k(\beta 2\gamma )}M_{\beta 2\gamma }`$ will be too as $`(\beta 11\gamma )^{}=(\beta 2\gamma )^{}`$. Suppose $`\beta m`$ and let $`j`$ be such that $$\alpha _1+\alpha _2+\mathrm{}+\alpha _{j1}m<\alpha _1+\alpha _2+\mathrm{}+\alpha _j.$$ Then $`M_{\beta i(2ni)\gamma }`$ is in the support of $`\theta _\alpha `$ if and only if $`\alpha _1+\alpha _2+\mathrm{}+\alpha _j`$ is $`m+i`$ or $`m+i+1`$. If it is $`m+i`$, then for $`i1`$ $$\theta _\alpha (S^\beta X_{2n}S^\gamma )=\theta _\alpha ((1)^iS^{\beta i(2ni)\gamma }+(1)^{(i1)}S^{\beta i1(2ni+1)\gamma })=0.$$ If $`i=1`$, then $`(\beta 1(2n1)\gamma )^{}\alpha `$. Since $`\theta _\alpha `$ is a peak function and $`n>1`$, we must have $`\alpha _{j+1}>1`$. This implies that $`2n1\alpha _{j+1}`$, hence $`(\beta 2n\gamma )^{}\alpha `$ and $`2^{k(\beta 2n\gamma )}M_{\beta 2n\gamma }`$ is a summand of $`\theta _\alpha `$, which contradicts our assumption. A similar argument for $`m+i+1`$ completes the proof of the lemma. Proof of Theorem 5.4. By Lemma 5.5, $`\mathrm{\Pi }(NC/)^{}\left(y_1,y_3,y_5,\mathrm{}\right)^{}`$. This containment is an equality since the dimension of the $`i`$th homogeneous component of both $`\mathrm{\Pi }`$ and $`y_1,y_3,y_5,\mathrm{}`$ is the $`i`$th Fibonacci number. ###### Definition 5.6. Pieri operators are symmetric if the image of $`K`$ lies within the algebra $`\mathrm{\Lambda }`$ of symmetric functions. Similarly, Pieri operators are Eulerian if the image of $`K`$ lies within $`\mathrm{\Pi }_{}`$, the $``$-span of $`\mathrm{\Pi }`$. This occurs if there is some scalar multiple $`\alpha _k\overline{h}_k`$ of each Pieri operator such that the $`\alpha _k\overline{h}_k`$ satisfy the Euler relations 5.2. We solve the conjecture presented in related to the general functions $`\theta _\alpha `$ introduced by Billey and Haiman . Let $`\mathrm{\Xi }`$ be the $``$-linear span of all the $`\theta _\alpha `$. ###### Theorem 5.7. The space $`\mathrm{\Xi }`$ is a Hopf subalgebra of $`𝒬\text{sym}`$. Moreover the set $$𝒥=\{\mathrm{\Psi }NC\|\theta (\mathrm{\Psi })=0\text{ for all }\theta \mathrm{\Xi }\}$$ is the principal ideal generated by $`X_2=2h_2h_1h_1`$. ###### Proof. We first show that $`𝒥`$ is an ideal and it is included in $`=X_{2n}`$, the ideal generated by the Euler relations. By Theorem 3.2 of , $`\mathrm{\Xi }`$ is a coalgebra. Hence $`\mathrm{\Xi }^{}=NC/𝒥`$ is an algebra, which shows that $`𝒥`$ is an ideal. Since $`\mathrm{\Pi }\mathrm{\Xi }`$ we have that $`𝒥`$. Now it is straightforward to check that $`X_2𝒥`$. Let $`\widehat{𝒥}𝒥`$ be the principal ideal generated by $`X_2`$. Since $`\mathrm{\Delta }(X_2)=1X_2+X_21`$ we have that $`\widehat{𝒥}`$ is a Hopf ideal and $`NC/\widehat{𝒥}`$ is a Hopf algebra. Its dual $`\left(NC/\widehat{𝒥}\right)^{}`$ is a Hopf subalgebra of $`𝒬\text{sym}`$ contained in $`\mathrm{\Xi }`$. To conclude our argument, we show that the dimension of the homogeneous components of degree $`n`$ in $`NC/\widehat{𝒥}`$ and $`\mathrm{\Xi }`$ are equal for all $`n`$. In $`NC/\widehat{𝒥}`$, the homogeneous component of degree $`n`$ has dimension given by the number of compositions of $`n`$ that contain no component equal to $`2`$. This satisfies the recurrence $`\pi _n=\pi _{n1}+\pi _{n2}+\pi _{n4}`$ with initial conditions $`\pi _1=1`$, $`\pi _2=1`$, $`\pi _3=2`$ and $`\pi _4=4`$. This is exactly the recurrence of Theorem 4.3 in given for calculating the dimension of the homogeneous component of degree $`n`$ in $`\mathrm{\Xi }`$. Hence $`\left(NC/\widehat{𝒥}\right)^{}=\mathrm{\Xi }`$ is a Hopf algebra and $`\widehat{𝒥}=𝒥`$. ## 6. Descent Pieri operators ###### Definition 6.1. An (edge)-labelled poset is a graded poset $`P`$ whose covers (edges of its Hasse diagram) are labelled with integers. To enumerate chains according to the descents in their sequence of (edge) labels, we use the descent Pieri operator $$x.\overline{h}_k:=\underset{\omega }{}\mathrm{end}(\omega ),$$ where the sum is over all chains $`\omega `$ of length $`k`$ starting at $`x`$, $$\omega :x\stackrel{b_1}{-\to }x_1\stackrel{b_2}{-\to }\mathrm{}\stackrel{b_k}{-\to }x_k=:\mathrm{end}(\omega ),$$ with no descents, that is $`b_1b_2\mathrm{}b_k`$. The resulting quasi-symmetric function $`K_P`$ was studied in , where (with some effort), it was shown to give a Hopf morphism from a reduced incidence Hopf algebra to $`𝒬\text{sym}`$. We may likewise have edge-labelled graphs, and define descent Pieri operators in that context. To a subset $`\{j_1<j_2<\mathrm{}<j_k\}`$ of $`[n1]`$, we associate the composition $`(j_1,j_2j_1,\mathrm{},nj_k)`$. Given a saturated chain $`\omega `$ in $`P`$ with labels $`b_1,b_2,\mathrm{},b_n`$, let $`D(\omega )`$ be the descent composition of $`\omega `$, that is the composition associated to the descent set $`\{ib_i>b_{i+1}\}`$ of $`\omega `$. Then (Equation 4 of ) we have (6.1) $$K_{[x,y]}=F_{D(\omega )},$$ where the sum is over all saturated chains $`\omega `$ in the interval $`[x,y]`$, and $`F_\alpha `$ is the complete (or fundamental) quasi-symmetric function. If we label a cover $`\mu \lambda `$ in Young’s lattice consistently by either the column or content of the box in $`\lambda /\mu `$, then the descent Pieri operator coincides with the Pieri operator of Example 4.2. ###### Example 6.2. $`k`$-Bruhat order and skew Schubert functions. The Pieri-type formula for the classical flag manifold suggests a symmetric Pieri operator on a suborder of the Bruhat order on the symmetric group, which encodes the structure of the cohomology of the flag manifold as a module over the ring of symmetric polynomials. Let $`𝒮_n`$ denote the symmetric group on $`n`$ elements and let $`\mathrm{}(w)`$ be the length of a permutation $`w`$ in this Coxeter group. We define the $`k`$-Bruhat order $`<_k`$ by its covers. Given permutations $`u,w𝒮_n`$, we say that $`u_kw`$ if $`\mathrm{}(u)+1=\mathrm{}(w)`$ and $`u^1w=(i,j)`$, where $`(i,j)`$ is a reflection with $`ik<j`$. When $`u_kw`$, we write $`wu^1=(a,b)`$ with $`a<b`$ and label the cover $`u_kw`$ in the $`k`$-Bruhat order with the integer $`b`$. The descent Pieri operators on this labelled poset are symmetric as $`\overline{h}_m`$ models the action of the Schur polynomial $`h_m(x_1,\mathrm{},x_k)`$ on the basis of Schubert classes (indexed by $`𝒮_n`$) in the cohomology of the flag manifold $`SL(n,)/B`$. We also have $$K_{[u,w]}=\underset{\lambda }{}c_{u,(\lambda ,k)}^wS_\lambda $$ where $`c_{u,(\lambda ,k)}^w`$ is the coefficient of the Schubert polynomial $`𝔖_w`$ in the product$`𝔖_uS_\lambda (x_1,\mathrm{},x_k)`$. This is the skew Schubert function $`S_{wu^1}`$ of . Geometry shows these coefficients $`c_{u,(\lambda ,k)}^w`$ are non-negative. It is an important open problem to give a combinatorial or algebraic proof of this fact. ###### Example 6.3. The weak order on $`𝒮_n`$ and Stanley symmetric functions. The weak order on the symmetric group $`𝒮_n`$ is the labelled poset whose covers are $`ww(i,i+1)`$, with label $`i`$ if $`\mathrm{}(w)+1=\mathrm{}(w(i,i+1))`$. In , it is shown that the descent Pieri operators on this labelled poset are symmetric and $`K_{[u,w]}`$ is the Stanley symmetric function or stable Schubert polynomial $`F_{wu^1}`$, introduced by Stanley to study reduced decompositions of chains in the weak order on $`𝒮_n`$ . ###### Example 6.4. noncommutative Schur functions of Fomin and Greene. Fomin and Greene have a theory of combinatorial representations of certain noncommutative Schur functions . These are a different noncommutative version of symmetric functions than $`NC`$. Using the Cauchy element in their algebra, they obtain symmetric functions $`F_{y/x}`$ which include Schur functions, Stanley symmetric functions, stable Grothendieck polynomials, and others. A combinatorial representation gives rise to an edge-labelled directed graph so that the functions $`F_{y/x}`$ of Fomin-Greene are the functions $`K_{[x,y]}`$ coming from the descent Pieri operators on this structure. Let $`FG_n`$ be the quotient of the free associative algebra $`u_1,u_2,\mathrm{},u_n`$ by the two-sided ideal generated by the following relations (6.2) $$\begin{array}{ccccc}\hfill u_iu_ku_j& =& u_ku_iu_j,\hfill & ij<k\hfill & \|ik\|2\hfill \\ \hfill u_ju_iu_k& =& u_ju_ku_i,\hfill & i<jk\hfill & \|ik\|2\hfill \\ \hfill (u_i+u_{i+1})u_{i+1}u_i& =& u_{i+1}u_i(u_i+u_{i+1}).\hfill & & \end{array}$$ In $`FG_n\times [z_1,z_2,\mathrm{},z_m]`$ define the noncommutative Cauchy element to be $$\psi :=\underset{i=1}{\overset{m}{}}\underset{j=n}{\overset{1}{}}(1+z_iu_j).$$ Let $`R`$ be any set whose cardinality is at most countable, and let $`R`$ be the free abelian additive group with basis consisting of the elements of $`R`$. A representation of $`FG_n`$ on $`R`$ is combinatorial if for all $`xR`$, we have $`x.u_iR\{0\}`$. Given a combinatorial representation of $`FG_n`$ on $`R`$ and $`x,yR`$, set $$F_{y/x}:=x.\psi ,y.$$ We define an edge-labelled directed multigraph $``$ with vertex set $`R`$ for which $`F_{y/x}`$ is the quasi-symmetric function coming from the descent Pieri operator on that structure. We construct $``$ by drawing an edge with label $`i`$ from $`x`$ to $`x.u_i`$ $$x\stackrel{i}{–\to }x.u_i$$ if $`x.u_i0`$. Considering the descent Pieri operators on $``$, we have the following. ###### Theorem 6.5. For every $`x,yR`$, $`F_{y/x}=K_{[x,y]}(z_1,\mathrm{},z_m,0)`$. ###### Remark 6.6. We identify the generators $`z_i`$ in $`[z_1,\mathrm{},z_m]`$ with those in the algebra $`𝒬sym`$ generated by the indeterminates $`z_1,z_2,\mathrm{}`$. Before we prove Theorem 6.5, we recall some results from . Define $$e_k(𝐮):=\underset{i_1>i_2>\mathrm{}>i_k}{}u_{i_1}u_{i_2}\mathrm{}u_{i_k}$$ and for a partition $`\lambda =(\lambda _1,\mathrm{},\lambda _m)`$ set $`e_\lambda (𝐮):=e_{\lambda _1}(𝐮)\mathrm{}e_{\lambda _m}(𝐮)`$. ###### Proposition 6.7 (Fomin-Greene). 1. For any positive integers $`a,b`$, we have $`e_a(𝐮)e_b(𝐮)=e_b(𝐮)e_a(𝐮)`$. 2. $`\psi =_\lambda m_\lambda (z)e_\lambda (𝐮)`$. Here, $`m_\lambda (z):=m_\lambda (z_1,\mathrm{},z_m)`$ is the monomial symmetric polynomial. Proof of Theorem 6.5. Observe that for $`xR`$, $$x.e_k(𝐮)=\underset{\stackrel{x\stackrel{i_1}{–\to }\mathrm{}\stackrel{i_k}{–\to }y}{i_1>\mathrm{}>i_k}}{}y=x.\overline{h}_k.$$ From this, and Proposition 6.7 it follows that the Pieri operators are symmetric, that is $`\overline{h}_a\overline{h}_b=\overline{h}_b\overline{h}_a`$ for all $`a,b^+`$. Hence, as in Example 5.2, we have $`x.S^\alpha =x.S^\beta `$ whenever $`\alpha `$ and $`\beta `$ are two compositions that determine the same partition. Then $`F_{y/x}=x.\psi ,y`$ $`=`$ $`{\displaystyle \underset{\lambda }{}}m_\lambda (z)x.e_\lambda (\text{u}),y`$ $`=`$ $`{\displaystyle \underset{\lambda }{}}m_\lambda (z)x.S^\lambda ,y`$ $`=`$ $`x.{\displaystyle \underset{\alpha }{}}M_\alpha (z)S^\alpha ,y`$ $`=`$ $`{\displaystyle \underset{r}{}}x.{\displaystyle \underset{\alpha r}{}}M_\alpha (z)S^\alpha ,y`$ $`=`$ $`{\displaystyle \underset{r}{}}K_{[x,y]^{(r)}}(z_1,\mathrm{},z_m,0)`$ by Theorem 3.3, which by definition is equal to $`K_{[x,y]}(z_1,\mathrm{},z_m,0)`$. ###### Example 6.8. $`P`$-Partitions. Let $`P`$ be a poset and consider any (vertex) labelling $`\gamma :P`$ of $`P`$. A $`(P,\gamma )`$-partition is an order preserving function $`f:P`$ such that if $`x<y`$ and $`\gamma (x)>\gamma (y)`$, then $`f(x)<f(y)`$. It is sufficient to check these conditions for covers $`xy`$ in $`P`$. Let $`𝒜(P,\gamma )`$ be the set of all $`(P,\gamma )`$-partitions. The weight enumerator $`\mathrm{\Gamma }(P,\gamma )`$ of the labelled poset $`(P,\gamma )`$ is $$\mathrm{\Gamma }(P,\gamma ):=\underset{f𝒜(P,\gamma )}{}\underset{xP}{}z_{f(x)}.$$ This is obviously quasi-symmetric. Properties of this weight enumerator are tied up with Stanley’s Fundamental Theorem of $`P`$-partitions . Let $`(P)`$ be the set of all linear extensions of $`P`$. A linear extension $`w`$ of $`P`$ lists the elements of $`P`$ in order $`w_1,w_2,\mathrm{},w_n`$, with $`w_i<w_j`$ (in $`P`$) implying $`i<j`$. Here $`n=\|P\|`$. Let $`D(w,\gamma )`$ be the descent composition of $`n`$ associated to the descent set of the sequence of integers $`\gamma (w_1),\gamma (w_2),\mathrm{},\gamma (w_n)`$. For a linear ordering $`w`$ of $`P`$, the set $`𝒜(w,\gamma )`$ may be identified with the set of all weakly increasing functions $`f:[n]`$ where if $`\gamma (w_i)>\gamma (w_{i+1})`$ then $`f(w_i)<f(w_{i+1})`$. Thus $`\mathrm{\Gamma }(w,\gamma )`$ is Gessel’s fundamental quasi-symmetric function $`F_{D(w,\gamma )}`$. The Fundamental Theorem of $`P`$-partitions notes that $$𝒜(P,\gamma )=\underset{w(P)}{}𝒜(w,\gamma ).$$ This implies that (6.3) $$\mathrm{\Gamma }(P,\gamma )=\underset{w(P)}{}\mathrm{\Gamma }(w,\gamma )=\underset{w(P)}{}F_{D(w,\gamma )}.$$ We show that $`\mathrm{\Gamma }(P,\gamma )`$ is given by descent Pieri operators on the (graded) poset $`P`$ of lower order ideals of P with (edge) labelling induced from the vertex labelling of $`P`$. A subset $`IP`$ is a lower order ideal of $`P`$ if whenever $`xI`$ and $`y<x`$, then $`yI`$. The set $`P`$ of lower order ideals of $`P`$ is ordered by inclusion. We label a cover $`IJ`$ in $`P`$ with $`\gamma (x)`$, where $`x`$ is the unique element $`xJI`$. Then $`(P)`$ is in bijection with the maximal chains of $`P`$. Using the descent Pieri operators for this structure, Equation 6.1 shows that every maximal chain of $`P`$ contributes the summand $`F_{D(w,\gamma )}`$ to $`K_P`$ where $`w`$ is the linear extension of that chain. Thus $$K_P=\mathrm{\Gamma }(P,\gamma ).$$ The Hopf structure of $`(P)`$ was studied by Malvenuto in . ###### Example 6.9. Quantum cohomology of Grassmannian, fusion coefficients, and the Hecke algebra at a root of unity. Let $`m,p`$ be positive integers and let $`𝒞_{p,m}`$ be the set of sequences $`\alpha :0<\alpha _1<\mathrm{}<\alpha _p`$ which also satisfy $`\alpha _p\alpha _1<m+p`$. We order this set of sequences by componentwise comparison to obtain a ranked poset. Given a cover $`\alpha \beta `$, there is a unique index $`i`$ with $`\alpha _i+1=\beta _i`$ and $`\alpha _j=\beta _j`$ for $`ij`$. We label such a cover with $`\beta _i`$. The elements of the poset $`𝒞_{m,p}`$ may alternately be described by pairs $`(a,\lambda )`$, where $`a`$ is a positive integer and $`\lambda `$ is a partition with $`\lambda _{p+1}=0`$ and $`\lambda _1m`$. We obtain $`(a,\lambda )`$ from the sequence $`\alpha `$ by $`\{\lambda _1+p,\mathrm{},\lambda _p+1\}`$ $``$ $`\{\alpha _1,\mathrm{},\alpha _p\}\text{mod}(m+p),`$ $`a(m+p)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{p}{}}}\alpha _i\lambda _ii.\text{}`$ We may likewise pass from the indexing scheme $`(a,\lambda )`$ to sequences $`\alpha `$, as this association is invertible (see ). For $`x𝒞_{m,p}`$ and $`0<kp`$, consider the Pieri operator (6.4) $$x.\overline{h}_k:=\underset{\omega }{}\text{end}(\omega ),$$ where the sum is over all chains $`\omega `$ of length $`k`$ starting at $`x`$, $$\omega :x\stackrel{b_1}{-\to }x_1\stackrel{b_2}{-\to }\mathrm{}\stackrel{b_k}{-\to }x_k=:\mathrm{end}(\omega ),$$ with no descents, that is $`b_1b_2\mathrm{}b_k`$, and also satisfying the restriction $`b_kb_1<m+p`$. Thus these operators $`\overline{h}_k`$ are not an instance of rank-selection or descent Pieri operators as previously introduced. These Pieri operators $`\overline{h}_k`$ are symmetric as they model the Pieri formula in the quantum cohomology ring of the Grassmannian of $`p`$-planes in $`^{m+p}`$. This commutative quantum cohomology ring has a basis $`q^a\sigma _\lambda `$ for $`(a,\lambda )𝒞_{m,p}`$, and $$q^a\sigma _\lambda \sigma _k=\underset{(b,\mu )}{}q^b\sigma _\mu ,$$ where the sum is over all indices $`(b,\mu )`$ appearing in the product $`(a,\lambda ).\overline{h}_k`$ (6.4), when it is written in terms of pairs. Thus we have the following formula $$K_{[(b,\mu ),(a,\lambda )]}=\underset{\nu }{}c_{\mu ,\nu }^\lambda S_\nu ,$$ where the sum is over all partitions $`\nu `$ of $`rk(a,\lambda )rk(b,\mu )`$ with $`\nu _{p+1}=0`$ and $`m\nu _1`$. Here, $`c_{\mu ,\nu }^\lambda `$ is the quantum Littlewood-Richardson coefficient , the coefficient of $`q^a\sigma _\lambda `$ in the product $`q^b\sigma _\mu \sigma _\nu `$. These Pieri operators also model the fusion product in the Verlinde algebra (see for a discussion), and the Pieri formula in the representation rings of Hecke algebras at roots of unity . Geometry and representation theory show that these coefficients $`c_{\mu ,\nu }^\lambda `$ are non-negative, but a combinatorial proof of this fact is lacking. ## 7. Peak enumeration and Eulerian Pieri operators ###### Definition 7.1. Let $`\omega `$ be a labelled ordered chain, that is $$\omega :x_0\stackrel{b_1}{-\to }x_1\stackrel{b_2}{-\to }\mathrm{}\stackrel{b_k}{-\to }x_k.$$ We say that $`\omega `$ has a peak at $`i`$ if $`b_{i1}b_i>b_{i+1}`$. Let $`\mathrm{\Lambda }(\omega )`$ be the peak composition of $`\omega `$, that is the composition of $`k`$ associated to the peak set $`\{i\|b_{i1}b_i>b_{i+1}\}`$ of $`\omega `$. Let $`P`$ be a labelled poset. To enumerate chains in intervals $`[x,y]`$ of $`P`$ according to their peaks, we use the peak enumerator $$\stackrel{~}{K}_{[x,y]}:=\underset{\omega }{}\theta _{\mathrm{\Lambda }(\omega )},$$ where the sum is over all saturated chains $`\omega `$ in the interval $`[x,y]`$. We show this peak enumerator is the quasi-symmetric function $`K_{\delta [x,y]}`$ associated to the descent Pieri operators on an enriched structure $`\delta P`$ defined on the labelled poset $`P`$. Given a labelled poset $`P`$, where (for simplicity) we assume that the labels $`b_i`$ are positive integers, we define $`\delta P`$, the doubling of $`P`$, to be the labelled directed graph with vertex set $`P`$, where every edge $`x\stackrel{b}{-\to }y`$ of $`P`$ is doubled, but with one label the negative of the original label, that is Such a poset whose Hasse diagram has multiple edges is called a réseau. The réseau $`\delta P`$ is the doubled réseau of $`P`$. To define descent Pieri operators on the réseau $`\delta P`$, we say that there is a descent at $`i`$ if consecutive labels $`b_i,b_{i+1}`$ satisfy either $`b_i>b_{i+1}`$ or else $`b_i=b_{i+1}<0`$. We then adjust the definitions of descent set and descent composition accordingly. The following Theorem is a generalisation of \[28, Theorem 3.6\], as will becom apparent from Example 7.5. ###### Theorem 7.2. Let $`P`$ be any labelled poset and $`\delta P`$ its doubled réseau. Then the modified descent Pieri operators on $`\delta P`$ are Eulerian, and we have $$K_{\delta [x,y]}=\stackrel{~}{K}_{[x,y]}=c_{x,\alpha }^y\theta _\alpha .$$ where the sum is only over peak compositions $`\alpha `$. These combinatorial invariants $`c_{x,\alpha }^y`$ of $`\delta P`$ enumerate the chains of $`P`$ whose peak sets have composition $`\alpha `$. Before we prove Theorem 7.2, we make some definitions and prove two auxiliary lemmas. For a composition $`\alpha `$ of $`n`$, let $`\alpha ^+`$ be the composition of $`n+1`$ obtained from $`\alpha `$ by increasing its last component by 1, and $`\alpha 1`$ be the composition of $`n+1`$ obtained by appending a component of size 1 to $`\alpha `$. Define linear maps $`\psi ,\phi :𝒬\text{sym}_n𝒬\text{sym}_{n+1}`$ by $`\psi (M_\beta )`$ $`:=`$ $`M_{\beta ^+}+2M_{\beta 1},`$ $`\phi (M_\beta )`$ $`:=`$ $`\delta _{1,\beta _l}M_{\beta ^+}+2M_{\beta 1},`$ where $`\beta _l`$ is the last component of $`\beta `$ and $`\delta _{1,\beta _l}`$ is the Kronecker delta function. Using the relation $`F_\beta =_{\alpha \beta }M_\alpha `$ between the two bases of $`𝒬\text{sym}`$, we see that $$\psi (F_\beta )=F_{\beta ^+}+F_{\beta 1}.$$ ###### Lemma 7.3. $`\psi (\theta _\alpha )=\theta _{\alpha ^+}`$ and $`\phi (\theta _\alpha )=\theta _{\alpha 1}`$. ###### Proof. The function $`\theta _{\alpha 1}`$ is the sum of terms $`2^{k(\beta )}M_\beta `$ for each $`\beta `$ satisfying $`\beta ^{}\alpha 1`$. Suppose $`\beta ^{}\alpha 1`$. If $`\beta =\gamma 1`$, then $`\beta ^{}=\gamma ^{}1`$ and we have $`\gamma ^{}\alpha `$. Conversely, if $`\gamma ^{}\alpha `$, then $`\beta :=\gamma 1`$ satisfies $`\beta ^{}\alpha 1`$. Thus every summand $`2^{k(\gamma )}M_\gamma `$ of $`\theta _\alpha `$ contributes a summand $`22^{k(\gamma )}M_{\gamma 1}`$ to $`\theta _{\alpha 1}`$. The other summands $`\beta `$ have $`\beta =\gamma \beta _l`$ with $`\beta _l>1`$. Then $`\beta ^{}=\gamma ^{}1(\beta _l1)`$. If $`\beta ^{}\alpha 1`$, then we must have $`\beta _l=2`$, so that $`\beta =\gamma 2`$. Then $`\beta ^{}=\gamma ^{}11\alpha 1`$, which implies that $`(\gamma 1)^{}\alpha `$. Conversely, if $`(\gamma 1)^{}\alpha `$, then $`(\gamma 2)^{}\alpha 1`$. Thus every summand $`2^{k(\gamma 1)}M_{\gamma 1}`$ of $`\theta _\alpha `$ contributes a summand $`2^{k(\gamma 2)}M_{\gamma 2}`$ to $`\theta _{\alpha 1}`$. This shows that $`\theta _{\alpha 1}=\phi (\theta _\alpha )`$. The arguments for $`\theta _{\alpha ^+}`$ are similar, but simpler. The key lemma relating the peak enumerator on $`P`$ and the modified descent Pieri operators on the réseau $`\delta P`$ concerns the case when $`P`$ is a chain. ###### Lemma 7.4. Suppose $`\omega `$ is a chain. Then $`K_{\delta \omega }=\theta _{\mathrm{\Lambda }(\omega )}`$. ###### Proof. We prove this by induction on the length of the chain $`\omega `$. The initial cases are easy calculations. Let $`b_1,\mathrm{},b_k`$ be the word of $`\omega `$, and set $`u`$ to be the truncation of $`\omega `$ at the penultimate cover, so that $`b_1,\mathrm{},b_{k1}`$ is the word of $`u`$. Consider first the case where $`b_{k1}b_k`$. Then every chain $`\gamma `$ in $`\delta u`$ gives two chains $`\gamma .b_k`$ and $`\gamma .\overline{b_k}`$ in $`\delta \omega `$. Since $`D(\gamma .b_k)=D(\gamma )^+`$ and $`D(\gamma .\overline{b}_k)=D(\gamma )1`$, we see that $`K_{\delta \omega }=\psi (K_{\delta u})`$. Similarly, if $`b_{k2}>b_{k1}>b_k`$, then considering the last three labels of a chain in $`\delta \omega `$ show that $`K_{\delta \omega }=\psi (K_{\delta u})`$. In both cases, $`\mathrm{\Lambda }(\omega )=\mathrm{\Lambda }(u)^+`$ (as the peak sets are the same), and the lemma follows by Lemma 7.3. Now suppose $`b_{k2}b_{k1}>b_k`$. Let $`v`$ be the truncation of $`\omega `$ at the $`(k2)`$th position. Let $`\gamma `$ be a chain of $`\delta v`$ with descent composition $`\alpha `$. Then $`\gamma `$ has 4 extensions to chains in $`\delta \omega `$, and 2 have descent composition $`\alpha ^+1`$ and 2 have descent composition $`\alpha 2`$. Thus if we define $`\varphi :F_\alpha 2F_{\alpha ^+1}+2F_{\alpha 2}`$, then $`\varphi (K_{\delta v})=K_{\delta \omega }`$. A straightforward calculation shows $`\varphi (M_\beta )=2M_{\beta ^+1}+2M_{\beta 2}+4M_{\beta 11}`$, which is $`\phi (\psi (M_\beta ))`$. Thus $`K_{\delta \omega }=\phi (\psi (K_{\delta v}))=\phi (K_{\delta u})`$. Since $`\omega `$ has a peak at $`n1`$, we have $`\mathrm{\Lambda }(\omega )=\mathrm{\Lambda }(u)1`$, and so this case follows by Lemma 7.3. Proof of Theorem 7.2. Given an interval $`[x,y]`$ in a poset or réseau, let $`\text{ch}[x,y]`$ be the set of saturated chains in $`[x,y]`$. Let $`K`$ be the quasi-symmetric function given by the descent Pieri operators on the réseau $`\delta P`$. Let $`xy`$ in $`P`$. Given a chain $`\omega \text{ch}\delta [x,y]`$, we obtain a chain $`\|\omega \|\text{ch}[x,y]`$ by replacing each cover in $`\delta [x,y]`$ with a negative integer label by the corresponding cover in $`[x,y]`$ whose label is positive. Then, by Equation 6.1, we have $`K_{\delta [x,y]}`$ $`=`$ $`{\displaystyle \underset{\omega \text{ch}\delta [x,y]}{}}F_{D(\omega )}`$ $`=`$ $`{\displaystyle \underset{\beta \text{ch}[x,y]}{}}{\displaystyle \underset{\omega :\|\omega \|=\beta }{}}F_{D(\omega )}`$ $`=`$ $`{\displaystyle \underset{\beta \text{ch}[x,y]}{}}K_{\delta \beta }={\displaystyle \underset{\beta \text{ch}[x,y]}{}}\theta _{\mathrm{\Lambda }(\beta )}=\stackrel{~}{K}_{[x,y]}.\text{}`$ ###### Example 7.5. Enriched $`P`$-partitions. Stembridge enriches the theory of $`P`$-partitions giving a new class of quasi-symmetric generating functions. Let $`(P,\gamma )`$ be a labelled poset and let $`=\{\overline{1},1,\overline{2},2,\overline{3},3,\mathrm{}\}`$ be two copies of the positive integers ordered as follows: $`\overline{1}<1<\overline{2}<2<\overline{3}<3<\mathrm{}`$. An enriched $`(P,\gamma )`$-partition is an order-preserving map $`f:P`$ such that for $`x<y`$ in $`P`$ and $`k^+`$ * if $`f(x)=f(y)=\overline{k}`$, then $`\gamma (x)<\gamma (y)`$, * if $`f(x)=f(y)=k`$, then $`\gamma (x)>\gamma (y)`$. Let $`(P,\gamma )`$ be the set of all enriched $`(P,\gamma )`$-partitions and define the weight enumerator $$\mathrm{\Delta }(P,\gamma )=\underset{f(P,\gamma )}{}\underset{xP}{}z_{f(x)}$$ where $`z_{\overline{k}}=z_k`$ for all positive integers $`k`$. The analogue of Equation 6.3 for enriched $`P`$-partitions is $$\mathrm{\Delta }(P,\gamma )=\underset{w(P)}{}\mathrm{\Delta }(w,\gamma ).$$ We thus need to characterise the quasi-symmetric function corresponding to a linear extension $`(w,\gamma )`$ of $`(P,\gamma )`$. A peak of the linear extension $`(w,\gamma )`$ is an index $`i`$ with $`1<i<\|w\|`$ where $`\gamma (w_{i1})<\gamma (w_i)>\gamma (w_{i+1})`$. Stembridge shows that $$\mathrm{\Delta }(w,\gamma )=\theta _{\mathrm{\Lambda }(w,\gamma )}$$ where $`\mathrm{\Lambda }(w,\gamma )`$ is the peak composition associated to the peak set of the linear extension $`(w,\gamma )`$. We can then generalise the construction we have for $`P`$-partitions. This time we proceed as in Definition 7.1 and consider the descent Pieri operators on the doubled réseau $`\delta P`$. By Lemma 7.4, every maximal chain of $`P`$ contributes exactly $`\mathrm{\Delta }(w,\gamma )`$ to $`K_{\delta P}`$, where $`w`$ is the linear extension of $`(P)`$ corresponding to that chain. This shows that $$K_{\delta P}=\mathrm{\Delta }(P,\gamma ).$$ ###### Example 7.6. Isotropic Pieri formula. The Pieri-type formulas for the flag manifolds $`SO(2n+1,)`$ and $`Sp(2n,)`$ of each give symmetric Eulerian Pieri operators. These are defined on enrichments of the same subposet of the Bruhat order on the group $`_n`$ of signed permutations. For an integer $`i`$, let $`\overline{ı}`$ denote $`i`$. We regard $`_n`$ as a subgroup of the group of permutations on $`\{\overline{n},\mathrm{},\overline{2},\overline{1},1,\mathrm{},n\}`$. Let $`\mathrm{}`$ be the length function on the Coxeter group $`_n`$. The $`0`$-Bruhat order $`<_0`$ on $`_n`$ is the labelled poset $`_n^0`$ with covers $`u_0w`$ if $`\mathrm{}(u)+1=\mathrm{}(w)`$ and $`u^1w`$ is a reflection with either the form $`(\overline{ı},i)`$ or the form $`(\overline{ı},j)(\overline{ȷ},i)`$ for some $`0<i,j`$. When $`u_0w`$, either $`wu^1=(\overline{\beta },\beta )`$ for some $`0<\beta `$ or else $`wu^1=(\overline{\beta },\overline{\alpha })(\alpha ,\beta )`$ for some $`0<\alpha <\beta n`$. We label a such a cover with the (positive) integer $`\beta `$. Consider the Eulerian descent Pieri operators on the doubled réseau $`\delta _n^0`$. These operators are symmetric, as $`\frac{1}{2}\overline{h}_k`$ models the action of the Schur P-polynomial $`p_k`$ on the basis of Schubert classes (indexed by $`_n`$) in the cohomology of the flag manifold $`SO(2n+1,)/B`$ . (This is because there are twice as many increasing chains in a doubled interval $`\delta [x,y]`$ as peakless chains in the interval $`[x,y]`$, and the coefficient of $`y`$ in $`x.p_k`$ is this number of peakless chains.) We modify this descent action of $`NC`$ on $`\delta _n^0`$ by identifying $`h_k`$ with $`\frac{1}{2}\overline{h}_k`$, which is still integral. These new Pieri operators are symmetric, as they model the action of $`p_k`$, and they are Eulerian, as $`2h_k`$ satisfies the Euler relations 5.2. In exact analogy to how the Skew Schubert functions are shown in to be the generating functions for the coefficients $`c_{u,(\lambda ,k)}^w`$, given by descent Pieri operators, we have the following formula $$K_{[u,w]}=\underset{\lambda }{}b_{u,\lambda }^wQ_\lambda ,$$ where the sum is over all strict partitions $`\lambda `$ of $`\mathrm{}(w)\mathrm{}(u)`$. Here $`b_{u,\lambda }^w`$ is the coefficient of the Schubert class $`𝔅_w`$ in the product $`𝔅_uP_\lambda `$, and $`P_\lambda ,Q_\lambda `$ are Schur P- and Q-polynomials, which form dual bases for the self dual symmetric Hopf algebra $`\mathrm{\Pi }_{}\mathrm{\Lambda }`$. The polynomials $`Q_\lambda `$ appear as $`_\lambda P_\lambda Q_\lambda `$ is the Cauchy element of $`\mathrm{\Pi }_{}\mathrm{\Lambda }`$. For the symplectic flag manifold, we modify the réseau $`\delta _n^0`$ by erasing the negative edge in a cover when $`wu^1=(\overline{\beta },\beta )`$. Write $`_n^0`$ for the resulting réseau. It is a slight modification of the 0-Bruhat réseau of , and may be used in its place for the combinatorics therin. Let $`\{\overline{h}_k\}`$ be the descent Pieri operators on $`_n^0`$. This family of Pieri operators is symmetric and Eulerian, as $`\overline{h}_k`$ models the action of the Schur Q-polynomial $`q_k`$ on the Schubert basis of the cohomology of the flag manifold $`Sp(2n,)`$, and the Schur Q-polynomials $`q_k`$ satisfy the Euler relations. As before, we have the following formula $$K_{[u,w]}=\underset{\lambda }{}c_{u,\lambda }^wP_\lambda ,$$ where the sum is over all strict partitions $`\lambda `$ of $`\mathrm{}(w)\mathrm{}(u)`$. Here $`c_{u,\lambda }^w`$ is the coefficient of the Schubert class $`_w`$ in the product $`_uQ_\lambda `$. Since every chain in an interval of $`_n^0`$ has the same number of covers of the form $`(\overline{\beta },\beta )`$—these count the number $`s(wu^1)`$ of sign changes between $`u`$ and $`w`$, we have $$K_{\delta [u,w]}=2^{s(wu^1)}K_{[u,w]}.$$ Geometry shows these coefficients $`b_{u,\lambda }^w`$ and $`c_{u,\lambda }^w`$ are non-negative. It is an important open problem to give a combinatorial or algebraic proof of this fact. ###### Example 7.7. Stanley symmetric functions of types $`B`$, $`C`$, and $`D`$. In , Billey and Haiman describe the Stanley symmetric functions of types $`B`$ and $`D`$ in terms of peaks of reduced words of elements in the corresponding Coxeter groups. For $`_n`$, the simple transpositions are $`s_0,s_1,\mathrm{},s_{n1}`$, where $`s_0=(\overline{1},1)`$ and if $`i>0`$, then $`s_i=(\overline{i+1},\overline{ı})(i,i+1)`$. The weak order on $`_n`$ is the labelled poset whose covers are $`wws_i`$ with label $`i+1`$ if $`\mathrm{}(w)+1=\mathrm{}(ws_i)`$. A reduced word $`𝐚`$ for $`w`$ is sequence of labels of a chain in $`_n`$ from the identity $`e`$ to $`w`$. Billey and Haiman define the Stanley symmetric function of type $`B`$ to be $$F_w^B:=\underset{𝐚R(w)}{}\theta _{\mathrm{\Lambda }(𝐚)},$$ where $`R(w)`$ is the set of reduced words for $`w`$ and $`\mathrm{\Lambda }(𝐚)`$ is the peak composition of the reduced word $`𝐚`$. By Theorem 7.2, $`F_w^B`$ is the function $`K_{\delta [e,w]}`$ obtained from the Eulerian descent operators on the doubled réseau $`\delta _n`$. Billey and Haiman establish the formula $$F_w^B=\underset{\lambda }{}f_\lambda ^wQ_\lambda ,$$ where the sum is over all strict partitions $`\lambda `$ of $`\mathrm{}(w)`$, and $`f_\lambda ^w`$ counts the reduced words that satisfy a condition imposed by the partition $`\lambda `$ (coming from the shifted Edelmann-Greene correspondence ). Thus the Eulerian descent Pieri operators on $`\delta _n`$ are also symmetric. While Billey and Haiman do not define Stanley functions of type $`C`$, one reasonably sets $`F_w^C:=2^{s(w)}F_w^B`$, where $`s(w)`$ is the number of sign changes in the permutation $`w`$. This is just the number of $`s_0`$’s appearing in any reduced word of $`w`$. Let the réseau $`_n`$ be the modification of the doubled réseau $`\delta _n`$ where we erase the edge with negative label $`\overline{1}`$ for covers $`wws_0`$. Then every chain in an interval $`[e,w]`$ of $`_n`$ gives rise to $`2^{s(w)}`$ chains in $`\delta [x,y]`$, each with the same descents as the original chain. Then Equation 6.1 and Theorem 7.2 show that $$K_{[x,y]}=2^{s(w)}K_{\delta [x,y]}=2^{s(w)}F_w^B=F_w^C.$$ This shows these descent Pieri operators are Eulerian and symmetric. The Coxeter group $`𝒟_n`$ has simple reflections $`s_1,s_{\widehat{1}},s_2,\mathrm{},s_{n1}`$. The weak order on $`𝒟_n`$ is the labelled poset with cover $`wws_i`$ labelled by $`i`$ if $`\mathrm{}(w)+1=\mathrm{}(ws_i)`$. Here, we set $`\widehat{1}<1`$. A reduced word $`𝐚`$ for $`w`$ as before is a chain in $`𝒟_n`$ from $`e`$ to $`w`$. Since $`s_1`$ and $`s_{\widehat{1}}`$ commute, there are no occurences of $`1\widehat{1}1`$ or $`\widehat{1}1\widehat{1}`$ in a reduced word, so changing all occurrences of $`\widehat{1}`$ to $`1`$ does not change the peaks in a reduced word, and the type $`D`$ Stanley symmetric functions of Billey and Haiman satisfy $$F_w^D=\underset{𝐚R(w)}{}2^{o(𝐚)}\theta _{\mathrm{\Lambda }(𝐚)},$$ where $`o(𝐚)`$ counts the number of occurrences of $`1`$ and $`\widehat{1}`$ in the reduced word $`𝐚`$. Let $`\delta 𝒟_n`$ and $`𝒟_n`$ be the doubled réseau and its modification, erasing all edges with (negative) labels $`1`$ and $`\widehat{1}`$. ###### Theorem 7.8. $`F_w^D=K_{[e,w]}`$. ###### Proof. By Theorem 7.2, we have $$F_w^D=\underset{𝐚R(w)}{}2^{o(𝐚)}K_{\delta 𝐚}.$$ The theorem follows from Equation 6.1 and the following 1 to $`2^{o(𝐚)}`$ map from chains in $`𝐚`$ to chains in $`\delta 𝐚`$, which preserves descents. When there are no subwords $`1\widehat{1}`$ or $`\widehat{1}1`$ in a chain in $`𝐚`$, simply make all possible substitutions of negative and positive labels for each occurrence of $`1`$ and $`\widehat{1}`$. If however, there is a subword $`1\widehat{1}`$, then there is another chain differing from the first only in that subword (having $`\widehat{1}1`$ instead), and the map uses the substitutions in both chains $`\widehat{1}1`$ $``$ $`\widehat{1}1,\overline{\widehat{1}}1,\overline{1}\widehat{1},\overline{1}\overline{\widehat{1}}`$ $`1\widehat{1}`$ $``$ $`1\widehat{1},1\overline{\widehat{1}},\widehat{1}\overline{1},\overline{\widehat{1}}\overline{1}.\text{}`$ Lastly, we remark that these descent operators on $`𝒟_n`$ are Eulerian and symmetric, as Billey and Haiman give a formula $$F_w^D=\underset{\lambda }{}e_w^\lambda Q_\lambda ,$$ where $`e_w^\lambda `$ is a rational number that counts certain weighted reduced words. ## Appendix A Hopf algebras A Hopf algebra is an algebra whose linear dual is also an algebra, with some compatibility conditions. They are important in representation theory (and in this paper) because they act on tensor products of their representations. The usefulness of Hopf algebras in combinatorics is apparent from the ubiquity of their applications. In this section, we summarize the basic notions of Hopf algebras. A $``$-module $``$ is a coalgebra if there are two maps $`\mathrm{\Delta }:`$ (coproduct) and $`ϵ:`$ (counit or augmentation) such that the following diagrams commute $$\text{}\text{},$$ where $`1`$ is the identity map on $``$. ###### Remark A.1. The first of these diagrams is the coassociativity property, which is the statement that the dual of $`\mathrm{\Delta }`$ defines an associative product on the linear dual of $``$, and the second asserts this linear dual has a unit, induced by the dual of $`ϵ`$. If $``$ is also an algebra, then it is a bialgebra if $`\mathrm{\Delta },ϵ`$ are algebra morphisms. While some authors call this structure a Hopf algebra, we define a Hopf algebra to be a bialgebra with a map $`s:`$ (coinverse or antipode) such that the following diagram commutes. $$\text{}.$$ Here $`\mu :`$ is the map induced by the multiplication of $``$ and $`u:`$ is the map induced by mapping $`1`$ to the unit of $``$. The above diagram implies that $`s`$ is an algebra antihomomorphism, i.e. $`s(hh^{})=s(h^{})s(h)`$ for all $`h,h^{}`$. The existence of an antipode $`s`$ may seem to be a strong restriction on a bialgebra, however as we will see, it is no restriction for graded bialgebras. A graded bialgebra is a graded algebra $`=_n_n`$ where $`\mathrm{\Delta }`$ is graded and $`_0=`$. Given $`x_n`$, the $`n`$th graded component, we have $$\mathrm{\Delta }(x)=x1+\underset{i=1}{\overset{n}{}}y_iz_{ni},$$ where $`y_i`$ and $`z_i`$ have degree $`i`$. The first term is always present due to the counit diagram. With this in mind, Ehrenborg proved the following. ###### Proposition A.2 (Lemma 2.1 ). Given a graded bialgebra $``$ there is a unique Hopf algebra with antipode $`s`$ defined recursively by $`s(1)=1`$, and for $`x_n`$, $`n1`$, $$s(x)=\underset{i=1}{\overset{n}{}}s(y_i)z_i.$$ Lastly, we remark on the useful Sweedler notation, which is an elegant solution to the following quandary. Given $`h`$, how do you efficiently represent $`\mathrm{\Delta }h`$ as an element of $``$? Carefully indexing this element would confuse even the writer. Sweedler notation sidesteps this by omitting the indices of summation entirely, $$\mathrm{\Delta }h=h_1h_2.$$ It is this notation that is normally used when dealing with Hopf algebras.
warning/0002/quant-ph0002093.html	ar5iv	text	# Double jumps and transition rates for two dipole-interacting atoms ## I Introduction The dipole-dipole interaction between two atoms can be understood through the exchange of virtual photons and depends on the transition dipole moment of the levels involved. It can be characterized by complex coupling constants, or by their real and imaginary parts, where the former affect decay constants and the latter lead to level shifts aga . Cooperative effects in the radiative behavior of atoms which may arise from their mutual dipole-dipole interaction have attracted considerable interest in the literature aga -Ho . Two of the present authors BeHe4 have investigated in detail the transition from anti-bunching to bunching with decreasing atomic distance for two dipole-dipole interacting two-level atoms. The striking phenomenon of macroscopic quantum jumps (electron shelving or macroscopic dark and light periods) can occur for a multi-level system where the electron is essentially shelved for seconds or even minutes in a metastable state without photon emissions SauterL -HePle1 . For two such systems the fluorescence behavior would, without cooperative effects, be just the sum of the separate photon emissions, with dark periods of both atoms, light periods of a single atom and of two atoms. In Ref. Sauter the fluorescence intensity of three such ions in a Paul trap was measured and a large fraction of double or triple jumps was reported, i.e. jumps by two or three intensity steps within the short resolution time. This fraction was orders of magnitudes larger than that expected for independent ions. A quantitative explanation of such a large cooperative effect for distances of the order of ten wave lengths of the strong transition has been found to be difficult hendriks ; Java ; Aga88 ; Lawande89 ; Chun . Other experiments at larger distances and with different ions showed no cooperative effects Itano88 ; Thom . Quite recently, two of the present authors BeHe5 investigated for two such systems cooperative effects in the mean duration, $`T_0`$, $`T_1`$, and $`T_2`$, of the dark, single-intensity, and double-intensity periods, respectively. This was done by simulations for two atoms in a $`V`$ configuration. The mean duration of the single- and double-intensity periods depended sensitively on the dipole-dipole interaction and thus on the atomic distance $`r`$. They exhibited noticeable oscillations which decreased in amplitude when $`r`$ increased. These oscillations seemed to continue up to a distance of well over five wave lengths of the strong transition and they were opposite in phase with those of Re$`C_3(r)`$, where $`C_3`$ is the complex dipole-dipole coupling constant associated with the strong transitions. In this paper we present an analytic approach to study cooperative effects for atoms with macroscopic quantum jumps. This is explained for two atoms, but is easily generalized. The approach is based on an explicit calculation of transition rates between the various intensity periods. From the transition rates all interesting statistical quantities can be determined, such as double jump rates and mean duration of different periods. We predict, to our knowledge for the first time, cooperative effects in the double jumps of two dipole-dipole interacting atoms. These results are for atoms in the $`V`$ configuration (see Fig. 1) and are verified by simulations. As a function of the atomic distance, the double jump rates show marked oscillations, with a maximal difference of up to 30%, decreasing as $`1/r`$. Most surprising is a change in the oscillatory behavior of the double jump rate from in phase with Re$`C_3(r)`$ to opposite in phase when the detuning of the laser driving the weak atomic transition is increased. For the mean durations $`T_1`$ and $`T_2`$ there can be a change in behavior from opposite in phase to in phase with Re$`C_3(r)`$. Moreover, for a particular value of the detuning, which depends on the other parameters, the double jump rate becomes constant in $`r`$ and the cooperative effects disappear. This is true also for the mean period durations and for their mean occurrences, with different values of the detuning, though. Typically, for nonzero detuning the oscillation amplitudes do not exceed those found for zero detuning. The experiments of Ref. Sauter exhibited extremely large cooperative effects, in fact up to three orders of magnitude. Since this was for a different atomic level configuration and for three ions in a trap our results do not apply directly. In principle, however, our analytic approach can be carried over to the experimental situation of Ref. Sauter , although the calculations become algebraically more involved and have not been carried out so far. The plan of the paper is as follows. In Section II the fluorescence with its three different intensity periods is treated as a three-step telegraph process and the Bloch equations are used to derive the transition rates between the periods. In Sections III and IV expressions for the double jump rate and the mean duration of the three types of intensity periods are obtained by means of these transition rates. The results are compared with simulations in which photon intensities are obtained by averaging photon numbers over a small time window. It turns out that this data-smoothing procedure can affect the results, and we show how this can be corrected for quantitatively. A similar effect can also occur when photon detectors measure the intensity of light by averaging over a small time window. In the last section the results are discussed. It is suggested that the mean rate of double-intensity periods is an experimentally more easily accessible candidate for exhibiting cooperative effects arising from the dipole-dipole interaction. ## II Transition rates ### II.1 Prerequisites We consider two atoms, at a fixed distance r, each a $`V`$ configuration as shown in Fig. 1. We assume the laser radiation normal to this line, and for the Einstein coefficients, the Rabi frequencies and the detuning we assume the relations $`\mathrm{\Omega }_2\mathrm{\Omega }_3,\mathrm{\Omega }_2\mathrm{\Omega }_3^2/A_3,A_20,\mathrm{\Delta }_3=0,`$ (1) $`\mathrm{\Delta }_2`$ arbitrary. The Dicke states are defined as $`\|g`$ $`=`$ $`\|1\|1,\|e_2=\|2\|2,\|e_3=\|3\|3`$ $`\|s_{jk}`$ $`=`$ $`\{\|j\|k+\|j\|k\}/\sqrt{2}`$ $`i\|a_{jk}`$ $`=`$ $`\{\|j\|k\|j\|k\}/\sqrt{2}`$ They are symmetric and antisymmetric, respectively, under permutation of the two atoms. The Dicke states and the possible transitions are displayed in Fig. 2. Solid single and double arrows indicate decay and strong driving by laser 3, respectively, while dashed double arrows indicate the weak driving by laser 2. For $`\mathrm{\Omega }_2=0`$, i.e. with the dashed arrows absent, the states decompose into three non-connected subsets, namely $`\|e_2`$, the four states of the inner ring and the four states of the outer ring in Fig. 2, and the subspaces spanned by these states will be denoted by dark, inner, and outer subspace, respectively. As in Ref. BeHe5 they will be associated in the following with the fluorescence periods of intensity 0, 1, and 2: $`\mathrm{dark}\mathrm{state}`$ $`:`$ $`\|e_2`$ (2) $`\mathrm{inner}\mathrm{states}(\mathrm{intensity}1)`$ $`:`$ $`\|s_{12},\|s_{23},\|a_{12},\|a_{23}`$ (3) $`\mathrm{outer}\mathrm{states}(\mathrm{intensity}2)`$ $`:`$ $`\|g,\|s_{13},\|e_3,\|a_{13}`$ (4) The weak laser will lead to slow transitions between the subspaces. The Bloch equations can be written, with the conditional Hamiltonian $`H_{\mathrm{cond}}`$ and the reset operation $``$ of Appendix A, in the compact form BeHe4 ; He $$\dot{\rho }=\frac{\mathrm{i}}{\mathrm{}}[H_{\mathrm{cond}}\rho \rho H_{\mathrm{cond}}^{}]+(\rho ).$$ (5) The operator $`H_{\mathrm{cond}}`$ is of the form $$H_{\mathrm{cond}}=H_{\mathrm{cond}}^0+H_{\mathrm{cond}}^1(\mathrm{\Omega }_2)$$ (6) where the operator $`H_{\mathrm{cond}}^0`$ depends on $`\mathrm{\Omega }_3`$ and on the dipole-dipole coupling constant $`C_3(r)`$, while $`H_{\mathrm{cond}}^1`$ is linear in $`\mathrm{\Omega }_2`$ and does not depend on $`C_3`$ and $`\mathrm{\Omega }_3`$. The super-operator $``$ depends on $`C_3`$. ### II.2 Intensity periods and subspaces For a single atom as in Fig. 1, with macroscopic light and dark periods, the stochastic sequence of individual photon emissions can be directly analyzed by the quantum jump approach HeWi ; Wi ; He ; HeSo ; MC ; QT ; PleKni , using the existence of different time scales. To high precision it yields a telegraph process and the transition rates between the periods BeHe1 ; BeHeSo . A more heuristic approach assumes that during a light period the density matrix of the atom lies in the subspace spanned by $`\|1`$ and $`\|3`$ and that during a dark period the state is given by $`\|2`$ Cook . One can then use the Bloch equations to calculate the build-up, during a time $`\mathrm{\Delta }t`$, of a population outside the respective subspace and obtains from this the probability of leaving the subspace. This probability is then interpreted as the transition probability from one period to the other. The results agree with those of the more microscopic quantum jump approach HeWi ; Wi ; HePle1 ; HePle2 . This idea will be used here for two dipole-interacting $`V`$ systems. We associate each of the three types of fluorescence periods with one of the subspaces spanned by the states in Eq. (2) - (4) and model transitions between periods as transitions between the corresponding subspaces. Without dipole interaction this is the same assumption as for a single atom, and with the interaction it has been tested numerically in Ref. beige to hold as long as the atomic separation is larger than a third wavelength of the strong transition. Thus, at a particular time $`t_0`$, the density matrix $`\rho (t_0)`$ of the two atoms is assumed to lie in one of the subspaces. Then, during a short time $`\mathrm{\Delta }t`$, satisfying $$\mathrm{\Omega }_3^1,A_3^1\mathrm{\Delta }t\mathrm{\Omega }_2^1,$$ (7) the system will go over to a density matrix $`\rho (t_0+\mathrm{\Delta }t)`$ which contains small populations in the other subspaces, due to the driving by $`\mathrm{\Omega }_20`$. The time derivatives of these populations at $`t_0+\mathrm{\Delta }t`$ give the transition rates to these subspaces because, as shown in Appendix B, they are independent of the particular choice of $`\mathrm{\Delta }t`$ and of the particular density matrix $`\rho (t_0)`$, as long as Eq. (7) is fulfilled. These rates can be interpreted as transition rates between corresponding intensity periods, just as in the one-atom case. A straightforward calculation using Eq. (5) yields the exact relations $`{\displaystyle \frac{\mathrm{d}}{\mathrm{dt}}}{\displaystyle \underset{\mathrm{outer}}{}}\mathrm{outer}\|\rho \|\mathrm{outer}`$ (8) $`=\mathrm{\Omega }_2\mathrm{Im}\left\{\sqrt{2}s_{12}\|\rho \|g+s_{23}\|\rho \|s_{13}+a_{23}\|\rho \|a_{13}\right\}`$ $`{\displaystyle \frac{\mathrm{d}}{\mathrm{dt}}}e_2\|\rho \|e_2=\sqrt{2}\mathrm{\Omega }_2\mathrm{Im}s_{12}\|\rho \|e_2`$ (9) $`{\displaystyle \frac{\mathrm{d}}{\mathrm{dt}}}{\displaystyle \underset{\mathrm{inner}}{}}\mathrm{inner}\|\rho \|\mathrm{inner}`$ (10) $`={\displaystyle \frac{\mathrm{d}}{\mathrm{dt}}}\left\{e_2\|\rho \|e_2+{\displaystyle \underset{\mathrm{outer}}{}}\mathrm{outer}\|\rho \|\mathrm{outer}\right\}`$ where $`\|\mathrm{outer}`$ stands for $`\|g`$, $`\|s_{13}`$, $`\|e_3`$, $`\|a_{13}`$ and $`\|\mathrm{inner}`$ for $`\|s_{12},\|s_{23},\|a_{12},\|a_{23}`$. Thus one has to calculate the coherences on the right-hand side at time $`t_0+\mathrm{\Delta }t`$ to first order in $`\mathrm{\Omega }_2`$, with the appropriate initial condition at time $`t_0`$, to obtain the transition rate to second order in $`\mathrm{\Omega }_2`$. If $`\rho (t_0)`$ lies in one of the above subspaces then by time $`t_0+\mathrm{\Delta }t`$ the system has reached a quasi-stationary state satisfying $$\dot{\rho }(t_0+\mathrm{\Delta }t)=0\mathrm{to}\mathrm{first}\mathrm{order}\mathrm{in}\mathrm{\Omega }_2,$$ (11) as shown in Appendix B. This is also true for a single atom and is the decisive equation. To obtain from this the coherences to first order in $`\mathrm{\Omega }_2`$ we write $$\rho (t_0+\mathrm{\Delta }t)=\rho ^0+\rho ^1+\mathrm{}$$ where $`\rho ^k`$ is of order $`\mathrm{\Omega }_2^k`$. Putting $`\dot{\rho }=0`$ in Eq. (5) and inserting the expansion for $`\rho `$ one obtains in zeroth order $`0`$ $`=`$ $`{\displaystyle \frac{\mathrm{i}}{\mathrm{}}}\left[H_{\mathrm{cond}}^0\rho ^0\rho ^0H_{\mathrm{cond}}^0\right]+(\rho ^0)`$ (12) and in first order in $`\mathrm{\Omega }_2`$ $`0={\displaystyle \frac{\mathrm{i}}{\mathrm{}}}\left[H_{\mathrm{cond}}^0\rho ^1\rho ^1H_{\mathrm{cond}}^0+H_{\mathrm{cond}}^1\rho ^0\rho ^0H_{\mathrm{cond}}^1\right]`$ (13) $`+(\rho ^1).`$ Thus $`\rho ^0`$ is an equilibrium state for $`\mathrm{\Omega }_2=0`$, taken to lie in the appropriate subspace. For the dark state and the subspace spanned by the inner states one has $`\rho ^0\rho _{\mathrm{dark}}^0`$ $`=`$ $`\|e_2e_2\|`$ (14) $`\rho ^0\rho _{\mathrm{inner}}^0={\displaystyle \frac{1}{2}}\left\{\rho _{\mathrm{ss}}^{(A)}\|22\|+\|22\|\rho _{\mathrm{ss}}^{(B)}\right\}`$ (15) $`={\displaystyle \frac{1}{4}}{\displaystyle \frac{A_3^2+\mathrm{\Omega }_3^2}{A_3^2+2\mathrm{\Omega }_3^2}}\{\|s_{12}s_{12}\|+\|a_{12}a_{12}\|\}`$ $`+{\displaystyle \frac{1}{4}}{\displaystyle \frac{\mathrm{\Omega }_3^2}{A_3^2+2\mathrm{\Omega }_3^2}}\{\|s_{23}s_{23}\|+\|a_{23}a_{23}\|\}`$ $`+{\displaystyle \frac{i}{2}}{\displaystyle \frac{\mathrm{\Omega }_3A_3}{A_3^2+2\mathrm{\Omega }_3^2}}\{\|s_{12}s_{23}\|\|a_{12}a_{23}\|\}+\mathrm{H}.\mathrm{c}.`$ by symmetry, independently of $`C_3`$, where $`\rho _{\mathrm{ss}}^{(A,B)}`$ are the steady states of the individual atoms in the $`1`$-$`3`$ subspace (for $`\mathrm{\Omega }_2=0`$ and $`C_3=0`$). For the subspace spanned by the outer states one calculates $`\rho ^0\rho _{\mathrm{outer}}^0[\{(A_3^2+\mathrm{\Omega }_3^2)^2+A_3^2\|C_3\|^2+2A_3^3\mathrm{Re}C_3\}\|gg\|`$ (16) $`+\{i\sqrt{2}A_3\mathrm{\Omega }_3(A_3^2+\mathrm{\Omega }_3^4+A_3C_3)\|gs_{13}\|+\mathrm{H}.\mathrm{c}.\}`$ $`\{A_3\mathrm{\Omega }_3^2(A_3+C_3)\|ge_3\|+\mathrm{H}.\mathrm{c}.\}`$ $`+\mathrm{\Omega }_3^2(2A_3^2+\mathrm{\Omega }_3^2)\|s_{13}s_{13}\|+\mathrm{\Omega }_3^4\left\{\|e_3e_3\|+\|a_{13}a_{13}\|\right\}`$ $`+\{i\sqrt{2}A_3\mathrm{\Omega }_3^3\|s_{13}e_3\|+\mathrm{H}.\mathrm{c}.\}]`$ One checks that for $`C_3=0`$ this becomes $`\rho _{\mathrm{outer}}^0\rho _{\mathrm{ss}}^{(A)}\rho _{\mathrm{ss}}^{(B)}`$, the expression for two independent atoms. We will denote the transition rates between the subspaces by $`p_{ij}`$. Here $`i,j=0,1,2`$ refer to the dark, inner and outer subspace, respectively, (and thus to the corresponding intensities). The $`p_{ij}`$ will be determined to second order in $`\mathrm{\Omega }_2`$. As expected, $`p_{02}`$ and $`p_{20}`$ will turn out to be zero. ### II.3 Calculation of $`p_{12}`$ We start from $`\rho ^0=\rho _{\mathrm{inner}}^0`$ in Eq. (15) as initial state. For the the transition rate $`p_{12}`$ to the outer subspace one needs, in view of Eq. (8), three coherences of $`\rho ^1`$ between the inner and outer subspace. To obtain these we write $`\{\|x_i\}=\{\|s_{12},\|s_{23},\|a_{12},\|a_{23}\}`$ (inner states) and $`\{\|y_j\}=\{\|g,\|s_{13},\|e_3,\|a_{13}\}`$ (outer states) for the corresponding bases and decompose $`\rho ^1`$ $`=`$ $`{\displaystyle \underset{i,j}{}}\rho _{ij}^1\|x_iy_j\|+\rho _{ij}^1\|y_jx_i\|+\mathrm{other}\mathrm{terms}.`$ (17) Inserting this into Eq. (13) and taking matrix elements with $`x_{i_0}\|`$ on the left and $`\|y_{j_0}`$ on the right gives $`0`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{}}}x_{i_0}\|\rho _{\mathrm{inner}}^0H_{\mathrm{cond}}^1\|y_{j_0}{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle \underset{i}{}}\rho _{ij_0}^1x_{i_0}\|H_{\mathrm{cond}}^0\|x_i`$ (18) $`+{\displaystyle \frac{i}{\mathrm{}}}{\displaystyle \underset{j}{}}\rho _{i_0j}^1y_j\|H_{\mathrm{cond}}^0\|y_{j_0}`$ $`+{\displaystyle \underset{i,j}{}}(A_3+\mathrm{Re}C_3)\rho _{ij}^1x_{i_0}\|R_+\|x_iy_j\|R_+^{}\|y_{j_0}`$ $`+{\displaystyle \underset{i,j}{}}(A_3\mathrm{Re}C_3)\rho _{ij}^1x_{i_0}\|R_{}\|x_iy_j\|R_{}^{}\|y_{j_0}.`$ This is a system of 16 linear equations for the 16 coherences $`\rho _{ij}^1`$, of which only three are needed in Eq. (8). Due to the symmetry of $`H_{\mathrm{cond}}`$ and $`R_+`$ and antisymmetry of $`R_{}`$ under the interchange of the two atoms, the system decouples. Taking for $`\|x_{i_0}`$ and $`\|y_{j_0}`$ either both symmetric or both antisymmetric states and putting the eight coherences into the column vector $$\stackrel{~}{\rho }(\rho _{s_{12}g}^1,\rho _{s_{12}s_{13}}^1,\rho _{s_{12}e_3}^1,\rho _{s_{23}g}^1,\rho _{s_{23}s_{13}}^1,\rho _{s_{23}e_3}^1,\rho _{a_{12}a_{13}}^1,\rho _{a_{23}a_{13}}^1)^T$$ (19) one obtains the equation $$(𝐀i\mathrm{\Delta }_2\mathrm{𝟏})\stackrel{~}{\rho }=𝐚_1$$ (20) where $`𝐀=`$ (30) $`\left[\begin{array}{cccccccc}0& i\mathrm{\Omega }_3/\sqrt{2}& 0& i\mathrm{\Omega }_3/2& (A_3+\mathrm{Re}C_3)/\sqrt{2}& 0& 0& (\mathrm{Re}C_3A_3)/\sqrt{2}\\ i\mathrm{\Omega }_3/\sqrt{2}& (A_3+C_3^{})/2& i\mathrm{\Omega }_3/\sqrt{2}& 0& i\mathrm{\Omega }_3/2& (A_3+\mathrm{Re}C_3)/\sqrt{2}& 0& 0\\ 0& i\mathrm{\Omega }_3/\sqrt{2}& A_3& 0& 0& i\mathrm{\Omega }_3/2& 0& 0\\ i\mathrm{\Omega }_3/2& 0& 0& A_3/2& i\mathrm{\Omega }_3/\sqrt{2}& 0& 0& 0\\ 0& i\mathrm{\Omega }_3/2& 0& i\mathrm{\Omega }_3/\sqrt{2}& (A_3+C_3^{}/2)& i\mathrm{\Omega }_3/\sqrt{2}& 0& 0\\ 0& 0& i\mathrm{\Omega }_3/2& 0& i\mathrm{\Omega }_3/\sqrt{2}& 3A_3/2& 0& 0\\ 0& 0& 0& 0& 0& (A_3\mathrm{Re}C_3)/\sqrt{2}& (A_3C_3^{})/2& i\mathrm{\Omega }_3/2\\ 0& 0& 0& 0& 0& 0& i\mathrm{\Omega }_3/2& (A_3C_3^{}/2)\end{array}\right]`$ $`𝐚_1={\displaystyle \frac{i\mathrm{\Omega }_2\mathrm{\Omega }_3}{4(A_3^2+2\mathrm{\Omega }_3^2)}}[\sqrt{2}{\displaystyle \frac{\mathrm{\Omega }_3^2+A_3^2}{\mathrm{\Omega }_3}},iA_3,0,i\sqrt{2}A_3,\mathrm{\Omega }_3,0,iA_3,\mathrm{\Omega }_3]^T.`$ (31) Inverting the $`8\times 8`$ matrix $`𝐀i\mathrm{\Delta }_2\mathrm{𝟏}`$ by Maple yields $`\stackrel{~}{\rho }`$ and the coherences. The result is complicated and not illuminating. Inserting the required coherences into Eq. (8) one obtains, to first order in Re$`C_3`$ and Im$`C_3`$ and to second order in $`\mathrm{\Omega }_2`$, $`p_{12}=\mathrm{\Omega }_2^2\{{\displaystyle \frac{A_3\mathrm{\Omega }_3^2}{\mathrm{\Omega }_3^48\mathrm{\Delta }_2^2\mathrm{\Omega }_3^2+4A_3^2\mathrm{\Delta }_2^2+16\mathrm{\Delta }_2^4}}`$ (32) $`+\mathrm{Re}C_3(r){\displaystyle \frac{2A_3^2\mathrm{\Omega }_3^2(\mathrm{\Omega }_3^44A_3^2\mathrm{\Delta }_2^216\mathrm{\Delta }_2^4)}{(A_3^2+2\mathrm{\Omega }_3^2)(\mathrm{\Omega }_3^48\mathrm{\Delta }_2^2\mathrm{\Omega }_3^2+4A_3^2\mathrm{\Delta }_2^2+16\mathrm{\Delta }_2^4)^2}}\}.`$ Note that only Re$`C_3`$ appears and that the terms linear in Im$`C_3`$ have canceled. ### II.4 Calculation of $`p_{10}`$ To determine $`p_{10}`$ we use Eq. (9) and start again from $`\rho ^0=\rho _{\mathrm{inner}}^0`$ as initial condition in Eq. (13), but now have to determine $`s_{12}\|\rho ^1\|e_2`$. Replacing $`\{\|y_i\}`$ by $`\|e_2`$ and choosing $`\|x_{i_0}=\|s_{12},\|s_{23}`$ in Eq. (18), one obtains two inhomogeneous linear equations for $`s_{12}\|\rho ^1\|e_2`$ and $`s_{23}\|\rho ^1\|e_2`$. These equations do not depend on $`C_3`$, since $`R_\pm `$ and $`R_\pm ^{}`$ vanish on $`\|e_2`$ and since $`C_3`$ does not appear in the part of $`H_{\mathrm{cond}}`$ acting on the inner states. Therefore $`s_{12}\|\rho ^1\|e_2`$ and $`s_{23}\|\rho ^1\|e_2`$ are independent of $`C_3`$. By a simple calculation one obtains $`s_{12}\|\rho ^1\|e_2`$ and inserting this into Eq. (9) yields, to second order in $`\mathrm{\Omega }_2`$, $$p_{10}=\mathrm{\Omega }_2^2\frac{A_3\mathrm{\Omega }_3^2(A_3^2+4\mathrm{\Delta }_2^2)}{(A_3^2+2\mathrm{\Omega }_3^2)[(\mathrm{\Omega }_3^24\mathrm{\Delta }_2^2)^2+4\mathrm{\Delta }_2^2A_3^2]}.$$ (33) This is independent of $`C_3`$ and is the same as for two independent atoms, namely the transition rate for a single atom from a light to a dark period Wi . ### II.5 Calculation of $`p_{01}`$ To determine $`p_{01}`$ we use Eq. (10). One also needs $`s_{12}\|\rho ^1\|e_2`$, as seen from Eq. (9), but in this case one has to start from $`\rho ^0=\rho _{\mathrm{dark}}^0`$ as initial condition in Eq. (13). Therefore one obtains the same equations for $`s_{12}\|\rho ^1\|e_2`$ and $`s_{23}\|\rho ^1\|e_2`$ as before, except for the inhomogeneous part. One has independence of $`C_3`$ and easily solves for $`s_{12}\|\rho ^1\|e_2`$. For the remaining coherences needed in Eq. (10), i.e. those in Eq. (8), one obtains the same form as in Eq. (20), with the same $`𝐀`$, but now with $`𝐚_1=0`$ since the term containing $`\rho ^0`$ vanishes. Therefore these coherences vanish here and hence $`p_{02}=0`$. Physically this means that in our formulation of the problem there are no direct transitions from a dark period to a period of intensity 2. From Eq. (10) one now obtains, to second order in $`\mathrm{\Omega }_2`$, $$p_{01}=2\mathrm{\Omega }_2^2\frac{A_3\mathrm{\Omega }_3^2}{(\mathrm{\Omega }_3^24\mathrm{\Delta }_2^2)^2+4\mathrm{\Delta }_2^2A_3^2}$$ (34) This is independent of $`C_3`$ and is the same as for two independent atoms, namely twice the transition rate for a single atom from a dark to a light period. ### II.6 Calculation of $`p_{21}`$ The transition rate $`p_{21}`$ is obtained from Eq. (10) and the required coherences are again those appearing in Eq. (8) and (9), now with $`\rho ^0=\rho _{\mathrm{outer}}^0`$ as initial condition. For $`s_{12}\|\rho ^1\|e_2`$ and $`s_{23}\|\rho ^1\|e_2`$ one obtains the same two equations as before, except for the inhomogeneous part which now vanishes. Hence these two coherences vanish now and as a consequence $`p_{20}=0`$. Physically this means that in our formulation there are no direct transitions from a period of intensity 2 to a dark period. For the coherences in Eq. (19) one has the same equation as Eq. (20), with the same matrix $`𝐀`$ but with $`𝐚_1`$ replaced by $`𝐚_2=i\mathrm{\Omega }_2/\left\{\sqrt{2}\left(4\mathrm{\Omega }_{3}^{}{}_{}{}^{4}+4\mathrm{\Omega }_{3}^{}{}_{}{}^{2}A_{3}^{}{}_{}{}^{2}+A_{3}^{}{}_{}{}^{2}\mathrm{Re}C_{3}^{}{}_{}{}^{2}+2A_{3}^{}{}_{}{}^{3}\mathrm{Re}C_3+A_{3}^{}{}_{}{}^{2}\mathrm{Im}C_{3}^{}{}_{}{}^{2}+A_{3}^{}{}_{}{}^{4}\right)\right\}`$ (35) $`\times \left[\begin{array}{c}\mathrm{\Omega }_{3}^{}{}_{}{}^{4}+2\mathrm{\Omega }_{3}^{}{}_{}{}^{2}A_{3}^{}{}_{}{}^{2}+A_{3}^{}{}_{}{}^{2}\mathrm{Re}C_{3}^{}{}_{}{}^{2}+2A_{3}^{}{}_{}{}^{3}\mathrm{Re}C_3+A_{3}^{}{}_{}{}^{2}\mathrm{Im}C_{3}^{}{}_{}{}^{2}+A_{3}^{}{}_{}{}^{4}\\ \\ i\mathrm{\Omega }_3\sqrt{2}A_3\left(A_{3}^{}{}_{}{}^{2}+A_3\mathrm{Re}C_3+i\mathrm{Im}C_3A_3+\mathrm{\Omega }_{3}^{}{}_{}{}^{2}\right)\\ \\ \mathrm{\Omega }_{3}^{}{}_{}{}^{2}\left(A_3+\mathrm{Re}C_3+i\mathrm{Im}C_3\right)A_3\\ \\ i\mathrm{\Omega }_3A_3\left(\mathrm{\Omega }_{3}^{}{}_{}{}^{2}A_{3}^{}{}_{}{}^{2}A_3\mathrm{Re}C_3+i\mathrm{Im}C_3A_3\right)\\ \\ \mathrm{\Omega }_{3}^{}{}_{}{}^{2}\left(\mathrm{\Omega }_{3}^{}{}_{}{}^{2}+2A_{3}^{}{}_{}{}^{2}\right)/\sqrt{2}\\ \\ i\mathrm{\Omega }_{3}^{}{}_{}{}^{3}A_3\\ \\ 0\\ \\ \mathrm{\Omega }_{3}^{}{}_{}{}^{4}/\sqrt{2}\end{array}\right]`$ (44) Inserting the resulting coherences into Eq. (10) gives, to first order in Re$`C_3`$ and $`\mathrm{Im}C_3`$ and to second order in $`\mathrm{\Omega }_2`$, $`p_{21}=\mathrm{\Omega }_2^2\{{\displaystyle \frac{2A_3\mathrm{\Omega }_3^2(A_3^2+4\mathrm{\Delta }_2^2)}{(\mathrm{\Omega }_3^48\mathrm{\Delta }_2^2\mathrm{\Omega }_3^2+16\mathrm{\Delta }_2^4+4A_3^2\mathrm{\Delta }_2^2)(A_3^2+2\mathrm{\Omega }_3^2)}}`$ (45) $`+\mathrm{Re}C_3(r){\displaystyle \frac{4A_3^2\mathrm{\Omega }_3^2(A_3^4\mathrm{\Omega }_3^4+4A_3^2\mathrm{\Omega }_3^612A_3^2\mathrm{\Delta }_2^2\mathrm{\Omega }_3^464A_3^2\mathrm{\Delta }_2^64A_3^6\mathrm{\Delta }_2^232A_3^4\mathrm{\Delta }_2^464\mathrm{\Delta }_2^4\mathrm{\Omega }_3^4+16\mathrm{\Delta }_2^2\mathrm{\Omega }_3^6)}{(A_3^2+2\mathrm{\Omega }_3^2)^3(\mathrm{\Omega }_3^48\mathrm{\Delta }_2^2\mathrm{\Omega }_3^2+4A_3^2\mathrm{\Delta }_2^2+16\mathrm{\Delta }_2^4)^2}}\}`$ where again the terms containing Im$`C_3`$ have canceled. ### II.7 Discussion. If one computes the coherences in Eqs. (8) and (10) to second order in $`C_3`$ one obtains $`p_{12}`$ and $`p_{21}`$ to second order in $`C_3`$. The resulting expressions are not enlightening and therefore not given here, but they do depend on $`(\mathrm{Im}C_3)^2`$. Fig. 3 shows how small the second-order dipole-dipole contribution to $`p_{21}`$ is for the parameters of the simulations and for distances larger than half a wave length. For smaller distances the results are probably not applicable anyway, as discussed in Ref. BeHe5 . For $`\mathrm{\Delta }_2=0`$ the rates $`p_{12}`$ and $`p_{21}`$ simplify to $`p_{12}`$ $`=`$ $`\mathrm{\Omega }_2^2\left\{{\displaystyle \frac{A_3}{\mathrm{\Omega }_3^2}}+\mathrm{Re}C_3(r){\displaystyle \frac{2A_3^2}{\mathrm{\Omega }_3^2(A_3^2+2\mathrm{\Omega }_3^2)}}\right\}`$ (46) $`p_{21}=\mathrm{\Omega }_2^2\{{\displaystyle \frac{2A_3^3}{\mathrm{\Omega }_3^2(A_3^2+2\mathrm{\Omega }_3^2)}}+`$ (47) $`\mathrm{Re}C_3(r){\displaystyle \frac{4A_3^4(A_3^2+4\mathrm{\Omega }_3^2)}{\mathrm{\Omega }_3^2(A_3^2+2\mathrm{\Omega }_3^2)^3}}\}`$ and one sees that the coefficients of the Re$`C_3`$ term in Eqs. (46) and (47) are positive. For $`\mathrm{\Delta }_2=0`$, therefore, $`p_{12}`$ and $`p_{21}`$ vary with the atomic distance in phase with Re$`C_3`$. For $`\mathrm{\Delta }_20`$, however, the coefficients of Re$`C_3`$ in Eqs. (32) or (45) can become zero or negative. In the first case $`p_{12}`$ or $`p_{21}`$ become constant in $`r`$, while in the second case they vary opposite in phase to Re$`C_3`$. It will be shown in the next sections that this dependence of $`p_{12}`$ and $`p_{21}`$ on the detuning of the weak laser entails a corresponding behavior of the double jump rate and an opposite behavior of the mean durations $`T_1`$ and $`T_2`$. This opposite behavior of $`T_1`$ and $`T_2`$ is easy to understand since they are related to the inverse of the transition rates. ## III Double jumps: Comparison of simulations with theory A double jump is defined as a transition from a double-intensity period to dark period, or vice versa, within a prescribed time interval $`\mathrm{\Delta }T_{\mathrm{DJ}}`$. Now, to distinguish different periods in experiments and in simulations one has to use an average photon intensity, obtained e.g. by means of averaging over a time window. This window has to be large enough to contain enough emissions, but must not be too large in order not to overlook too many short periods. Our simulations employ a procedure similar to that in Ref. BeHe5 and use a moving window window of fixed width, denoted by $`\mathrm{\Delta }T_\mathrm{w}`$. The time interval $`\mathrm{\Delta }T_{\mathrm{DJ}}`$ should be larger than $`\mathrm{\Delta }T_\mathrm{w}`$. We consider the fluorescence periods as a telegraph process with three steps and use the $`p_{ij}`$ of the last section as transition rates. At first the influence of the averaging window $`T_\mathrm{w}`$ will be neglected. The rate of downward double jumps is obtained as follows. For $`i`$ = 0, 1, 2, let $`n_i`$ be the mean number of periods of intensity $`i`$ per unit time. For a long path of length $`T`$ the total number of periods of intensity $`i`$ is then $`N_i(T)=n_iT`$. At the end of each period of intensity 2 there begins a period of intensity 1, and the probability for this period of intensity 1 to be shorter than $`\mathrm{\Delta }T_{\mathrm{DJ}}`$ is given by $$1\mathrm{exp}\{(p_{10}+p_{12})\mathrm{\Delta }T_{\mathrm{DJ}}\}.$$ At the end of a period of intensity 1 the branching ratio for a transition to a period of intensity 0 is $`p_{10}/(p_{10}+p_{12})`$. Thus during time $`T`$ the total number of such downward double jumps, denoted by $`N_{\mathrm{DJ}}^{20}(T)`$, is $$N_{\mathrm{DJ}}^{20}(T)=N_2(T)\frac{p_{10}}{(p_{10}+p_{12})}\left\{1\mathrm{exp}\{(p_{10}+p_{12})\mathrm{\Delta }T_{\mathrm{DJ}}\}\right\}$$ and therefore the rate, $`n_{\mathrm{DJ}}^{20}`$, of downward double jumps within $`\mathrm{\Delta }T_{\mathrm{DJ}}`$ is $$n_{\mathrm{DJ}}^{20}=n_2\frac{p_{10}}{(p_{10}+p_{12})}\left\{1\mathrm{exp}\{(p_{10}+p_{12})\mathrm{\Delta }T_{\mathrm{DJ}}\}\right\}.$$ (48) In a similar way one finds that the rate, $`n_{\mathrm{DJ}}^{02}`$, of upward double jumps within $`\mathrm{\Delta }T_{\mathrm{DJ}}`$ is $$n_{\mathrm{DJ}}^{02}=n_0\frac{p_{12}}{(p_{10}+p_{12})}\left\{1\mathrm{exp}\{(p_{10}+p_{12})\mathrm{\Delta }T_{\mathrm{DJ}}\}\right\}.$$ (49) It remains to determine $`n_0`$ and $`n_2`$. Since a period of intensity 1 ends with a transition to a period of either intensity 0 or intensity 2 one has, with the respective branching ratios, $$n_0=\frac{p_{10}}{p_{10}+p_{12}}n_1$$ (50) $$n_2=\frac{p_{12}}{p_{10}+p_{12}}n_1.$$ (51) If one denotes by $`T_i`$ the mean durations of a period of intensity $`i`$, one has $$\underset{i=0}{\overset{2}{}}n_iT_i=1.$$ (52) Moreover, one has $$T_0=1/p_{01},T_1=1/(p_{10}+p_{12}),T_2=1/p_{21}$$ (53) and this then gives $$n_0=\frac{p_{01}p_{21}}{p_{01}p_{21}+p_{21}p_{10}+p_{01}p_{12}}p_{10}$$ (54) $$n_2=\frac{p_{01}p_{21}}{p_{01}p_{21}+p_{21}p_{10}+p_{01}p_{12}}p_{21}.$$ (55) From this, together with Eqs. (48) and (49), one sees immediately that the rates of upward and downward double jumps are equal, $$n_{\mathrm{DJ}}^{02}=n_{\mathrm{DJ}}^{20}.$$ (56) This fact was also observed in the simulations. The combined number of double jumps therefore equals $`n_{\mathrm{DJ}}n_{\mathrm{DJ}}^{02}+n_{\mathrm{DJ}}^{20}`$ (57) $`=2{\displaystyle \frac{p_{01}p_{10}p_{12}p_{21}}{(p_{01}p_{21}+p_{21}p_{10}+p_{01}p_{12})(p_{01}+p_{12})}}`$ $`\times \left\{1\mathrm{exp}\{(p_{10}+p_{12})\mathrm{\Delta }T_{\mathrm{DJ}}\}\right\}.`$ For $`\mathrm{\Delta }T_{\mathrm{DJ}}T_1`$ and by expanding the exponential, this gives for the combined double jump rate, without correction for the averaging window, $$n_{\mathrm{DJ}}=2\frac{p_{01}p_{10}p_{12}p_{21}}{p_{01}p_{21}+p_{21}p_{10}+p_{01}p_{12}}\mathrm{\Delta }T_{\mathrm{DJ}}.$$ (58) Fig. 4 shows a comparison of this result with data from the simulations. Except for atomic distances less than about three quarters of the wave length of the strong transition the agreement appears as quite reasonable, and the disagreement for small distances is not unexpected since there the intensities start to decrease and a description by a telegraph process may be no longer a good approximation, as pointed out in Ref. BeHe5 . But one observes that the theoretical result is systematically above the simulated curve. This seeming disagreement, however, is easily explained and can be taken care of as follows. ### III.1 Corrections for averaging window We recall that the simulated data were obtained by averaging the numerical photon emission times with a moving window of length $`\mathrm{\Delta }T_\mathrm{w}`$. Then, roughly, periods which are shorter than about two thirds of the window length are overlooked, and therefore the number of recorded (or observed) periods of type 2, which enters Eq. (48), is smaller than that given by Eq. (55). The recorded or observed number is denoted by $`n_{2,\mathrm{cor}}`$. It is approximately given by $$n_{2,\mathrm{cor}}=n_2\mathrm{exp}\{p_{21}\frac{2}{3}\mathrm{\Delta }T_\mathrm{w}\},$$ (59) and this expression should be inserted into Eq. (48) for $`n_2`$. In this way one obtains the corrected theoretical curve in Fig. 5. The curve changes very little if instead of two thirds one takes 60% or 70% of $`\mathrm{\Delta }T_\mathrm{w}`$. It is seen that the agreement with the simulated data is much improved for distances greater than three quarters of a wave length of the strong transition. It still appears, however, that the oscillation amplitudes of the theoretical curve are somewhat larger than those of the simulated curve. This is again understandable as an effect of the averaging procedure. In the simulations it was noticed numerically that the $`r`$ dependence of the double jump rate depended somewhat on the length of the averaging window $`T_\mathrm{w}`$ and distinct features tended to be somewhat washed out for larger $`\mathrm{\Delta }T_\mathrm{w}`$, in particular the oscillation amplitudes of the simulated data decreased with the length of the averaging window. A larger $`\mathrm{\Delta }T_\mathrm{w}`$ gave a smoother intensity curve, but made the determination of the transition times between different periods more difficult, while a shorter averaging window introduced more noise. We found the use of $`\mathrm{\Delta }T_\mathrm{w}=114A_3^1`$ to be a good compromise. If it were possible to choose smaller averaging window the amplitudes should increase, as predicted by the theory. ### III.2 Detuning One can explicitly insert the expressions for $`p_{ij}`$ of the last section into Eq. (58), but the result becomes unwieldy. One can show that in an expansion of Eq. (58) with respect to Re$`C_3`$ to first order the coefficient of Re$`C_3`$ is positive for zero detuning. This implies that the double jump rate is in phase with Re$`C_3(r)`$ for the atomic distances under consideration and for zero detuning. For increasing detuning the double jump rate can become constant in $`r`$ and then change its oscillatory behavior to that of $``$Re$`C_3`$. An example for the latter is shown in Fig. 6. ## IV Duration of fluorescence periods: Effect of averaging window The mean durations, $`T_0`$, $`T_1`$, and $`T_2`$, of the three periods were investigated for cooperative effects in Ref. BeHe5 by simulations with averaging windows at discrete times. Here we have performed similar simulations with a moving averaging window. It turns out that both the present and the previous simulation for $`T_i`$ are about 15% higher than those predicted by Eq. (53), using the expressions for $`p_{ij}`$ of Section II and without correcting for the use of the averaging window due to which short periods are not recorded. We will now show how this can be taken into account in the theory. As in Section III we consider a three-step telegraph process with periods of type 0, 1, and 2, whose mean durations are denoted by $`T_0`$, $`T_1`$, and $`T_2`$, respectively. We assume that periods of length $`\mathrm{\Delta }\tau `$ or less are not recorded. Fig. 7 shows periods of type 1 which are interrupted by a short period of type 0 and 2, respectively. If the respective short periods are not recorded, then the two periods of type 1 in the left part of the figure are recorded as a single longer period, and similarly for the right part of the figure. This leads to an apparent decrease of shorter periods of type 1 and to a corresponding increase of longer periods. To make this quantitative we put $`\lambda _i1/T_i`$ and denote the number per unit time of periods of type $`i`$, whose duration is less than $`\mathrm{\Delta }\tau `$, by $`n_i^{\mathrm{\Delta }\tau }`$, i.e. $$n_i^{\mathrm{\Delta }\tau }=n_i\left\{1\mathrm{exp}\{\lambda _i\mathrm{\Delta }\tau \}\right\}.$$ (60) Per unit time, one has $`n_0^{\mathrm{\Delta }\tau }`$ occurrences of the situation in the left part of Fig. 7 and $`n_2^{\mathrm{\Delta }\tau }`$ occurrences of that in the right part. The probability for one of the periods of type 1 in the left or right part of Fig. 7 to have a length lying in the time interval $`(t_1,t_1+dt_1)`$ is $`2\lambda _1\mathrm{exp}\{\lambda _1t_1\}dt_1`$, where the factor of 2 comes from the two possible situations. Therefore, the recorded number, per unit time, of periods of type 1 with duration in $`(t_1,t_1+dt_1)`$ is changed (decreased) by $$2(n_0^{\mathrm{\Delta }\tau }+n_2^{\mathrm{\Delta }\tau })\lambda _1\mathrm{exp}\{\lambda _1t_1\}dt_1.$$ (61) Similarly, the apparent increase of the number, per unit time, of periods of type 1 with duration in $`(t_1,t_1+dt_1)`$ is, by Fig. 7, $`(n_0^{\mathrm{\Delta }\tau }+n_2^{\mathrm{\Delta }\tau }){\displaystyle \underset{t_1t_1^{}+t_1^{\prime \prime }t_1+dt_1}{}𝑑t_1^{}𝑑t_1^{\prime \prime }}`$ (62) $`\times \lambda _1\mathrm{exp}\{\lambda _1t_1^{}\}\lambda _1\mathrm{exp}\{\lambda _1t_1^{\prime \prime }\}`$ $`=(n_0^{\mathrm{\Delta }\tau }+n_2^{\mathrm{\Delta }\tau })\lambda _1^2t_1\mathrm{exp}\{\lambda _1t_1\}dt_1.`$ Denoting by $`\nu _{1\mathrm{r}\mathrm{e}\mathrm{c}}(t_1)dt_1`$ the actually recorded number, per unit time, of periods of type 1 with duration in $`(t_1,t_1+dt_1)`$ one obtains from the two previous expressions $`\nu _{1\mathrm{r}\mathrm{e}\mathrm{c}}(t_1)dt_1=n_1\lambda _1\mathrm{exp}\{\lambda _1t_1\}dt_1`$ (63) $`+(n_0^{\mathrm{\Delta }\tau }+n_2^{\mathrm{\Delta }\tau })(\lambda _1^2t_12\lambda _1)\mathrm{exp}\{\lambda _1t_1\}dt_1.`$ The average duration of the recorded periods of type 1 will be denoted by $`T_{1,\mathrm{cor}}`$, and it is given by $$T_{1,\mathrm{cor}}=_{\mathrm{\Delta }\tau }^{\mathrm{}}𝑑t_1t_1\nu _{1\mathrm{r}\mathrm{e}\mathrm{c}}(t_1)/_{\mathrm{\Delta }\tau }^{\mathrm{}}𝑑t_1\nu _{1\mathrm{r}\mathrm{e}\mathrm{c}}(t_1).$$ (64) Using Eq. (63) for $`\nu _{1\mathrm{r}\mathrm{e}\mathrm{c}}(t_1)`$ one obtains, after an elementary calculation and for $`\mathrm{\Delta }\tau `$ satisfying $`\mathrm{\Delta }\tau /T_11`$, $$T_{1,\mathrm{cor}}=\frac{1}{p_{10}+p_{12}}+\mathrm{\Delta }\tau \left\{1+\frac{p_{01}p_{10}+p_{12}p_{21}}{(p_{10}+p_{12})^2}\right\}.$$ (65) The first term is the ideal theoretical value, $`T_1`$, and the remainder is the correction due to non-recorded short periods. In a similar way one obtains $$T_{0,\mathrm{cor}}=\frac{1}{p_{01}}+\mathrm{\Delta }\tau \left\{1+\frac{p_{10}}{p_{01}}\right\}$$ (66) $$T_{2,\mathrm{cor}}=\frac{1}{p_{21}}+\mathrm{\Delta }\tau \left\{1+\frac{p_{12}}{p_{21}}\right\}$$ (67) where again the respective first terms are the ideal values, $`T_0`$ and $`T_2`$. To compare this with simulated data, obtained with a moving averaging window of length $`\mathrm{\Delta }T_\mathrm{w}=247A_3^1`$, we have taken $`\mathrm{\Delta }\tau =\frac{2}{3}\mathrm{\Delta }T_\mathrm{w}`$, as in the previous section, and have plotted the results together with the simulated data in Fig. 8. The agreement is very good. Quite generally, for zero detuning the oscillations of $`T_1`$ and $`T_2`$ are opposite in phase to those of Re$`C_3(r)`$, as already noted at the end of Section II. As in the case of the double jump rate, $`T_1`$ and $`T_2`$ can become constant in $`r`$ for particular values of the detuning (different for $`T_1`$ and $`T_2`$ ), and then change to a behavior in phase with Re$`C_3(r)`$. The above approach of taking the averaging window into account works for the following reason. For a single atom with macroscopic dark periods it is known that the emission of photons is describable, to high accuracy, by an underlying two-step telegraph process. For two independent atoms with macroscopic dark periods the emissions are therefore described by an underlying three-step telegraph process. For two atoms interacting by a weak dipole-dipole interaction the actual emission process of photons should therefore still have, at least approximately, an underlying three-step telegraph process. What we have done above is replacing the actual emission process by this underlying three-step telegraph process and then incorporating the averaging window by taking into account the influence of the overlooked short periods on the statistics. ## V Discussion of results We have investigated cooperative effects in the fluorescence of two dipole-dipole interacting atoms in a $`V`$ configuration. One of the excited states of the $`V`$ configuration is assumed to be metastable, i.e. with a weak transition to the ground state. When driven by two lasers, a single such configuration exhibits macroscopic dark periods and periods of fixed intensity, like a two-step telegraph process. A system of two such atoms exhibits three fluorescence types, i.e. dark periods and periods of single and double intensity, like a three-step telegraph process. For large atomic distances, when the dipole-dipole interaction is negligible, the total fluorescence just consists of the sum of the individual atomic contributions. We have shown that for smaller atomic distances the fluorescence modified by the dipole-dipole interaction which depends on the atomic distance $`r`$. In particular we have, to our knowledge for the first time, explicitly demonstrated cooperative effects in the rate of double jumps from a period of double intensity to a dark period or vice versa, both analytically and by simulations. By means of an analytical theory we have obtained the r-dependent transition rates , $`p_{ij}`$, between the three intensity periods. These were then used to calculate the rate of double jumps and in the mean period durations $`T_0,T_1,`$ and $`T_2`$. When comparing with the simulations it turned out that one had to take into account the averaging window used for obtaining an intensity curve from the individual photon emissions. With this the agreement between simulation and analytic theory became excellent. For zero laser detuning, for which the simulations were performed, the double jump rates are in phase with and $`T_1`$ and $`T_2`$ opposite in phase to Re$`C_3(r)`$. The theoretical expressions, however, allow general detuning, $`\mathrm{\Delta }_2`$, of the laser which drives the weak transition. It has been shown that for a particular $`\mathrm{\Delta }_2`$, which depends on the other parameters, the double jump rate becomes constant and, for larger $`\mathrm{\Delta }_2`$, varies opposite in phase to Re$`C_3(r)`$. A similar change of characteristic behavior also occurs for $`T_1`$ and $`T_2`$, for different values of $`\mathrm{\Delta }_2`$ though. The amplitude of the oscillations with the atomic distance remain in the same region of magnitude as for zero detuning. As pointed out in Ref. BeHe5 , a dependence of the oscillations on Re$`C_3(r)`$ is not unexpected since Re$`C_3(r)`$ affects the decay rates of the excited Dicke states of the combined system. But an intuitive argument why the above change of behavior occurs for increased detuning is at present not apparent. We have pointed out in Section III that there is another statistical property of the fluorescence which can serve as an indicator of the influence of the dipole-dipole interaction and which is probably not too difficult to determine experimentally. This quantity is the rate with which fluorescence periods of definite type occur, in particular the rate of periods with double intensity. Our theoretical results show that this rate behaves similar to the double jump rate, as regards the variation with the atomic distance, and an example is shown in Fig. 9. This quantity is probably much easier to measure than the double jump rate or the mean duration $`T_2`$. Our theoretical approach can be carried over to other level configurations and to more than two atoms. For given parameters the evaluation should be not too difficult. If, however, one is interested in closed algebraic expressions the effort will increase considerably with the number of atoms. In particular, it would be interesting to apply our approach to the situation of the experiment of Ref. Sauter with its different level configuration and its three ions in the trap. ## Appendix A Dipole-dipole interaction in the Bloch equations The dipole-dipole interaction enters the Bloch equations through $`r`$-dependent complex coupling constants (cf. Ref. BeHe5 ) $`C_j={\displaystyle \frac{3A_j}{2}}\mathrm{e}^{ik_{j1}r}[{\displaystyle \frac{1}{ik_{j1}r}}(1\mathrm{cos}^2\vartheta _j)`$ (68) $`+({\displaystyle \frac{1}{(k_{j1}r)^2}}{\displaystyle \frac{1}{i(k_{j1}r)^3}})(13\mathrm{cos}^2\vartheta _j)].`$ Here $`\vartheta _j`$ is the angle between the transition dipole moment $`𝐃_{1j}`$ and the line connecting the atoms and $`k_{j1}=2\pi /\lambda _{j1}`$, where $`\lambda _{j1}`$ is the wavelength of the j-1 transition for an atom. For $`A_20`$ one has $`C_20`$. Thus one can neglect the dipole interaction when one atom is in state $`\|2`$. The dependence of $`C_3`$ on $`r`$ is maximal for $`\vartheta _3=\pi /2`$ and the corresponding $`C_3`$ is plotted in Fig. 10. For atomic distances greater than about three quarters of a wave length of the strong transition, $`\|C_3\|`$ is less than $`0.2A_3`$, but for smaller distances Re$`C_3`$ approaches $`A_3`$ and Im$`C_3`$ diverges. The reset operation $``$ and $`H_{\mathrm{cond}}`$ are given by the same expressions as in Ref. BeHe5 , except for the detuning. One has $$(\rho )=\left(A_3+\mathrm{Re}C_3\right)R_+\rho R_+^{}+\left(A_3\mathrm{Re}C_3\right)R_{}\rho R_{}^{}$$ (69) where $`R_+`$ $`=`$ $`\left(S_{13}^{}+S_{23}^{}\right)/\sqrt{2}`$ $`=`$ $`\|gs_{13}\|+\|s_{13}e_3\|+\left(\|s_{12}s_{23}\|\|a_{12}a_{23}\|\right)/\sqrt{2},`$ $`R_{}`$ $`=`$ $`\left(S_{13}^{}S_{23}^{}\right)/\sqrt{2}`$ (70) $`=`$ $`\|ga_{13}\|+\|a_{13}e_3\|+\left(\|s_{12}a_{23}\|+\|a_{12}s_{23}\|\right)/\sqrt{2}.`$ The summands in Eq. (6) are given by $`H_{\mathrm{cond}}^0`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{2i}}\left[A_3\left(\|s_{23}s_{23}\|+\|a_{23}a_{23}\|\right)+(A_3+C_3)\|s_{13}s_{13}\|+(A_3C_3)\|a_{13}a_{13}\|+2A_3\|e_3e_3\|\right]`$ (71) $`+{\displaystyle \frac{\mathrm{}}{2}}[\sqrt{2}\mathrm{\Omega }_3(\|gs_{13}\|+\|s_{13}e_3\|)+\mathrm{\Omega }_3(\|s_{12}s_{23}\|\|a_{12}a_{23}\|)+\mathrm{H}.\mathrm{c}.]`$ $`\mathrm{}\mathrm{\Delta }_2\left[2\|e_2e_2\|+\|s_{12}s_{12}\|+\|a_{12}a_{12}\|+\|s_{23}s_{23}\|+\|a_{23}a_{23}\|\right]`$ $`H_{\mathrm{cond}}^1(\mathrm{\Omega }_2)={\displaystyle \frac{\mathrm{}}{2}}[\sqrt{2}\mathrm{\Omega }_2(\|gs_{12}\|+\|s_{12}e_2\|)+\mathrm{\Omega }_2(\|s_{13}s_{23}\|+\|a_{13}a_{23}\|)+\mathrm{H}.\mathrm{c}.]`$ (72) From Eq. (71) one sees that Re$`C_3`$ changes the spontaneous decay rates and that Im$`C_3`$ leads to level shifts. Therefore, for small $`r`$, the decay rate of $`\|a_{13}`$ approaches 0 in this case and the large level shifts cause a decrease of fluorescence associated with the levels $`\|s_{13}`$ and $`\|a_{13}`$. ## Appendix B calculation of $`\rho (t_0+\mathrm{\Delta }t)`$ to first order in $`\mathrm{\Omega }_2`$ We write the Bloch equations of Eq. (5) in the form $$\dot{\rho }=\rho $$ (73) where the Liouvillean $`(A_3,\mathrm{\Omega }_3,\mathrm{\Delta }_2,C_3,\mathrm{\Omega }_2)`$, a super-operator, can be read off from Eqs. (5) and (69) - (72). One can decompose $``$ as $$=_0+_{\mathrm{\Omega }_2}$$ (74) where $`_0=(A_3,\mathrm{\Omega }_3,\mathrm{\Delta }_2,C_3,0)`$ and $`_{\mathrm{\Omega }_2}\rho =\mathrm{i}[H_{\mathrm{cond}}^1(\mathrm{\Omega }_2),\rho ]/\mathrm{}`$. We note that $`H_{\mathrm{cond}}^1(\mathrm{\Omega }_2)`$ is Hermitian and that $`_0`$ can be considered as a Liouvillean of Bloch equations. Choosing an initial density matrix $`\rho (t_0)`$ lying in one of the subspaces in Eqs. (2) - (4) one obtains, to first order in $`\mathrm{\Omega }_2`$, $`\rho (t_0+\mathrm{\Delta }t)`$ $`=`$ $`e^{\mathrm{\Delta }t}\rho (t_0)`$ (75) $`=`$ $`e^{_0\mathrm{\Delta }t}\rho (t_0)`$ $`+{\displaystyle _0^{\mathrm{\Delta }t}}𝑑\tau e^{_0(\mathrm{\Delta }t\tau )}_{\mathrm{\Omega }_2}e^{_0\tau }\rho (t_0),`$ just as with usual quantum mechanical perturbation theory in the interaction picture. Now we use the fact that $`_0`$, as a Liouvillean of Bloch equations, has an eigenvalue $`0`$ (corresponding to steady states) and eigenvalues with negative real parts of the order of $`\mathrm{\Omega }_3`$ and $`A_3`$. Therefore, if $`\mathrm{\Delta }t`$ satisfies Eq. (7), the first term on the right-hand side of Eq. (75) gives one of the equilibrium states, $`\rho ^0`$, of $`_0`$ given in Eqs. (14) - (16), to high accuracy, while the term $`e^{_0\tau }\rho (t_0)`$ under the integrand also rapidly approaches $`\rho ^0`$. After a change of integration variable one therefore has to first order in $`\mathrm{\Omega }_2`$ $$\rho (t_0+\mathrm{\Delta }t)=\rho ^0+_0^{\mathrm{\Delta }t}𝑑\tau e^{_0\tau }_{\mathrm{\Omega }_2}\rho ^0.$$ (76) It can be shown that $`_{\mathrm{\Omega }_2}\rho ^0`$ has no components in the zero-eigenvalue subspace of $`_0`$ zero . Therefore, the integrand in Eq. (76) is rapidly damped, and since $`\mathrm{\Delta }t\mathrm{\Omega }_3^1,A_3^1`$, the upper integration limit can be extended to infinity. Hence we can write, to first order in $`\mathrm{\Omega }_2`$, $$\rho (t_0+\mathrm{\Delta }t)=\rho ^0+_0^{\mathrm{}}𝑑\tau e^{_0\tau }_{\mathrm{\Omega }_2}\rho ^0.$$ (77) Thus, if $`\mathrm{\Delta }t`$ satisfies Eq. (7) then, to first order in $`\mathrm{\Omega }_2`$, $`\rho (t_0+\mathrm{\Delta }t)`$ is independent of $`\mathrm{\Delta }t`$, and one has $$\rho (t_0+\mathrm{\Delta }t)=\rho ^0+(ϵ_0)^1_{\mathrm{\Omega }_2}\rho ^0$$ (78) to first order in $`\mathrm{\Omega }_2`$, where the limit $`ϵ+0`$ is understood. Multiplying this by $`ϵ`$ gives $$\rho (t_0+\mathrm{\Delta }t)=_{\mathrm{\Omega }_2}(ϵ_0)^1_{\mathrm{\Omega }_2}\rho ^0=𝒪(\mathrm{\Omega }_2^2)$$ (79) which is Eq. (11). That the transition rates are independent of the particular choice of $`\mathrm{\Delta }t`$ follows from Eqs. (78) and (8).
warning/0002/cond-mat0002328.html	ar5iv	text	# Anderson localization in bipartite lattices ## I Introduction An interesting and still debated issue in the physics of the Anderson’s localization concerns the existence of delocalized states in dimensions $`d2`$, the conditions under which they appear, and their properties. This problem, which is, for instance, of relevance in the theory of the integer quantum Hall effect, got recently a renewed interest after evidences of a metal–insulator transition in two dimensions have been discovered. One of the cases in which localization does not occur in any dimension is at the band center energy of a tight binding model on a bipartite lattice, when both the regular hopping and the disorder only couple one sublattice to the other, i.e. in the so called two sublattice model. Although this is not a common physical situation, its consequences are surprising, and seem to escape any quasi-classical interpretation, which, on the contrary, provides simple physical explanations of other delocalization mechanisms . Already in 1976, Theodorou and Cohen realized that a one dimensional tight binding model with nearest neighbor random hopping has a single delocalized state at the band center (see also Ref.). Afterwards, Wegner and Oppermann and Wegner showed that a delocalized state indeed exists under the above conditions in any dimension, within a large–$`n`$ expansion, being $`n`$ the number of orbitals per site. Later on, Wegner and Gade argued that these models correspond to a particular class of non-linear $`\sigma `$-models for matrices in the zero replica limit. They were able to show that the quantum corrections to the $`\beta `$-function which controls the scaling behavior of the conductance vanish at the band center at all orders in the disorder strength. thus implying a metallic behavior at this value of the chemical potential. Moreover, they showed that, contrary to the standard case, the $`\beta `$-function of the density of states is finite. These results were based upon the non-linear $`\sigma `$-model derived by Gade by means of a boson-replica trick method, in a particular two sublattice Hamiltonian with broken time reversal invariance. More recently, the model without time reversal invariance has got a renewed interest for its implications in different physical contexts, for instance models with non-Hermitean stochastic operators, or random flux models in two dimensions (see e.g. Refs.). In this paper, we present an analysis of a generic disordered tight-binding Hamiltonian on a generic bipartite lattice. The starting model is therefore of quite general validity, also describing systems with time reversal invariance, and reduces in particular cases to the model discussed by Gade, or, in the honeycomb lattice, to models of Dirac fermions, or, finally, to random flux models. By means of a fermionic-replica trick method, we derive the generic non-linear $`\sigma `$-model describing the diffusive modes, which we analyse by the Renormalization Group (RG). Needless to say, the effective model belongs to the same class of non-linear $`\sigma `$-models identified by Wegner and Gade, demonstrating once more the universality of this description in the theory of Anderson localization. Since the work is quite technical, we prefer to give in the following section a short summary of the main results. ### A Summary of the main results In this section we shortly present the main results, with particular emphasis to the connections and differences with the standard theory of the Anderson’s localization. We consider a generic bipartite lattice and work with a unit cell which contains two sites from opposite sublattices. The Pauli matrices $`\sigma `$’s act on the two components of the wavefunction, corresponding to the two sites within each unit cell. In this lattice, we study a disordered tight-binding Hamiltonian which has the peculiar property of involving only both Pauli matrices $`\sigma _1`$ and $`\sigma _2`$. In other words, $`H`$ satisfies the conditions $$\{H,\sigma _3\}=0,\{H,\sigma _1\}0,\{H,\sigma _2\}0,$$ (1) where $`\{\mathrm{},\mathrm{}\}`$ indicates the anticommutator. By means of a path-integral approach within a fermionic replica trick method, we find that the low-energy diffusive modes at the band center, $`E=0`$, can be represented by the non-linear $`\sigma `$-model $`S[U]`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{16}}{\displaystyle 𝑑RTr\left[\stackrel{}{}U(R)^2\stackrel{}{}U(R)^2\right]}`$ (2) $``$ $`{\displaystyle \frac{2\pi }{32}}{\displaystyle \frac{\mathrm{\Pi }}{2}}{\displaystyle 𝑑R\left\{Tr\left[U(R)^2\stackrel{}{}U(R)^2\right]\right\}^2},`$ (3) where $`U(R)`$ is unitary and belongs to the coset space U($`4m`$)/Sp($`2m`$), being $`m`$ the number of replicas. At finite energy $`E0`$, the symmetry of $`U(R)`$ gets reduced to Sp($`2m`$)/Sp($`m`$)$`\times `$ Sp($`m`$), as in the standard case. The enlarged symmetry is accompanied by new diffusive modes which appear in the retarded-retarded and advanced-advanced channels, which are instead massive in the standard case. In (3), $`\sigma _{xx}`$ is the Kubo conductivity (in units of $`e^2/\mathrm{}`$) in the Drude approximation. We find that the new coupling constant, $`\mathrm{\Pi }`$, is proportional to the fluctuations of the staggered density of states, i.e. to the following correlation function $$\frac{1}{V}\underset{RR^{}}{}\mathrm{e}^{iq(RR^{})}(\overline{\rho _s(E,R)\rho _s(E,R^{})},$$ (4) where the bar indicates the impurity average, and $`\rho _s(E,R)`$ is the staggered density of states at energy $`E`$, $`\rho _s(E,R)={\displaystyle \underset{n}{}}\varphi _n(R)^{}\sigma _3\varphi _n(R)\delta (Eϵ_n),`$ being $`\varphi _n(R)`$ the two-component eigenfunction of energy $`ϵ_n`$. The structure of the above action was derived by Gade and Wegner for a particular Hamiltonian. Here, we derive it for a generic bipartite lattice and random hopping. Moreover, we provide a simple physical interpretation of $`\mathrm{\Pi }`$. Going back to (3), Gade and Wegner gave a beautiful proof, based just on symmetry considerations, that the quantum corrections to the $`\beta `$-function of $`\sigma _{xx}`$ vanish in the zero replica limit. In Appendix E, we outline how their proof works in our case, which is sligthly, but not qualitatively, different from the U(N)/SO(N) case they have considered. Essentially, one can show that the action (3) posseses an invariant coupling $`\sigma _{xx}+m\mathrm{\Pi }`$, which, in the $`m0`$ limit, implies that $`\sigma _{xx}`$ is not renormalized, apart from its bare dimensions. On the contrary, both the density of states and $`\mathrm{\Pi }`$ have non vanishing $`\beta `$-functions. In $`d=1`$, the system flows to strong coupling, hence we can not access the asymptotic infrared behavior. Nevertheless, the starting flow of the running variables indicates that the density of states diverges. In $`d=2,3`$, the system flows to weak coupling, hence we can safely assume that the infrared behavior is captured by the RG equations. Indeed, in two dimensions, the density of states diverges at $`E=0`$, while, in $`d=3`$, it saturates to a finite value, although exponentially increased in $`1/\sigma _{xx}`$. Moreover, $`\mathrm{\Pi }`$ has an anomalous behavior in $`d=2`$, where it is predicted to diverge logarithmically. We explicitly estimate how these quantities behave as $`E0`$, by means of a two-cutoff scaling approach, as discussed by Gade. We have also analysed various symmetry breaking terms. The simplest ones are those which spoil the particular symmetry Eq.(1) of the model at $`E=0`$, i.e. an on-site disorder or a same-sublattice regular hopping. These perturbations bring the symmetry of $`U(R)`$ down to Sp($`2m`$)/Sp($`m`$)$`\times `$ Sp($`m`$), as in the standard localization problem. However, by evaluating the anomalous dimensions of these terms, we can estimate the cross-over lengths above which the symmetry reduction is effective. While in $`d=1,2`$ these terms always lead to a localized behavior also at the band center, in $`d=3`$ the vicinity to the band center leads to an increase of the window in which delocalized states exist. Finally, if the impurity potential breaks time-reversal symmetry, the matrix field $`U(R)`$ is shown to belong to the coset space U$`(2m)`$, which indeed agrees with the analysis of Gade. The paper is organized as follows. In section II, we introduce the Hamiltonian. In section III, we derive the path-integral representation of the model, by using Grassmann variables within the replica trick method, and, in section IV, we study the symmetry properties of the action. In section V, we evaluate the saddle point of the action, while in sections VI, VII, VIII we derive the effective non-linear $`\sigma `$-model describing the long-wavelength fluctuations around the saddle point. The Renormalization Group analysis is presented in section IX, and the behavior in the presence of on-site disorder, of a same-sublattice regular hopping or with broken time reversal invariance is studied in sections X, XI, and XII, respectively. Finally, section XIV is devoted to a discussion of the results. We have also included several appendices containing more technical parts. ## II The Model We consider a tight binding Hamiltonian on a bipartite lattice, of the form $$H=\underset{RA}{}\underset{R^{}B}{}h_{RR^{}}\left(c_R^{}c_R^{}^{}+c_R^{}^{}c_R^{}\right),$$ (5) where $`A`$ and $`B`$ label the two sublattices and the hopping matrix elements $`h_{RR^{}}`$ are randomly distributed. We take a unit cell which includes two sites from different sublattices. In some cases, like the honeycomb lattice, this is indeed the primitive unit cell. In other cases, like the square lattice, it is not. In this representation, the Hamiltonian can be written as $$H=\underset{RR^{}}{}h_{RR^{}}^{12}(c_{1R}^{}c_{2R^{}}^{}+H.c.)$$ (6) where $`1`$ and $`2`$ label now the two sites in the unit cell, while $`R`$ and $`R^{}`$ refer to the unit cells, and $`h_{RR^{}}^{12}=h_{R^{}R}^{21}`$. By introducing the two component operators $`c_R^{}=\left(\begin{array}{c}c_{1R}^{}\\ c_{2R}^{}\end{array}\right),`$ we can also write the Hamiltonian as $`H`$ $`=`$ $`{\displaystyle \underset{R,R^{}}{}}c_R^{}H_{RR^{}}c_R^{}^{}`$ (7) $`=`$ $`{\displaystyle \underset{R,R^{}}{}}{\displaystyle \frac{1}{2}}\left(h_{RR^{}}^{12}+h_{RR^{}}^{21}\right)c_R^{}\sigma _1c_R^{}^{}+{\displaystyle \frac{i}{2}}\left(h_{RR^{}}^{12}h_{RR^{}}^{21}\right)c_R^{}\sigma _2c_R^{}^{},`$ (8) where the $`\sigma _i`$’s ($`i=1,2,3`$) are the Pauli matrices. We notice that, quite generally, the Hamiltonian involves both $`\sigma _1`$ and $`\sigma _2`$, but neither $`\sigma _3`$ nor $`\sigma _0`$, so that it satisfies the conditions in Eq. (1). We can write $`h_{RR^{}}^{12}=t_{RR^{}}^{12}+\tau _{RR^{}}^{12}`$, where $`t_{RR^{}}^{12}`$ are the average values, which represent the regular (translationally invariant) hopping matrix elements, while $`\tau _{RR^{}}^{12}`$ are random variables with zero average, which we assume to be gaussian distributed with width $`\left(\tau _{RR^{}}^{12}\right)^2=u^2\left(t_{RR^{}}^{12}\right)^2.`$ The dimensionless parameter $`u`$ is a measure of the disorder strength in units of the regular hopping. In this way, the Hamiltonian is written as the sum of a regular part, $`H^{(0)}`$, plus a disordered part, $`H_{imp}`$. For the regular hopping, we define $`t_{RR^{}}={\displaystyle \frac{1}{2}}\left(t_{RR^{}}^{12}+t_{RR^{}}^{21}\right),w_{RR^{}}={\displaystyle \frac{1}{2}}\left(t_{RR^{}}^{12}t_{RR^{}}^{21}\right),`$ so that the non disorderd part, $`H^{(0)}`$, of the Hamiltonian is $$H^{(0)}=\underset{R,R^{}}{}c_R^{}H_{RR^{}}^{(0)}c_R^{}^{}=\underset{R,R^{}}{}c_R^{}\left(t_{RR^{}}\sigma _1+iw_{RR^{}}\sigma _2\right)c_R^{}^{}.$$ (9) Since, for any lattice vector $`R_0`$, $`t_{RR^{}}=t_{R+R_0R^{}+R_0}`$, as well as $`w_{RR^{}}=w_{R+R_0R^{}+R_0}`$, and, moreover, $`t_{RR^{}}=t_{R^{}R}`$ while $`w_{RR^{}}=w_{R^{}R}`$, the Fourier transforms satisfy $`t_k=t_k^{}`$ ($`t_k`$ real) and $`w_k=w_k^{}`$ ($`w_k`$ imaginary). In momentum space, the Hamiltonian matrix, $`H_k^{(0)}=t_k\sigma _1+iw_k\sigma _2`$, is diagonalized by the unitary transformation $`c_k^{}=U_kd_k^{}`$, with $$U_k=\mathrm{e}^{i\frac{\pi }{4}\sigma _2}\mathrm{e}^{i\frac{\theta _k}{2}\sigma _1},$$ (10) where $$\mathrm{tan}\theta _k=\frac{iw_k}{t_k}=\frac{mt_k^{12}}{et_k^{12}}.$$ (11) Indeed, $`U_k^{}H_k^{(0)}U_k=ϵ_k\sigma _3`$, where $`ϵ_k^2=t_k^2w_k^2=\|t_k\|^2+\|w_k\|^2`$. ### A Current operator The commutator of the density $`c_R^{}c_R^{}`$ with the non disordered Hamiltonian (9) is $`{\displaystyle \underset{R^{}}{}}c_R^{}\left(t_{RR^{}}\sigma _1+iw_{RR^{}}\sigma _2\right)c_R^{}^{}c_R^{}^{}\left(t_{R^{}R}\sigma _1+iw_{RR^{}}\sigma _2\right)c_R^{},`$ or, in Fourier space, $`{\displaystyle \underset{k}{}}\left(t_{k+q}t_k\right)c_k^{}\sigma _1c_{k+q}^{}+i\left(w_{k+q}w_k\right)c_k^{}\sigma _2c_{k+q}^{}.`$ Therefore, in the long wavelength limit, the current operator in the absence of disorder is $`\stackrel{}{J}_q^{(0)}`$ $`{\displaystyle \underset{k}{}}\stackrel{}{}_kt_kc_k^{}\sigma _1c_{k+q}^{}+i\stackrel{}{}_kw_kc_k^{}\sigma _2c_{k+q}^{}`$ (14) $`={\displaystyle \underset{k}{}}\stackrel{}{}_kϵ_kc_k^{}\stackrel{}{B}_k\stackrel{}{\sigma }c_k^{}+ϵ_k\stackrel{}{}\theta _kc_k^{}\stackrel{}{B}_{,k}\stackrel{}{\sigma }c_k^{}`$ $`={\displaystyle \underset{k}{}}\stackrel{}{}ϵ_kd_k^{}\sigma _3d_{k+q}^{}+ϵ_k\stackrel{}{}\theta _kd_k^{}\sigma _2d_{k+q}^{},`$ the last being the expression in the basis which diagonalizes the Hamiltonian $`H_0`$. In (14), $`\stackrel{}{\sigma }=(\sigma _1,\sigma _2,\sigma _3)`$, and the vectors $$\stackrel{}{B}_k=(\mathrm{cos}\theta _k,\mathrm{sin}\theta _k,0),\stackrel{}{B}_{,k}=(\mathrm{sin}\theta _k,\mathrm{cos}\theta _k,0),$$ (15) describe intra and inter-band contributions to the current vertex. Notice that the regular hopping Hamiltonian can be simply written as $`H^{(0)}=_kϵ_kc_k^{}\stackrel{}{B}_k\stackrel{}{\sigma }c_k^{}`$. Moreover, since also the impurity part of the whole Hamiltonian, (6), does not commute with the density operator, in the disordered model the current operator acquires an additional term proportional to the random hopping matrix elements, which we discuss later. ## III Path Integral The starting point of our analysis is a path-integral representation of the generating functional, in terms of Grassmann variables, following the work by Efetov, Larkin and Khmelnitskii. To this end, we introduce, for each unit cell $`R`$, the Grassmann variables $`c_{R;a,p,\alpha }`$ and their complex conjugates $`\overline{c}_{R;a,p,\alpha }`$, where $`a=1,2`$ is the sublattice index, $`p=\pm `$ is the index of the advanced (+) and retarded (-) components, and the index $`\alpha `$ runs over $`m`$ identical copies of the model, as in the usual replica trick method. In what follows, by convention, the Pauli matrices $`\sigma `$’s act in the two sublattice space, the $`\tau `$’s in the space of the Grassmann fields $`c`$ and $`\overline{c}`$, and the $`s`$’s in the $`\pm `$ space. In order to treat on equal footing both the particle-hole and the particle-particle diffusive modes (diffusons and cooperons, respectively), as implied by time-reversal invariance, it is convenient to introduce the Nambu spinors $`\mathrm{\Psi }_R`$ and $`\overline{\mathrm{\Psi }}_R`$ defined through $`\mathrm{\Psi }_R={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{c}\overline{c}_R\\ c_R\end{array}\right),`$ where $`\overline{c}_R`$ and $`c_R`$ are column vectors with components $`\overline{c}_{R;a,p,\alpha }`$ and $`c_{R;a,p,\alpha }`$, respectively, and $`\overline{\mathrm{\Psi }}_R=\left[c\mathrm{\Psi }\right]^t,`$ where $`c=i\tau _2`$ is the charge conjugation operator. The action in terms of the spinors is \[see Eq.(8)\] $`S`$ $`=`$ $`{\displaystyle \underset{RR^{}}{}}\overline{\mathrm{\Psi }}_R\left(E+i{\displaystyle \frac{\omega }{2}}s_3H_{RR^{}}\right)\mathrm{\Psi }_R^{}`$ (16) $`=`$ $`{\displaystyle \underset{RR^{}}{}}\overline{\mathrm{\Psi }}_R\left(E+i{\displaystyle \frac{\omega }{2}}s_3H_{RR^{}}^{(0)}\right)\mathrm{\Psi }_R^{}+{\displaystyle \underset{RR^{}}{}}\overline{\mathrm{\Psi }}_RH_{imp,RR^{}}\mathrm{\Psi }_R^{}`$ (17) $`=`$ $`S_0+S_{imp},`$ (18) where $`E\pm i\omega /2`$ are the complex energies of the advanced/retarded components. ### A Disorder average Before taking the disorder average, we notice that, in the spinor notation, $`c_{1R}^{}c_{2R^{}}^{}+H.c.2\overline{\mathrm{\Psi }}_{1R}\mathrm{\Psi }_{2R^{}}=2\overline{\mathrm{\Psi }}_{2R^{}}\mathrm{\Psi }_{1R}.`$ Therefore, the impurity part of the action can be written as $`S_{imp}={\displaystyle \underset{R,R^{}}{}}2\tau _{RR^{}}^{12}\overline{\mathrm{\Psi }}_{1R}\mathrm{\Psi }_{2R^{}}.`$ The generating functional, within the replica method, is $$𝒵=𝒟\overline{\mathrm{\Psi }}𝒟\mathrm{\Psi }𝒟\tau P[\tau ]\mathrm{e}^{S_0S_{imp}},$$ (19) where $`P[\tau ]`$ is the gaussian probability distribution of the random bonds $`\tau _{RR^{}}^{12}`$. The average over disorder changes the impurity action into $`S_{imp}`$ $`=`$ $`{\displaystyle \underset{R,R^{}}{}}2u^2\left(t_{RR^{}}^{12}\right)^2\left(\overline{\mathrm{\Psi }}_{1R}\mathrm{\Psi }_{2R^{}}\right)^2`$ (20) $`=`$ $`{\displaystyle \underset{R,R^{}}{}}2u^2\left(t_{RR^{}}^{12}\right)^2\left(\overline{\mathrm{\Psi }}_{1R}\mathrm{\Psi }_{2R^{}}\right)\left(\overline{\mathrm{\Psi }}_{2R^{}}\mathrm{\Psi }_{1R}\right).`$ (21) We define $`W_{RR^{}}=2u^2\left(t_{RR^{}}^{12}\right)^2e,`$ so that $`W_q^{}=W_q`$, and introduce $`X_{1R}^{\alpha \beta }=\overline{\mathrm{\Psi }}_{1R}^\alpha \mathrm{\Psi }_{1R}^\beta ,`$ where $`\alpha `$ is a multilabel for Nambu, advanced/retarded and replica components, and analogously $`X_{2R}^{\alpha \beta }`$, as well as their Fourier transforms. By these definitions, $$S_{imp}=\frac{1}{V}\underset{q}{}\underset{\alpha ,\beta }{}W_qX_{1,q}^{\alpha \beta }X_{2,q}^{\beta \alpha }.$$ (22) This form, as compared to (21), has the advantage to allow a simple Hubbard–Stratonovich transformation. Notice that the use of Nambu spinors has the great advantage to involve just a single Fourier component of $`W_{RR^{}}`$. If we write $`W_q=\omega _q\mathrm{e}^{i\varphi _q},`$ where $`\omega _q>0`$ and $`\varphi _q=\varphi _q`$, and define $`Y_{1q}=\mathrm{e}^{i\frac{\varphi _q}{2}}X_{1q},Y_{2q}=\mathrm{e}^{i\frac{\varphi _q}{2}}X_{2q},`$ $`S_{imp}`$ takes the simple form $`S_{imp}`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}\omega _qY_{1,q}^{\alpha \beta }Y_{2,q}^{\beta \alpha }={\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{\omega _q}{2}}Tr\left[Y_{1,q}Y_{2,q}+Y_{2,q}Y_{1,q}\right]`$ (23) $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{\omega _q}{4}}Tr\left[Y_{0,q}Y_{0,q}Y_{3,q}Y_{3,q}\right],`$ (24) where we have introduced $`Y_0=Y_1+Y_2`$ as well as $`Y_3=Y_1Y_2`$. Moreover, our choice of the impurity potential, which does not break, on average, the spatial symmetries of the lattice, implies that $`\varphi _q=\theta _q`$, see Eq.(11). We notice that, if a term is written as $`\lambda {\displaystyle \frac{A^2}{4}}{\displaystyle \underset{\alpha \beta }{}}X^{\alpha \beta }X^{\beta \alpha }=\lambda {\displaystyle \frac{A^2}{4}}Tr\left(X^2\right),`$ where $`X=X^{}`$ and $`\lambda =\pm 1`$, one can always decouple it, by introducing an hermitean matrix $`Q`$, by the following Hubbard–Stratonovich transformation $$\mathrm{exp}\left[\lambda \frac{A^2}{4}Tr\left(X^2\right)\right]=N𝒟Q\mathrm{exp}\left[A^2Tr(Q^2)+\sqrt{\lambda }Tr\left(QX^t\right)\right],$$ (25) where the normalization factor $`N^1=𝒟Q\mathrm{exp}\left[A^2Tr(Q^2)\right]`$. In the specific example, (24) can be transformed into $`S_{imp}`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{\omega _q}}Tr\left[Q_{0q}Q_{0q}+Q_{3q}Q_{3q}\right]{\displaystyle \frac{i}{V}}{\displaystyle \underset{q}{}}Tr\left[Q_{0q}Y_{0q}^t+iQ_{3q}Y_{3q}^t\right]`$ (26) $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{\omega _q}}Tr\left[Q_{0q}Q_{0q}+Q_{3q}Q_{3q}\right]`$ (27) $``$ $`{\displaystyle \frac{i}{V}}{\displaystyle \underset{q}{}}Tr\left[Q_{0q}\left(\mathrm{cos}{\displaystyle \frac{\varphi _q}{2}}X_{0q}^t+i\mathrm{sin}{\displaystyle \frac{\varphi _q}{2}}X_{3q}^t\right)\right]`$ (28) $``$ $`{\displaystyle \frac{i}{V}}{\displaystyle \underset{q}{}}iTr\left[Q_{3q}\left(\mathrm{cos}{\displaystyle \frac{\varphi _q}{2}}X_{3q}^t+i\mathrm{sin}{\displaystyle \frac{\varphi _q}{2}}X_{0q}^t\right)\right].`$ (29) If we define $`Q_q=Q_{0q}\sigma _0+iQ_{3q}\sigma _3`$, we obtain $`S_{imp}`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{2\omega _q}}Tr\left[Q_q^{}Q_q^{}\right]{\displaystyle \frac{i}{V}}{\displaystyle \underset{p,q}{}}\overline{\mathrm{\Psi }}_pQ_q\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _3}\mathrm{\Psi }_{p+q}`$ (30) $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{2\omega _q}}Tr\left[Q_q^{}Q_q^{}\right]i{\displaystyle \underset{R}{}}\overline{\mathrm{\Psi }}_RQ_R\mathrm{\Psi }_R+{\displaystyle \frac{i}{V}}{\displaystyle \underset{p,q}{}}\overline{\mathrm{\Psi }}_p\left(1\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _3}\right)Q_q\mathrm{\Psi }_{p+q},`$ (31) where the last term vanishes at $`q=0`$. Notice that the two-sublattice symmetry properties of the Hamiltonian, see Eq. (1), are reflected in the particular form of the $`Q`$-matrix, which only contains the Pauli matrices $`\sigma _0`$ and $`\sigma _3`$. In particular, since the tensor $`Q(R)\mathrm{\Psi }(R)\overline{\mathrm{\Psi }}(R)`$, $`\sigma _0`$ selects the uniform component while $`\sigma _3`$ the staggered component of the product of the two Grassmann fields. Notice that the electron-$`Q`$ coupling in the $`\mathrm{\Psi }_R`$ basis can be written as a local coupling but for the last term in (31). To simplify the notation of this term we find it useful to define the operator $`\widehat{L}`$ through $$\widehat{L}Q(R)=\frac{1}{V}\underset{q}{}\mathrm{e}^{iqR}\left(1\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _3}\right)Q_q.$$ (32) If we now transform the spinors in Eq.(30) to the diagonal basis, the coupling to the $`Q_q`$ matrix transforms as $`U_{p+q}^{}Q_q\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _3}U_p^{}`$. The simplest consequence is that, in the diagonal basis, $`Q_{k_1,k_2}=Q_{0,k_1,k_2}\sigma _0+iQ_{1,k_1,k_2}\sigma _1`$ depends on two wavevectors. However, in the case of cubic lattices, see Eq.(A1), $`U_{p+q}^{}Q_q\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _3}U_{p+q}^{}=\mathrm{e}^{\frac{i}{2}\theta _{p+q}\sigma _1}\left(Q_{0,q}iQ_{3,q}\sigma _1\right)\mathrm{e}^{\frac{i}{2}\theta _p\sigma _1}\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _1}`$ $`=\mathrm{e}^{\frac{i}{2}(\theta _q+\varphi _q)\sigma _1}\left(Q_{0,q}iQ_{3,q}\sigma _1\right)=Q_{0,q}iQ_{3,q}\sigma _1,`$ since $`\varphi _q=\theta _q`$. Therefore, for cubic lattices, we can also write $$S_{imp}=\frac{1}{V}\underset{q}{}\frac{1}{2\omega _q}Tr\left[Q_q^{}Q_q^{}\right]i\underset{R}{}\overline{\mathrm{\Phi }}_RQ(R)\mathrm{\Phi }_R,$$ (33) where now $`\mathrm{\Phi }_R`$ is the Grassmann field of the $`d_R^{}`$ operators, and the matrix $`Q(R)=Q_0(R)\sigma _0+iQ_1(R)\sigma _1`$, with $`Q_1(R)=Q_3(R)`$ defined above. Finally, we notice that, in the case of the honeycomb lattice, at the wavevector $`q_{}`$, connecting the two Dirac cones of the non-disordered dispersion band, $`\omega _q_{}=0`$. This observation will turn useful when discussing the long-wavelength behavior of the model. ## IV Symmetries The action (18), at $`E=\omega =0`$, i.e. at the band center with zero complex frequency, is invariant under a transformation $`\mathrm{\Psi }_RT\mathrm{\Psi }_R`$ if $`cT^tc^tH_{RR^{}}T=H_{RR^{}}.`$ Since the random matrix elements $`H_{RR^{}}`$ involve both Pauli matrices $`\sigma _1`$ and $`\sigma _2`$, $`T`$ has to satisfy at the same time $`cT^tc^t\sigma _1T=\sigma _1`$ and $`cT^tc^t\sigma _2T=\sigma _2`$. This implies that $`cT^tc^t`$ $`=`$ $`\sigma _1T^1\sigma _1`$ (34) $`\sigma _1T^1\sigma _1`$ $`=`$ $`\sigma _2T^1\sigma _2.`$ (35) The condition (35) can be fulfilled only by a transformation $`T=T_0\sigma _0+T_3\sigma _3`$ Under such a transformation $`QcT^tc^tQT=\sigma _1T^1\sigma _1QT\sigma _2T^1\sigma _2QT`$ Since $`\sigma _1Q\sigma _1=Q^{}`$, then $`T^1\sigma _1QT\sigma _1=T^{}Q^{}\sigma _1\left(T^1\right)^{}\sigma _1=T^{}\sigma _1Q\left(T^1\right)^{}\sigma _1.`$ Hence, the transformation is also unitary, $`T^{}=T^1`$. Moreover, such a transformation leaves the $`Q`$-manifold invariant, which implies that our Hubbard-Stratonovich decoupling scheme, which makes use of $`Q=Q_0\sigma _0+iQ_3\sigma _3`$, is exhaustive. The unitary transformation, $`T`$, can be written as $$T=\mathrm{exp}\left[\frac{W_0}{2}\sigma _0+\frac{W_3}{2}\sigma _3\right],$$ (36) where $`W_0^{}=W_0,W_3^{}=W_3.`$ In addition, we must impose the charge conjugacy invariance, which, through Eq.(34), implies that $`\begin{array}{ccccc}cW_0^tc^t& =& W_0& =& W_0^{},\\ cW_3^tc^t& =& W_3& =& W_3^{}.\end{array}`$ The number of independent parameters turns out to be $`16m^2`$, which suggests that $`T`$ is related to a unitary group, specificaly U$`(4m)`$, as argued by Gade and Wegner. In fact, we can alternatively write $$T=\left(\begin{array}{cc}\mathrm{e}^{\frac{W_0+W_3}{2}}& 0\\ 0& \mathrm{e}^{\frac{W_0W_3}{2}}\end{array}\right)\left(\begin{array}{cc}U& 0\\ 0& cU^{}c^t\end{array}\right),$$ (37) where $`U`$ is indeed a unitary transformation belonging to U$`(4m)`$. The invariance of (18) at finite frequency, $`\omega 0`$, implies the additional condition $`cT^tc^ts_3T=s_3`$, which reduces the number of independent parameters to $`8m^2+2m`$, lowering the symmetry of $`U`$ down to Sp$`(2m)`$. If $`E0`$, $`T`$ has to satisfy also $`cT^tc^tT=\sigma _1T^1\sigma _1T=1.`$ This implies that, at finite energy, i.e. away from the band center, $`T`$ does not contain anymore a $`\sigma _3`$-component. Indeed $`E`$ lowers the symmetry of $`U`$ down to Sp$`(2m)`$, which is further reduced to $`\mathrm{Sp}(m)\times \mathrm{Sp}(m)`$ by a finite frequency, as in the standard situation. ## V Saddle Point The full action $`S`$ $`=`$ $`{\displaystyle \underset{k,q}{}}\overline{\mathrm{\Psi }}_k\left(E\delta _{q0}+i{\displaystyle \frac{\omega }{2}}s_3\delta _{q0}H_k^{(0)}\delta _{q0}+{\displaystyle \frac{i}{V}}Q_q\mathrm{e}^{\frac{i}{2}\varphi _q\sigma _3}\right)\mathrm{\Psi }_{k+q}`$ (38) $`+`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{2\omega _q}}Tr\left[Q_q^{}Q_q^{}\right],`$ (39) by integrating over the Nambu spinors, transforms into $$S[Q]=\frac{1}{V}\underset{q}{}\frac{1}{2\omega _q}Tr\left[Q_qQ_q^{}\right]\frac{1}{2}Tr\mathrm{ln}\left[E+i\frac{\omega }{2}s_3H^{(0)}+iQi\widehat{L}Q\right].$$ (40) The saddle point equation for homogeneous solutions at $`E=0`$ and $`\omega =0^+`$ is $`Q=i{\displaystyle \frac{\omega _0}{4}}{\displaystyle \frac{d^2k}{4\pi ^2}\left(i0^+s_3ϵ_k+iQ\right)^1}+\left(i0^+s_3+ϵ_k+iQ\right)^1,`$ where $`\omega _0=_{RR^{}}2u^2(t_{RR^{}}^{12})^2`$. The general solution is $$Q_{sp}=\frac{\pi }{4}\omega _0\rho (0)s_3\mathrm{\Sigma }s_3,$$ (41) with $`\rho (0)`$ being the density of states at $`E=0`$. In order to distinguish transverse from longitudinal modes, it is convenient to parametrize the $`Q`$-matrix in the following way $$Q(R)_P=\sigma _1T(R)^1\sigma _1\left[Q_{sp}+P(R)\right]T(R)Q(R)+\sigma _1T(R)^1\sigma _1P(R)T(R).$$ (42) Here $`P(R)`$ describes the longitudinal modes, which we discuss more in detail in section VII, and $`T`$ the transverse modes. Namely, $`T`$ has the form given in Eq.(37), $$T(R)=\mathrm{exp}\left(\frac{W(R)}{2}\right)=\mathrm{exp}\left(\frac{W_0(R)}{2}\sigma _0+\frac{W_3(R)}{2}\sigma _3\right),$$ (43) with $`\mathrm{exp}[(W_0+W_3)/2]`$ belonging now to the coset space $`\mathrm{U}(4m)/\mathrm{Sp}(2m)`$. This amounts to impose that $`\{W_0,s\}=0,[W_3,s]=0,`$ by which it derives that $$Q(R)=\sigma _1T(R)^1\sigma _1Q_{sp}T(R)=Q_{sp}\mathrm{e}^{W_0(R)\sigma _0+W_3(R)\sigma _3}.$$ (44) In the $`\pm `$ space, we can write $$W_0=\left(\begin{array}{cc}0& B\\ B^{}& 0\end{array}\right),W_3=\left(\begin{array}{cc}iA& 0\\ 0& iC\end{array}\right),$$ (45) where $`A^{}=A`$, $`C^{}=C`$, and additionally, since $`cW^tc^t=\sigma _1W^{}\sigma _1=\sigma _1W\sigma _1`$, then $`cA^tc^t=A`$, $`cC^tc^t=C`$ and $`cB^tc^t=B^{}`$. By writing $`A`$ $`=`$ $`A_0\tau _0+i\left(A_1\tau _1+A_2\tau _2+A_3\tau _3\right),`$ (46) $`B`$ $`=`$ $`B_0\tau _0+i\left(B_1\tau _1+B_2\tau _2+B_3\tau _3\right),`$ (47) $`C`$ $`=`$ $`C_0\tau _0+i\left(C_1\tau _1+C_2\tau _2+C_3\tau _3\right),`$ (48) we find that the above conditions imply that, for $`i=0,\mathrm{},3`$, $$B_i,A_i,C_ie,$$ (49) and $$A_0=A_0^t,C_0=C_0^t,$$ (50) while, for $`j=1,2,3`$, $$A_j=A_j^t,C_j=C_j^t.$$ (51) ## VI Effective Action In this section, we derive the effective field theory describing the long wavelength transverse fluctuations of $`Q(R)`$ around the saddle point. In the case of honeycomb lattices, we should worry about the momentum component of $`Q`$ which couples the two Dirac cones. However, one can see that the free action of $`Q`$ diverges at this wavevector, so that we are allowed to ignore the fluctuations around this momentum. ### A Integration over the Grassmann fields As we said, by integrating (39) over the Grassmann variables, we obtain the following action of $`Q`$: $$S[Q]=\frac{1}{V}\underset{q}{}\frac{1}{2\omega _q}Tr\left[Q_qQ_q^{}\right]+\frac{1}{2}Tr\mathrm{ln}\left[E+i\frac{\omega }{2}s_3H^{(0)}+iQi\widehat{L}Q\right].$$ (52) We start by neglecting the longitudinal fluctuations. Then, since $`Q=\stackrel{~}{T}^{}Q_{sp}T`$, where we define $`\stackrel{~}{T}=\sigma _1T\sigma _1\sigma _2T\sigma _2`$, we can rewrite the second term of $`S[Q]`$ as $`{\displaystyle \frac{1}{2}}Tr\mathrm{ln}\left(E\stackrel{~}{T}T^{}+i{\displaystyle \frac{\omega }{2}}\stackrel{~}{T}s_3T^{}\stackrel{~}{T}H^{(0)}T^{}+iQ_{sp}V\right),`$ (53) where we define $`V=i\stackrel{~}{T}\widehat{L}QT^{}.`$ Since $`H_{RR^{}}^{(0)}`$ involves either $`\sigma _1`$ and $`\sigma _2`$, while $`T`$ involves $`\sigma _0`$ and $`\sigma _3`$, then $`H_{RR^{}}^{(0)}T(R^{})^{}=\stackrel{~}{T}(R^{})^{}H_{RR^{}}^{(0)}`$ $`=`$ $`\stackrel{~}{T}(R)^{}H_{RR^{}}^{(0)}+\left(\stackrel{~}{T}(R^{})^{}\stackrel{~}{T}(R)^{}\right)H_{RR^{}}^{(0)}`$ $`=`$ $`\stackrel{~}{T}(R)^{}H_{RR^{}}^{(0)}\stackrel{}{}\stackrel{~}{T}(R)^{}\left(\stackrel{}{R}\stackrel{}{R^{}}\right)H_{RR^{}}^{(0)}`$ $`+`$ $`{\displaystyle \frac{1}{2}}_{ij}\stackrel{~}{T}(R)^{}\left(R_iR_i^{}\right)\left(R_jR_j^{}\right)H_{RR^{}}^{(0)}.`$ Therefore the term $`\stackrel{~}{T}H^{(0)}T^{}`$ which appears in (53) can be written at long wavelengths as $`\stackrel{~}{T}(R)H_{RR^{}}^{(0)}T(R^{})^{}`$ $`=`$ $`H_{RR^{}}^{(0)}i\stackrel{~}{T}(R)\stackrel{}{}\stackrel{~}{T}^1(R)\stackrel{}{J}^{(0)}(R^{})_{RR^{}}`$ (54) $`+`$ $`{\displaystyle \frac{1}{2}}\stackrel{~}{T}(R)_{ij}\stackrel{~}{T}(R)^1\left(R_iR_i^{}\right)\left(R_jR_j^{}\right)H_{RR^{}}^{(0)}`$ (56) $`H_{RR^{}}^{(0)}+U_{RR^{}},`$ where we used the fact that the long-wavelength part of the current operator in real space is, see (B3), $`\stackrel{}{J}(R^{})=i{\displaystyle \underset{R}{}}\left(\stackrel{}{R}\stackrel{}{R^{}}\right)c_R^{}H_{RR^{}}c_R^{}^{}.`$ Notice that in (56) only the current vertex which derives from the regular hopping appears, which is not the full current operator. Moreover, a further current-like coupling will arise from the expansion in $`V`$ (see below). To this end, in the long-wavelength limit, the operator $`\widehat{L}`$, see (32), can be approximately written as $$\widehat{L}Q(R)\stackrel{}{\beta }\stackrel{}{}Q(R)\sigma _3\frac{1}{2}\left(\stackrel{}{\beta }\stackrel{}{}\right)^2Q(R),$$ (57) where $`\stackrel{}{\beta }=\stackrel{}{}\varphi (q=0)/2`$. Having defined $`U`$, (53) can be written as $`{\displaystyle \frac{1}{2}}Tr\mathrm{ln}\left(E\stackrel{~}{T}T^{}+i{\displaystyle \frac{\omega }{2}}\stackrel{~}{T}s_3T^{}UVH^{(0)}+iQ_{sp}\right)`$ (58) $`=`$ $`{\displaystyle \frac{1}{2}}Tr\mathrm{ln}G+{\displaystyle \frac{1}{2}}Tr\mathrm{ln}\left(1+GE\stackrel{~}{T}T^{}+Gi{\displaystyle \frac{\omega }{2}}\stackrel{~}{T}s_3T^{}GUGV\right),`$ (59) where $`G=(H^{(0)}+iQ_{sp})^1`$ is the Green’s function in the absence of transverse fluctuations. The effective field theory is then derived by expanding $`S[Q]`$ up to second order in $`U`$ and $`V`$, and first order in $`E`$ and $`\omega `$. In this way we get the following terms. ### B Expansion in the $`Q`$ free action The free part of the action $`S_0[Q]={\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{2\omega _q}}Tr\left[Q_qQ_q^{}\right],`$ can be expanded at small q. Since $`\omega _q=\omega _q`$, then $`\omega _q\omega _0(1\gamma q^2),`$ leading to $`S_0[Q]`$ $``$ $`{\displaystyle \frac{1}{2\omega _0}}{\displaystyle 𝑑RTr\left[Q(R)Q(R)^{}\right]}+{\displaystyle \frac{\gamma }{2V\omega _0}}{\displaystyle \underset{q}{}}q^2Tr\left[Q_qQ_q^{}\right]`$ (60) $`=`$ $`{\displaystyle \frac{1}{2\omega _0}}{\displaystyle 𝑑RTr\left[Q_{sp}^2\right]}+{\displaystyle \frac{\gamma }{2\omega _0}}{\displaystyle 𝑑RTr\left[\stackrel{}{}Q(R)\stackrel{}{}Q(R)^{}\right]}.`$ (61) The second term is a contribution to the current-current correlation function of the part of the current vertices proportional to the random hopping. ### C Expansion in $`E`$ Expansion of (59) in $`E`$ gives $$\frac{E}{2}Tr\left(G\stackrel{~}{T}T^{}\right)=i\frac{E}{\omega _0}Tr\left(Q_{sp}\stackrel{~}{T}T^{}\right)=i\frac{E}{\omega _0}TrQ.$$ (62) ### D Expansion in $`\omega `$ Expansion in $`\omega `$ gives $$i\frac{\omega }{4}Tr\left(G\stackrel{~}{T}\widehat{s}T^{}\right)=\frac{\omega }{2\omega _0}Tr\left(s_3Q\right).$$ (63) ### E Expansion in $`U`$ The second order expansion in $`U`$ contains the following terms: $$\frac{1}{2}Tr\left(GU\right),$$ (64) and $$\frac{1}{4}Tr\left(GUGU\right),$$ (65) Taking in (64), the component of $`U`$ containing second derivatives, we get $`{\displaystyle \frac{1}{4}}Tr\left\{\stackrel{~}{T}(R)_{ij}\stackrel{~}{T}(R)^1\left(R_iR_i^{}\right)\left(R_jR_j^{}\right)H_{RR^{}}^{(0)}G(R^{},R)\right\}.`$ By means of the Ward identity (B6), the above expression turns out to be $$\frac{\chi _{ij}^{++}}{8}Tr\left\{\stackrel{~}{T}(R)_{ij}\stackrel{~}{T}(R)^1\right\},$$ (66) which, integrating by part, is also equal to $``$ $`{\displaystyle \frac{\chi _{ij}^{++}}{8}}Tr\left\{\stackrel{~}{T}(R)_i\stackrel{~}{T}(R)^1\stackrel{~}{T}(R)_j\stackrel{~}{T}(R)^1\right\}`$ (67) $`=`$ $`{\displaystyle \frac{1}{16}}\chi _{ij}^{++}Tr\left[D_iD_jD_is_3\sigma _1D_js_3\sigma _1\right]`$ (68) $``$ $`{\displaystyle \frac{1}{16}}\chi _{ij}^{++}Tr\left[D_iD_j+D_is_3\sigma _1D_js_3\sigma _1\right].`$ (69) Here we have introduced a matrix $`\stackrel{}{D}(R)`$ with the $`i`$-th component $$D_i(R)=D_{0,i}(R)\sigma _0+D_{3,i}(R)\sigma _3\stackrel{~}{T}(R)_i\stackrel{~}{T}(R)^1.$$ (70) The part of (64) which contains first derivatives gives rise to a boundary term $$\left[\frac{1}{V}\underset{k}{}\stackrel{}{}\theta _k\frac{ϵ_k^2}{ϵ_k^2+\mathrm{\Sigma }^2}\right]𝑑RTr\left[\stackrel{}{}W(R)\sigma _3\right],$$ (71) where $`\mathrm{\Sigma }`$ has been defined by Eq.(41), which we discard by taking appropriate boundary conditions. Let us now analyse the term (65), where we have to keep of $`U`$ only the part containing first derivatives. By making use of (56), this term is, in momentum space, $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{kq}{}}Tr\left\{\left[\stackrel{~}{T}\stackrel{}{}\stackrel{~}{T}^1\right]_q\stackrel{}{J}_{k+q}^{(0)}G(k+q)\left[\stackrel{~}{T}\stackrel{}{}\stackrel{~}{T}^1\right]_q\stackrel{}{J}_k^{(0)}G(k)\right\}`$ (72) $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{kq}{}}Tr\left\{\left[\stackrel{~}{T}\stackrel{}{}\stackrel{~}{T}^1\right]_q\stackrel{}{J}_k^{(0)}G(k)\left[\stackrel{~}{T}\stackrel{}{}\stackrel{~}{T}^1\right]_q\stackrel{}{J}_k^{(0)}G(k)\right\}`$ (73) $`={\displaystyle \frac{1}{4}}{\displaystyle \underset{k}{}}{\displaystyle \underset{R}{}}Tr\left\{\stackrel{}{D}(R)\stackrel{}{J}_k^{(0)}G(k)\stackrel{}{D}(R)\stackrel{}{J}_k^{(0)}G(k)\right\},`$ (74) (75) valid for small $`q`$. Here the matrix $`\stackrel{}{J}_k^{(0)}=\stackrel{}{}_kt_k\sigma _1+i\stackrel{}{}_kw_k\sigma _2`$. The non vanishing terms in (75) have both $`\stackrel{}{D}`$’s either $`\stackrel{}{D}_0\sigma _0`$ or $`\stackrel{}{D}_3\sigma _3`$ \[see (70)\]. In the diagonal basis, upon defining, as we did in Eq.(41), $`Q_{sp}=\mathrm{\Sigma }s_3`$, with $`\mathrm{\Sigma }=\pi \omega _0\rho (0)/4`$, the Green’s function is $`𝒢(k)`$ $`=`$ $`{\displaystyle \frac{1}{ϵ_k\sigma _3+i\frac{\pi }{4}\omega _0\rho (0)s_3}}`$ (76) $`=`$ $`i{\displaystyle \frac{\mathrm{\Sigma }}{ϵ_k^2+\mathrm{\Sigma }^2}}\sigma _0s_3{\displaystyle \frac{ϵ_k}{ϵ_k^2+\mathrm{\Sigma }^2}}\sigma _3s_0`$ (77) $``$ $`G_0(k)\sigma _0s_3+G_3(k)\sigma _3s_0.`$ (78) Going back to the original basis, $$G(k)=U_k𝒢(k)U_k^{}=G_0(k)\sigma _0s_3+G_3(k)\stackrel{}{B}_k\stackrel{}{\sigma }s_0,$$ (79) where $`\stackrel{}{B}_k`$ has been defined in (15). Therefore, $$\sigma _3G(k)\sigma _3=G_0(k)\sigma _0s_3G_3(k)\stackrel{}{B}_k\stackrel{}{\sigma }s_0=s_1G(k)s_1.$$ (80) By means of (80), we find that $$\sigma _3\stackrel{}{J}_k^{(0)}G(k)^+\sigma _3=\stackrel{}{J}_k^{(0)}\sigma _3G(k)^+\sigma _3=\stackrel{}{J}_k^{(0)}G(k)^{},$$ (81) from which it derives that (75) can be written as the sum of two different terms $`{\displaystyle \frac{1}{16}}\chi _{ij}^+Tr\left[D_iD_jD_is_3\sigma _1D_js_3\sigma _1\right]`$ (82) $`+{\displaystyle \frac{1}{16}}\chi _{ij}^{++}Tr\left[D_iD_j+D_is_3\sigma _1D_js_3\sigma _1\right].`$ (83) By summing (82), (83), (68) and (69), we get $`{\displaystyle \frac{1}{16}}\left(\chi _{ij}^+\chi _{ij}^{++}\right)Tr\left[D_iD_jD_is_3\sigma _1D_js_3\sigma _1\right],`$ which is equal to $$\frac{2\pi }{32\mathrm{\Sigma }^2}\sigma _{ij}^{(0)}𝑑RTr\left(_iQ(R)_jQ(R)^{}\right),$$ (84) where $`\sigma _{ij}^{(0)}={\displaystyle \frac{1}{2\pi }}\left(\chi _{ij}^+\chi _{ij}^{++}\right)`$ is the Kubo conductivity with the current vertices which involve only the regular hopping \[cf. Eq.(B4\]. ### F Expansion in $`V`$ The expansion in $`V`$ up to second order, gives two terms $$\frac{1}{2}Tr\left(GV\right),$$ (85) and $$\frac{1}{4}Tr\left(GVGV\right),$$ (86) In addition, we must also consider the mixed term $$\frac{1}{2}Tr\left(GUGV\right).$$ (87) The first order term (85) is $`{\displaystyle \frac{1}{2}}Tr\left(GV\right)={\displaystyle \frac{i}{\omega _0}}Tr\left(Q_{sp}V\right)`$ (88) $`={\displaystyle \frac{1}{\omega _0}}{\displaystyle 𝑑RTr\left(Q(R)^{}\stackrel{}{\beta }\stackrel{}{}Q(R)\sigma _3\right)}`$ (89) $`{\displaystyle \frac{1}{2\omega _0}}{\displaystyle 𝑑RTr\left(\stackrel{}{\beta }\stackrel{}{}Q(R)^{}\stackrel{}{\beta }\stackrel{}{}Q(R)\right)}.`$ (90) The first term (89) is another boundary term, which we neglect. The second order term, Eq.(86), gives $$\frac{1}{4V}\underset{k}{}\left(G_0(k)^2+G_3(k)^2\right)𝑑RTr\left[\stackrel{}{\beta }\stackrel{}{}Q(R)^{}\stackrel{}{\beta }\stackrel{}{}Q(R)\right].$$ (91) Notice that the saddle point equation implies that $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{k}{}}G_3(k)^2G_0(k)^2={\displaystyle \frac{2}{\omega _0}}.`$ By using the above equation in (90), we find that (91) plus (90) give $`{\displaystyle \frac{1}{2V}}{\displaystyle \underset{k}{}}G_3(k)^2{\displaystyle 𝑑RTr\left[\stackrel{}{\beta }\stackrel{}{}Q(R)^{}\stackrel{}{\beta }\stackrel{}{}Q(R)\right]}`$ (92) $`={\displaystyle \frac{1}{2V}}{\displaystyle \underset{k}{}}{\displaystyle \frac{ϵ_k^2}{ϵ_k^2+\mathrm{\Sigma }^2}}{\displaystyle 𝑑RTr\left[\stackrel{}{\beta }\stackrel{}{}Q(R)^{}\stackrel{}{\beta }\stackrel{}{}Q(R)\right]}`$ (93) The contribution of the mixing term, (87), can be evaluated in a similar way, giving $$\frac{1}{2V}\underset{k}{}\frac{ϵ_k^2}{ϵ_k^2+\mathrm{\Sigma }^2}𝑑RTr\left[\stackrel{}{}\theta _k\stackrel{}{}Q(R)^{}\stackrel{}{\beta }\stackrel{}{}Q(R)\right].$$ (94) We notice that, since $`\varphi _q=\theta _q`$, then (93) and (94) can be included in (84) if the following redefinition of the current vertex is assumed: $$\stackrel{}{J}_k^{(0)}=\stackrel{}{}ϵ_k\stackrel{}{B}_k\stackrel{}{\sigma }+ϵ_k\left(\stackrel{}{}\theta _k\stackrel{}{}\theta _0\right)\stackrel{}{B}_{,k}\stackrel{}{\sigma }.$$ (95) Therefore the Kubo conductivity which appears in (84) has to be calculated with the above expression of the regular current vertex. This has notable consequences. First of all, for cubic lattices, the interband contribution vanishes, hence the Kubo conductivity coincides with the true one, since the enlargment of the unit cell was artificial. This is compatible with (33), where we showed that the impurity action is local in the basis which diagonalizes the regular hopping, implying that the regular current vertex contains only the intraband operator. To conclude, the effective action so far derived is therefore $`S[Q]`$ $`=`$ $`{\displaystyle \frac{1}{2\omega _0}}{\displaystyle 𝑑RTr\left[Q(R)Q(R)^{}\right]}`$ (97) $`+{\displaystyle \frac{\gamma }{2V\omega _0}}{\displaystyle \underset{q}{}}q^2Tr\left[Q_qQ_q^{}\right]`$ $`+`$ $`{\displaystyle \frac{2\pi }{32\mathrm{\Sigma }^2}}\sigma _{ij}^{(0)}{\displaystyle 𝑑RTr\left(_iQ(R)_jQ(R)^{}\right)}`$ (98) $`+`$ $`{\displaystyle 𝑑Ri\frac{E}{\omega _0}Tr\left(Q(R)\right)}{\displaystyle \frac{\omega }{2\omega _0}}Tr\left(s_3Q(R)\right).`$ (99) ## VII Longitudinal Fluctuations The full expression of the $`Q`$-matrix we must indeed consider is the one given by (42), $`Q(R)_P`$ $`=`$ $`\stackrel{~}{T}(R)^1\left[Q_{sp}+P(R)\right]T(R)Q(R)+\stackrel{~}{T}(R)^1P(R)T(R)`$ (100) $``$ $`Q(R)+S(R),`$ (101) where $`T(R)`$ involves transverse fluctuations and $$P(R)=\left(P_{00}s_0+P_{03}s_3\right)\sigma _0+i\left(P_{31}s_1+P_{32}s_2\right)\sigma _3,$$ (102) being all $`P`$’s hermitean. Charge conjugation implies that $`cP^tc^t=P`$. The field $`P(R)`$ takes into account longitudinal fluctuations which are massive. Let us define $`\mathrm{\Gamma }(RR^{})`$ the Fourier transform of $`\omega _q^1`$. Then, the free action of the $`Q_P`$ field is $`S[Q_P]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑R𝑑R^{}\mathrm{\Gamma }(RR^{})Tr\left[Q_P(R)Q_P(R^{})^{}\right]}`$ (103) $`=`$ $`{\displaystyle \frac{1}{2\omega _0}}{\displaystyle 𝑑RTr\left[Q_P(R)Q_P(R)^{}\right]}`$ (105) $`{\displaystyle \frac{1}{4}}{\displaystyle 𝑑R𝑑R^{}\mathrm{\Gamma }(RR^{})Tr\left[\left(Q_P(R)^{}Q_P(R^{})^{}\right)\left(Q_P(R)^{}Q_P(R^{})^{}\right)\right]}.`$ Since $`QQ^{}=Q_{sp}^2`$, (105) gives $`{\displaystyle \frac{1}{2\omega _0}}{\displaystyle 𝑑RTr\left[P(R)P(R)^{}+2Q_{sp}P(R)+Q_{sp}^2\right]}.`$ The second term cancels with the first order expansion of $`Tr\mathrm{ln}G_P`$, since $`Q_{sp}`$ is the saddle point solution. What is left, i.e. $$\frac{1}{2\omega _0}𝑑RTr\left[P(R)P(R)^{}\right],$$ (106) is actually the mass term of the longitudinal modes, since the second order expansion in $`P`$ of $`Tr\mathrm{ln}G_P`$ is zero. The other term, (105), can be analysed within a gradient expansion of $`Q_P(R)Q_P(R^{})=\stackrel{}{}Q_P(R)\left(\stackrel{}{R}\stackrel{}{R^{}}\right)+\mathrm{}`$. The details of the calculations are given in Appendix C 1, so that, in this section, we just present the final results. The free action of the longitudinal fields is found to be $$S_0[P]=\frac{1}{V}\underset{q}{}\frac{1}{2\omega _q}Tr(P_qP_q^{}).$$ (107) Here, we neglect the the contribution of the invariant measure, which, in the zero replica limit, gives rise to fluctuations smaller by a factor $`u^2`$ than Eq.(107). The integration over $`P`$ with the above action has several important consequences for the action of the transverse modes (see Appendix C 1). First of all, all the terms of the Kubo conductivity with the random current vertices are recovered. In addition, we find a new operator $$\frac{2\pi }{832\mathrm{\Sigma }^4}\mathrm{\Pi }𝑑RTr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right]Tr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right],$$ (108) which has contributions from two different terms. The first one is obtained by expanding each Green’s function in (65) at first order, $`G_P=G_0iG_0PG_0`$, and the second is derived by (105). They are analogous to the components of the Kubo conductivity with regular and with random current vertices, respectively. ## VIII Effective non–linear $`\sigma `$-model In conclusion the final expression of the action of the transverse modes in the long-wavelength limit is $`S[Q]`$ $`=`$ $`{\displaystyle \frac{2\pi }{32\mathrm{\Sigma }^2}}\sigma _{xx}{\displaystyle 𝑑RTr\left(\stackrel{}{}Q(R)\stackrel{}{}Q(R)^{}\right)}`$ (109) $`+`$ $`{\displaystyle 𝑑Ri\frac{E}{\omega _0}Tr\left(Q(R)\right)}{\displaystyle \frac{\omega }{2\omega _0}}Tr\left(s_3Q(R)\right)`$ (110) $``$ $`{\displaystyle \frac{2\pi }{832\mathrm{\Sigma }^4}}\mathrm{\Pi }{\displaystyle 𝑑RTr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right]Tr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right]},`$ (111) where we make use of the fact that, in the models we consider, $`\sigma _{ij}=\delta _{ij}\sigma _{xx}`$. Since $`Q(R)=Q_{sp}T(R)^2`$, see Eq.(44), expressing $`T(R)`$ by means of $`U(R)`$ as in Eq.(37), the action at $`E=\omega =0`$ can also be written as $`S[U]`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{16}}{\displaystyle 𝑑RTr\left[\stackrel{}{}U(R)^2\stackrel{}{}U(R)^2\right]}`$ (112) $``$ $`{\displaystyle \frac{2\pi }{32}}{\displaystyle \frac{\mathrm{\Pi }}{2}}{\displaystyle 𝑑R\left\{Tr\left[U(R)^2\stackrel{}{}U(R)^2\right]\right\}^2},`$ (113) as anticipated in the section I A. As compared to the non-linear $`\sigma `$-model which is obtained in the absence of the particle-hole symmetry, the above action differs first of all because of the symmetry properties of the matrix field $`U(R)`$, which describes now the Goldstone modes within the coset space $`\mathrm{U}(4m)/\mathrm{Sp}(2m)`$. Moreover, it also differs for the last term of (111), which, in the general case, even if present, is not related to massless modes. An analogous term was originally obtained by Gade in a two sublattice model described by two on-site levels with a regular hopping of the form $`H_{RR^{}}=t_{RR^{}}\sigma _1`$, and a local time-reversal symmetry breaking random potential $`H_{imp,R}=w_{1R}\sigma _1+w_{2R}\sigma _2`$, which we discuss in section XII. Although the action may be parametrized by a simple unitary field $`U(R)`$ as in (113), we prefer to work with the matrix $`Q(R)`$ which has the more transparent physical interpretation $`Q(R)\mathrm{\Psi }(R)\overline{\mathrm{\Psi }}(R)`$. Finally, it is important to notice that either $`\sigma _{xx}`$ and $`\mathrm{\Pi }`$ have contributions from both the regular and the random current vertices. This implies that, even in the limit of strong disorder, in which the average hopping is negligible with respect to its fluctuations, these constants are finite and become of order unity. ### A Gaussian Propagators At second order in $`W`$, the dispersion term $`{\displaystyle \frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^2}}{\displaystyle 𝑑RTr\left(\stackrel{}{}Q^{}\stackrel{}{}Q\right)}{\displaystyle \frac{2\pi \sigma _{xx}}{32}}{\displaystyle 𝑑RTr\left(\stackrel{}{}W\stackrel{}{}W\right)}`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{32}}{\displaystyle 𝑑R\mathrm{\hspace{0.33em}4}Tr\left(\stackrel{}{}B\stackrel{}{}B^{}\right)}+2Tr\left(\stackrel{}{}A\stackrel{}{}A+\stackrel{}{}C\stackrel{}{}C\right),`$ where $`A`$, $`B`$ and $`C`$ are defined through Eqs.(46), (47) and (48). For the $`B`$’s we find the quadratic action $`{\displaystyle \frac{\pi \sigma _{xx}}{2}}{\displaystyle \underset{i=0}{\overset{4}{}}}{\displaystyle \underset{k}{}}{\displaystyle \underset{ab}{}}k^2B_{i,ab}(k)B_{i,ab}(k),`$ so that $$B_{i,ab}(k)B_{j,cd}(k)=\delta _{ij}\delta _{ac}\delta _{bd}D(k),$$ (114) where $$D(k)=\frac{1}{\pi \sigma _{xx}}\frac{1}{k^2}.$$ (115) For the $`A`$’s we have to take into account also the disconnected term: $`{\displaystyle \frac{2\pi \mathrm{\Pi }}{328\mathrm{\Sigma }^4}}{\displaystyle 𝑑RTr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right]Tr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right]}`$ $`{\displaystyle \frac{2\pi \mathrm{\Pi }}{64}}{\displaystyle 𝑑RTr\left[\stackrel{}{}W_3\right]Tr\left[\stackrel{}{}W_3\right]}`$ $`={\displaystyle \frac{2\pi \mathrm{\Pi }}{64}}{\displaystyle 𝑑RTr\left[\stackrel{}{}A+\stackrel{}{}C\right]Tr\left[\stackrel{}{}A+\stackrel{}{}C\right]}`$ $`={\displaystyle \frac{2\pi \mathrm{\Pi }}{16}}{\displaystyle 𝑑RTr\left[\stackrel{}{}A_0+\stackrel{}{}C_0\right]Tr\left[\stackrel{}{}A_0+\stackrel{}{}C_0\right]}.`$ The non vanishing propagators are $`A_{0,ab}(k)A_{0,cd}(k)`$ $`=`$ $`D(k)\left(\delta _{ac}\delta _{bd}+\delta _{ad}\delta _{bc}\right)`$ (116) $``$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}\delta _{ab}\delta _{cd},`$ (117) $`C_{0,ab}(k)C_{0,cd}(k)`$ $`=`$ $`D(k)\left(\delta _{ac}\delta _{bd}+\delta _{ad}\delta _{bc}\right)`$ (118) $``$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}\delta _{ab}\delta _{cd},`$ (119) $`A_{0,ab}(k)C_{0,cd}(k)`$ $`=`$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}\delta _{ab}\delta _{cd},`$ (120) where $`m`$ is the number of replicas, while for $`i=1,2,3`$ $`A_{i,ab}(k)A_{i,cd}(k)`$ $`=`$ $`D(k)\left(\delta _{ac}\delta _{bd}\delta _{ad}\delta _{bc}\right)`$ (121) $`C_{i,ab}(k)C_{i,cd}(k)`$ $`=`$ $`D(k)\left(\delta _{ac}\delta _{bd}\delta _{ad}\delta _{bc}\right).`$ (122) Notice that the particular symmetry of the two-sublattice model leads to additional diffusive modes in the retarded-retarded and advanced-advanced channels, which are not massless in the standard case. ### B Physical Meaning of $`\mathrm{\Pi }`$ Let us introduce an external source which couples to the staggered density of states, which is accomplished by adding to the action a term $`{\displaystyle 𝑑R\overline{\mathrm{\Psi }}_Rs_3\sigma _3\widehat{\lambda }(R)\mathrm{\Psi }_R},`$ where, in the replica space, the source $`\widehat{\lambda }_{\alpha ,\beta }=\lambda _\alpha \delta _{\alpha ,\beta }`$. The fluctuations of the staggered density of states is obtained by the derivative of the partition function with respect, for instance, to $`\lambda _\alpha `$ and $`\lambda _\beta `$, with $`\alpha \beta `$. Inserting the source term in the action, and integrating over the Grassmann fields after introducing the matrix $`Q`$, leads to the following expression of the staggered density of states fluctuation in terms of $`Q`$: $$F(R,R^{})=\frac{1}{\pi ^2\omega _0^2}Tr\left[Q_{\alpha \alpha }(R)s_3\sigma _3\right]Tr\left[Q_{\beta \beta }(R^{})s_3\sigma _3\right].$$ (123) The gaussian estimate of the above correlation function at momentum $`k`$ is given by $$F(k)=16\frac{\mathrm{\Sigma }^2}{\pi ^2\omega _0^2}\left[A_{0,\alpha \alpha }(k)+C_{0,\alpha \alpha }(k)\right]\left[A_{0,\beta \beta }(k)+C_{0,\beta \beta }(k)\right]=\frac{64\mathrm{\Sigma }^2}{\pi ^2\omega _0^2}D(k)\frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}.$$ (124) Therefore, $`\mathrm{\Pi }`$ is directly related to the singular behavior of the staggered density of states fluctuations. ## IX Renormalization Group In this section, we will apply the Wilson–Polyakov Renormalization Group (RG) procedure to analyse the scaling behavior of the action. Indeed, some of the calculations which we present are redumdant, given the proof by Gade and Wegner that the $`\beta `$-function is zero (see Appendix E). Nevertheless, other results besides the conductance $`\beta `$-function are important, so that we describe the whole RG procedure. ### A RG equations In the spirit of Wilson–Polyakov RG approach, we assume that $`T(R)=T_f(R)T_s(R),`$ where $`T_f`$ involves fast modes with momentum $`q[\mathrm{\Lambda }/s,\mathrm{\Lambda }]`$, while $`T_s`$ involves slow modes with $`q[0,\mathrm{\Lambda }/s]`$, being $`\mathrm{\Lambda }`$ the higher momentum cut-off, and the rescaling factor $`s>1`$. Within an $`ϵ`$–expansion, where $`ϵ=d2`$, we define $`{\displaystyle _{\mathrm{\Lambda }/s}^\mathrm{\Lambda }}{\displaystyle \frac{d\stackrel{}{k}}{(2\pi )^d}}D(k)L={\displaystyle \frac{1}{4\pi ^2\sigma _{xx}}}\mathrm{ln}s+\mathrm{O}(ϵ).`$ It is straightforward to show the following result $`Tr\left[\stackrel{}{}Q^{}\stackrel{}{}Q\right]=Tr\left[\stackrel{}{}Q_f^{}\stackrel{}{}Q_f\right]`$ (125) $`+2Tr\left[\stackrel{}{D}_s\sigma _1Q_f\stackrel{}{D}_sQ_f^{}\sigma _1\right]2\mathrm{\Sigma }^2Tr\left[\stackrel{}{D}_s\stackrel{}{D}_s\right]`$ (126) $`+4Tr\left[\stackrel{}{D}_sQ_f^{}\stackrel{}{}Q_f\right],`$ (127) where $`Q_f=\stackrel{~}{T}_f^{}Q_{sp}T_f`$ and $`\stackrel{}{D}_s=T_s\stackrel{}{}T_s^{}`$. Moreover, $`{\displaystyle \frac{1}{\mathrm{\Sigma }^4}}Tr\left[Q^{}\stackrel{}{}Q\sigma _3\right]Tr\left[Q^{}\stackrel{}{}Q\sigma _3\right]`$ (128) $`=Tr\left[\left(\stackrel{}{}W_s+\stackrel{}{}W_f\right)\sigma _3\right]Tr\left[\left(\stackrel{}{}W_s+\stackrel{}{}W_f\right)\sigma _3\right].`$ (129) Since the fast and slow modes live in disconnected regions of momentum space, only the stiffness term (127) generates corrections. By expanding the terms coupling slow and fast modes up to second order in $`W_f`$, the stiffness generates an action term for the slow modes which, after averaging over the fast ones, is $$\frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^2}𝑑RTr\left[\stackrel{}{}Q_s^{}\stackrel{}{}Q_s\right]+F_1_f\frac{1}{2}F_2^2_f,$$ (130) where $`F_1`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^2}}{\displaystyle 𝑑R}2Tr\left[\stackrel{}{D}\sigma _1Q_{sp}W_f\stackrel{}{D}W_fQ_{sp}\sigma _1\right]`$ (132) $`+2Tr\left[\stackrel{}{D}\sigma _1Q_{sp}\stackrel{}{D}Q_{sp}\stackrel{~}{W}_f^2\sigma _1\right],`$ and $$F_2=4\frac{2\pi \sigma _{xx}}{32}𝑑RTr\left[\stackrel{}{D}W_f\stackrel{}{}W_f\right].$$ (133) The explicit evaluation of these terms is outlined in Appendix D. Here we just give the final results. The Kubo conductivity is renormalized according to $$\sigma _{xx}\left(14Lm\right)\sigma _{xx},$$ (134) while the $`\mathrm{\Pi }`$ factor $$\mathrm{\Pi }\mathrm{\Pi }+4L\sigma _{xx}.$$ (135) For what it regards the renormalization of $`E`$ and $`\omega `$, we notice that $`Q`$ $`=`$ $`Q_{sp}T_sT_f^2T_s`$ $``$ $`Q_{sp}T_s^2+Q_{sp}T_sW_fT_s+{\displaystyle \frac{1}{2}}Q_{sp}T_sW_f^2T_s.`$ Since the slow and fast degrees of freedom are defined in different regions of momentum space, only the second term is relevant. By means of Eq.(D28), we find that $$QQ_s\left\{1+\frac{1}{2}\left(28m+\frac{\mathrm{\Pi }}{\sigma _{xx}+m\mathrm{\Pi }}\right)L\right\}$$ (136) This leads to similar corrections to $`E`$ and $`\omega `$, which will have the same scaling behavior. Finally, to describe the cross-over behavior in the presence of symmetry breaking terms, we also need the scaling behavior of the operator $`Tr\left[Q(R)^2\right]`$. We find that $`TrQ^2_f`$ $`=`$ $`Tr\left[Q_{sp}T_sT_f^2T_sQ_{sp}T_sT_f^2T_s\right]_f`$ (137) $`=`$ $`\left[1+2\left(24m+{\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+m\mathrm{\Pi }}}\right)L\right]TrQ_s^2L\left(TrQ_s\right)^2.`$ (138) To implement the RG, we have to rescale the momenta in order to recover the original range $`[0,\mathrm{\Lambda }]`$. This is accomplished by the transformation $`kk/s`$, or, equivalently, $`RRs`$. Therefore, the stiffness as well as the fluctuation terms acquire a scaling factor $`s^ϵ1+ϵ\mathrm{ln}s`$, while the $`E`$ and $`\omega `$ terms a factor $`s^d`$. Hence, after defining $`t=1/(4\pi ^2\sigma _{xx})`$, $`c=1/(4\pi ^2\mathrm{\Pi })`$, and $`\lambda `$ the coupling constant of the operator $`TrQ^2`$, we get the following $`\beta `$-functions $`\beta _t`$ $`=`$ $`ϵt+4mt^2,`$ (139) $`\beta _c`$ $`=`$ $`ϵc4c^2,`$ (140) $`\beta _E`$ $`=`$ $`dE+E{\displaystyle \frac{t}{2}}\left(2+{\displaystyle \frac{t}{c+mt}}\right),`$ (141) $`\beta _\lambda `$ $`=`$ $`d\lambda +2\lambda t\left(2+{\displaystyle \frac{t}{c+mt}}\right).`$ (142) At finite energy, $`E0`$, we may use a two-cutoff scaling approach. Namely, we can follow the previous RG equations up to a cross-over scale, $`s_{cross}=s(E)`$, at which the energy has flowed to a value $`E_0`$ of order $`\mathrm{\Sigma }`$, which plays the role of the high-energy cut-off in the theory. Above this scale, we must neglect all contributions coming from the $`W_3`$ modes, which acquire a mass. That is, we must abandon the RG equations (139)–(142), and let the coupling constants flow in accordance with the standard RG equations, which amount only to a renormalization of $`t`$ according to $$\beta _t=ϵt+t^2.$$ (144) If, by integrating (144), the inverse conductance $`t(s)`$ flows to infinity, signalling an insulating behavior, then we can define a localization length $`\xi _{loc}(E)`$ as the scale at which $`t`$ has grown to a value of order unity. ### B RG in $`d=2`$ In $`d=2`$, the solution of the RG equations for $`t`$, $`c`$, and $`E`$ is $`t(s)`$ $`=`$ $`t(1)=t_0,`$ $`{\displaystyle \frac{1}{c(s)}}`$ $`=`$ $`{\displaystyle \frac{1}{c(1)}}+4\mathrm{ln}s={\displaystyle \frac{1}{c_0}}+4\mathrm{ln}s,`$ $`\mathrm{ln}{\displaystyle \frac{E(s)}{E}}`$ $`=`$ $`\left[2+{\displaystyle \frac{t_0}{2}}\left(2+{\displaystyle \frac{t_0}{c_0}}\right)\right]\mathrm{ln}s+{\displaystyle \frac{1}{2}}t_0^2\mathrm{ln}^2s`$ At finite energy, the crossover length, $`s(E)`$, to the standard, non particle-hole symmetric, model, which, as previously discussed, is defined through $`E[s(E)]=E_0\mathrm{\Sigma }`$, is given by $$s(E)=\mathrm{exp}\left\{\frac{1}{2A}\left(\sqrt{B^2+4A\mathrm{ln}\frac{E_0}{E}}B\right)\right\},$$ (145) being $`A={\displaystyle \frac{1}{2}}t_0^2,B=2+{\displaystyle \frac{t_0}{2}}\left(2+{\displaystyle \frac{t_0}{c_0}}\right).`$ Above $`s(E)`$, $`t`$ flows according to Eq.(144) with $`ϵ=0`$, hence it grows to infinity, implying that the wavefunctions are localized for any $`E0`$. The localization length $`\xi _{loc}(E)s(E)`$, apart from a multiplicative factor which is $`\mathrm{exp}[(1t_0)/t_0]`$ if $`t_01`$. We see that, for $`EE_0\mathrm{exp}\left\{{\displaystyle \frac{1}{2t_0^2}}\left[2+{\displaystyle \frac{t_0}{2}}\left(2+{\displaystyle \frac{t_0}{c_0}}\right)\right]^2\right\},`$ the localization length has a power law behavior, namely $`\xi _{loc}(E)\left({\displaystyle \frac{E_0}{E}}\right)^{\frac{1}{B}},`$ otherwise, at very small energy, it diverges slower, $`\xi _{loc}(E)\mathrm{exp}\sqrt{{\displaystyle \frac{\mathrm{ln}(E_0/E)}{A}}}.`$ The density of states renormalizes like $`\mathrm{ln}{\displaystyle \frac{\rho (s)}{\rho _0}}={\displaystyle \frac{t_0}{2}}\left(2+{\displaystyle \frac{t_0}{c_0}}\right)\mathrm{ln}s+{\displaystyle \frac{1}{2}}t_0^2\mathrm{ln}^2s.`$ At finite energy, the density of states flows until $`s<S(E)`$, after which it stays constant. This implies that the renormalized value is obtained by $`\mathrm{ln}{\displaystyle \frac{\rho (E)}{\rho _0}}=\mathrm{ln}{\displaystyle \frac{\rho (s(E))}{\rho _0}}=\mathrm{ln}{\displaystyle \frac{E_0}{E}}2\mathrm{ln}s(E),`$ leading to $$\rho (E)=\rho _0\left(\frac{E_0}{E}\right)\frac{1}{s(E)^2}\rho _0\left(\frac{E_0}{E}\right)\mathrm{exp}\left(\sqrt{\frac{4\mathrm{ln}(E_0/E)}{A}}\right),$$ (146) the last equality valid at small energy. ### C RG in $`d=3`$ In $`d=3`$, $`t(s)`$ $`=`$ $`t_0s^1,`$ $`c(s)`$ $`=`$ $`{\displaystyle \frac{c_0s^1}{1+4c_04c_0s^1}},`$ $`\mathrm{ln}{\displaystyle \frac{E(s)}{E}}`$ $`=`$ $`3\mathrm{ln}s+{\displaystyle \frac{t_0}{2}}\left(2+{\displaystyle \frac{t_0}{c_0}}+4t_0\right)\left(1{\displaystyle \frac{1}{s}}\right)t_0^2\left(1{\displaystyle \frac{1}{s^2}}\right),`$ hence both $`t`$ and $`c`$ running variables vanish for $`s\mathrm{}`$. The cross-over length diverges, as we approach $`E=0`$, approximately like $$s(E)\left(\frac{E_0}{E}\right)^{\frac{1}{3}}\mathrm{e}^{\frac{t_0}{6}\left(2+\frac{t_0}{c_0}+t_0\right)}.$$ (147) Hence, the density of states $$\rho (E)=\rho _0\left(\frac{E_0}{E}\right)\frac{1}{s(E)^3}\rho _0\mathrm{e}^{\frac{t_0}{2}\left(2+\frac{t_0}{c_0}+t_0\right)},$$ (148) saturates at $`E=0`$ to a value exponentially increased in $`t_0`$ with respect to the bare $`\rho _0`$. At $`E0`$, the inverse conductance above $`s_{cross}=s(E)`$ flows according to (144) with $`ϵ=1`$ and boundary condition $`t(s_{cross})=t_0/s_{cross}`$. We find that, if $`t_0<s_{cross},`$ the system is metallic, otherwise it is insulating, with a localization length $`\xi _{loc}(E){\displaystyle \frac{s(E)}{t_0s(E)}}.`$ This results implies that, for any amount of disorder, sufficiently close to $`E=0`$, all eigenfunctions are delocalized, in agreement with recent numerical results. However, if the disorder is gaussian, as we assumed, the random hopping model with zero regular hopping seems to be characterized by an inverse Drude conductance, $`t_0`$, which is an increasing function of $`\|E\|`$, being smaller than the critical value $`ϵ`$ at $`E=0`$ (see also Ref.). In this case, the presence of a finite mobility edge in $`d=3`$, even for zero regular hopping, would not depend crucially upon the intermediate RG flow in the vicinity to the band center. Nevertheless, we expect that $`t_0`$ at $`E=0`$ varies for different kinds of disorder, and eventually it may become greater than unity. In this case, it is just the vicinity to the band center which makes it possible a finite mobility edge. ## X On site disorder In this and in the following section, we analyse various symmetry breaking terms, which, in the two sublattice representation, contain $`\sigma _0`$ and $`\sigma _3`$, hence spoiling Eq.(1). We start by adding an onsite disorder $`\delta S_{imp}={\displaystyle \underset{R}{}}u_{1,R}c_{1,R}^{}c_{1,R}^{}+u_{2,R}c_{2,R}^{}c_{2,R}^{}={\displaystyle \underset{R}{}}c_R^{}\left({\displaystyle \frac{u_{1,R}+u_{2,R}}{2}}\sigma _0+{\displaystyle \frac{u_{1,R}u_{2,R}}{2}}\sigma _3\right)c_R^{},`$ where $`u_{i,R}=0`$ and $`u_{i,R}u_{j,R^{}}=\delta _{ij}\delta _{RR^{}}v^2`$. Within the path integral, this term becomes, once average over disorder is performed, $`\delta S_{imp}`$ $`=`$ $`{\displaystyle \frac{v^2}{2V}}{\displaystyle \underset{q}{}}Y_{1,q}^{\alpha \beta }Y_{1,q}^{\beta \alpha }+Y_{2,q}^{\alpha \beta }Y_{2,q}^{\beta \alpha }`$ $`=`$ $`{\displaystyle \frac{v^2}{4V}}{\displaystyle \underset{q}{}}Y_{0,q}^{\alpha \beta }Y_{0,q}^{\beta \alpha }+Y_{3,q}^{\alpha \beta }Y_{3,q}^{\beta \alpha }.`$ By adding this term to (24), we get $`S_{imp}+\delta S_{imp}`$ $`=`$ $`{\displaystyle \frac{1}{4V}}{\displaystyle \underset{q}{}}(\omega _q+v^2)Tr\left(Y_{0,q}Y_{0,q}\right)`$ (150) $`(\omega _qv^2)Tr\left(Y_{3,q}Y_{3,q}\right).`$ If we assume that the onsite disorder is weak, i.e. $`\omega _q>v^2`$ at small $`q`$, the consequence is that the $`Q`$ free action becomes $`S_{imp}^0`$ $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{\omega _q+v^2}}Tr\left[Q_{0,q}Q_{0,q}\right]+{\displaystyle \frac{1}{\omega _qv^2}}Tr\left[Q_{3,q}Q_{3,q}\right]`$ (151) $`=`$ $`{\displaystyle \frac{1}{V}}{\displaystyle \underset{q}{}}{\displaystyle \frac{1}{2(\omega _q+v^2)}}Tr\left[Q_qQ_q^{}\right]+{\displaystyle \frac{2v^2}{\omega _q^2v^4}}Tr\left[Q_{3,q}Q_{3,q}\right].`$ (152) Therefore, the on site disorder introduces a mass in the $`Q_3`$ propagators. Specifically, since $`2iQ_3\sigma _3=QQ^{}`$, the mass term can be written as $`{\displaystyle \frac{v^2}{4(\omega _q^2v^4)}}Tr\left[\left(Q_qQ_q^{}\right)\left(Q_qQ_q^{}\right)\right].`$ Close to the saddle point, $`QQ^{}=Q_{sp}^2`$, and, for small $`q`$, we get $`{\displaystyle \frac{v^2}{4(\omega _0^2v^4)}}{\displaystyle 𝑑RTr\left[Q(R)Q(R)+Q(R)^{}Q(R)^{}\right]},`$ which, at second order in $`W`$, reads $`{\displaystyle \frac{v^2\mathrm{\Sigma }^2}{2(\omega _0^2v^4)}}{\displaystyle 𝑑RTr\left[W(R)W(R)+W(R)s_3W(R)s_3\right]}`$ (153) $`={\displaystyle \frac{2v^2\mathrm{\Sigma }^2}{\omega _0^2v^4}}{\displaystyle 𝑑RTr\left[W_3(R)W_3(R)\right]}.`$ (154) In the presence of this term, we could proceed, as before, in the framework of two cutoff scaling theory. That is, we apply the previous RG equations until the above term becomes of the order $`\mathrm{\Sigma }\rho _0`$, i.e. up to the scale which, by Eq.(142), is $$s_{cross}=\mathrm{exp}\left\{\frac{1}{2A}\left(\sqrt{B^2+4A\mathrm{ln}\lambda }B\right)\right\},$$ (155) in $`d=2`$, where $`A=2t_0^2`$ and $`B=2+t_0(2+t_0/c_0)`$, while $`s_{cross}\lambda ^{1/d}`$ in $`d>2`$, where $`\lambda {\displaystyle \frac{\omega _0}{v^2}},`$ in the limit of small $`v`$. Above this scale, the $`W_3`$ propagator gets fully massive, and the inverse conductivity flows with the RG equation (144). In $`d=2`$, this implies that, ultimately, the system gets localized, although the density of states has increased in the first stage of the RG. ## XI Same–Sublattice Regular Hopping We can also introduce a particle–hole symmetry breaking term, by adding to the Hamiltonian a regular term connecting same sublattices, e.g. $`\delta H_{RR^{}}=t_{RR^{}}^{(0)}\sigma _0+t_{RR^{}}^{(3)}\sigma _3\delta H_k=t_k^{(0)}\sigma _0+t_k^{(3)}\sigma _3.`$ Expanding the action in $`\delta H`$, after integrating over the Nambu spinors, we get an additional term $$\delta S[Q]=\frac{1}{2}Tr\left[\stackrel{~}{T}\delta HT^{}G\right].$$ (156) We define $`{\displaystyle \frac{4}{\omega _0}}\lambda _0Q_{sp}`$ $``$ $`i{\displaystyle \frac{1}{V}}{\displaystyle \underset{k}{}}t_k^{(0)}G(k),`$ $`{\displaystyle \frac{4}{\omega _0}}\lambda _3Q_{sp}`$ $``$ $`i{\displaystyle \frac{1}{V}}{\displaystyle \underset{k}{}}t_k^{(3)}G(k),`$ so that (156) becomes $$i\frac{4}{\omega _0}Tr\left[Q\left(\lambda _0\sigma _0+\lambda _3\sigma _3\right)\right].$$ (157) The $`\lambda _0`$-term acts like an energy term. This implies that, if we just shift the chemical potential, we do recover the same scenario as in the absence of this term and at $`E=0`$. On the contrary, the $`\lambda _3`$-term is always a relevant perturbation, whose strength increases under RG iteration as the energy $`E`$. We can define a crossover scale $`s_{cross}`$, which has the same expression as $`s(E)`$ in (145) and (147), for $`d=2`$ and $`d=3`$, respectively, provided $`E\lambda _3`$. Above this scale, the $`W_3`$ modes get fully massive and their contribution to the RG flow drops out. Sometimes $`t_{RR^{}}^{(3)}=0`$, as, for instance, for next-nearest neighbor hopping in a square lattice. In this case, $`\lambda _3=0`$ and we need to evaluate the second order term $`\delta S[Q]={\displaystyle \frac{1}{4}}Tr\left[\stackrel{~}{T}\delta HT^{}G\stackrel{~}{T}\delta HT^{}G\right].`$ If we define $`F(R)=\stackrel{~}{T}(R)T(R)^{}`$, which contains either $`\sigma _0`$ and $`\sigma _3`$, this term can be written, at long wavelengths, as $`\delta S[Q]={\displaystyle \frac{1}{4}}{\displaystyle \underset{k}{}}\left(t_k^{(0)}\right)^2Tr\left[F_qG(k)F_qG(k)\right].`$ By introducing, $`\mathrm{\Sigma }_i^{pq}={\displaystyle \frac{1}{V}}{\displaystyle \underset{k}{}}\left(t_k^{(0)}\right)^2Tr\left[\sigma _iG^p\sigma _iG^q\right],`$ where $`p,q=\pm `$, the following results hold $`\mathrm{\Sigma }_0^{pq}={\displaystyle \frac{1}{2}}(1pq){\displaystyle \frac{4C}{\omega _0}},\mathrm{\Sigma }_3^{pq}={\displaystyle \frac{1}{2}}(1+pq){\displaystyle \frac{4C}{\omega _0}},`$ where $`C`$ is a constant of dimension energy square, with order of magnitude given by the typical value of $`\left(t_k^{(0)}\right)^2`$ close to the surface corresponding to $`E=0`$. Therefore we can write, $`\delta S[Q]`$ $`=`$ $`{\displaystyle \frac{C}{2\omega _0}}{\displaystyle 𝑑RTr\left[F(R)\stackrel{~}{F}(R)F(R)s_3\sigma _1\stackrel{~}{F}(R)s_3\sigma _1\right]}`$ (158) $`=`$ $`\mathrm{const}.+{\displaystyle \frac{C}{2\omega _0\mathrm{\Sigma }^2}}{\displaystyle 𝑑RTr\left[Q(R)^2\right]},`$ (159) which is a mass term for the $`W_3`$ propagators, similar to (154). Therefore, a same-sublattice hopping introduces a cross-over length analogous to (155), with $`\lambda {\displaystyle \frac{\mathrm{\Sigma }^2}{C}}.`$ Above this scale, the contribution of the $`W_3`$ modes to the RG flow has to be dropped out. ## XII Time-reversal symmetry breaking If the random hopping breaks time-reversal symmetry, i.e. $`H_{imp}={\displaystyle \underset{RR^{}}{}}\tau _{RR^{}}^{12}c_R^{}c_R^{}^{}+H.c.,`$ with both real and imaginary part of $`\tau _{RR^{}}^{12}`$ gaussian distributed, after averaging, the impurity action can be written as $$S_{imp}=\frac{1}{V}\underset{q}{}W_qTr\left[X_{1,0,q}X_{2,0,q}+X_{1,3,q}X_{2,3,q}\right],$$ (160) where $`X_{1,0,R}^{\alpha \beta }=\overline{\mathrm{\Psi }}_{1R}^\alpha \tau _0\mathrm{\Psi }_{1R}^\beta ,X_{1,3,R}^{\alpha \beta }=\overline{\mathrm{\Psi }}_{1R}^\alpha \tau _3\mathrm{\Psi }_{1R}^\beta ,`$ with the indices $`\alpha `$ and $`\beta `$ running only over the replicas and the advanced/retarded components. This implies that the manifold in which $`Q`$ varies contains in this case only $`\tau _0`$ and $`\tau _3`$ components. Indeed, as in the time reversal invariant case we are able to parametrize the $`8m\times 8m`$ matrix $`T`$ in terms of a $`4m\times 4m`$ matrix $`U\mathrm{U}(4m)/\mathrm{Sp}(2m)`$ \[see Eq.(37)\], similarly, without time-reversal symmetry, $`T`$ can be parametrized by means of a $`2m\times 2m`$ matrix $`U\mathrm{U}(2m)`$, in agreement with Gade. The effective non-linear $`\sigma `$-model is not modified, but the expressions (D1), (D2), (D4), (D6), and (D7) have to be substituted by $`B_{ab}P_{bc}B_{cd}^{}`$ $`=`$ $`2D(k)\delta _{ad}Tr\left(P_0\right),`$ (161) $`B_{ab}P_{bc}B_{cd}`$ $`=`$ $`0,`$ (162) $`A_{ab}P_{bc}A_{cd}`$ $`=`$ $`2D(K)\delta _{ad}Tr\left(P_0\right)D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}P_{ad},`$ (163) $`C_{ab}P_{bc}C_{cd}`$ $`=`$ $`2D(K)\delta _{ad}Tr\left(P_0\right)D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}P_{ad},`$ (164) $`A_{ab}P_{bc}C_{cd}`$ $`=`$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}P_{ad}.`$ (165) where $`P=P_0+iP_3\tau _3`$. Hence, the RG equations at $`m=0`$ are, in this case, $`\beta _t`$ $`=`$ $`ϵt,`$ (166) $`\beta _c`$ $`=`$ $`ϵc2c^2,`$ (167) $`\beta _E`$ $`=`$ $`dE+{\displaystyle \frac{Et^2}{2c}},`$ (168) which coincide with those obtained by Gade. ## XIII Discussion and comparison with the standard localization theory In this section, we summarize the main differences between the model (3) and the standard non-linear $`\sigma `$-model which is derived in the theory of Anderson localization , placing particular emphasis on the properties of the $`Q`$-matrix. The specific form of the off-diagonal disorder we consider, which only couples one sublattice to the other, leads, via the Hubbard-Stratonovich decoupling, to the introduction of a space-varying $`8m\times 8m`$ complex $`Q`$-matrix, $`Q=Q_0\sigma _0+iQ_3\sigma _3`$. Here, $`Q_0`$ and $`Q_3`$ are $`4m\times 4m`$ hermitean matrices, of which matricial structure refers to the retarded/advanced, spinor particle/hole and $`m`$ replica components. Contrary to the standard case, $`Q`$ is not hermitean. The evaluation of the saddle point, $`Q_{sp}=\sigma _0\tau _0s_3`$ (sectionV), as well as the derivation of the efective action (section VI) are analogous to the standard case . (We recall that $`\sigma _i`$, $`s_i`$ and $`\tau _i`$ indicate the Pauli matrices, including the unit matrix, acting on sublattice, advanced/retarded and spinor components, respectively.) The non-linear $`\sigma `$-model, Eq.(111), is obtained by integrating out the longitudinal massive $`Q`$-fluctuations and only keeping the transverse soft modes. The real novelty with respect to localization theory is not in the structure of the effective action. Indeed, the new term in (111), namely $`\mathrm{\Pi }{\displaystyle 𝑑R\left[Tr\left(Q(R)^{}\stackrel{}{}Q(R)\sigma _3\right)\right]^2},`$ even if present, would be irrelevant in the standard case. On the contrary, the essential difference, as expected, lies in the ensamble spanned by the soft modes at the particle-hole symmetry point $`E=0`$. We get $`Q_{Soft}=\stackrel{~}{T}^1Q_{sp}T`$, where the unitary matrix $`T`$ only contains $`\sigma _0`$ and $`\sigma _3`$, $`T=\mathrm{exp}\left[{\displaystyle \frac{W_0\sigma _0+W_3\sigma _3}{2}}\right],`$ and $`\stackrel{~}{T}=\sigma _1T\sigma _1=\sigma _2T\sigma _2=\mathrm{exp}\left[{\displaystyle \frac{W_0\sigma _0W_3\sigma _3}{2}}\right].`$ These expressions derive by the conditions (1), which fully specify the model, as shown in section IV. In that section, we also showed that the ensamble can be expressed in terms of unitary $`4m\times 4m`$ matrices $`U=\mathrm{exp}\left[{\displaystyle \frac{W_0+W_3}{2}}\right],`$ as argued by Gade and Wegner. Selecting the subset which leaves the saddle point invariant gives Eq.(3) with $`U\mathrm{U}(4m)/\mathrm{Sp}(2m)`$. In terms of $`T`$, the condition $`\stackrel{~}{T}^1Q_{sp}TQ_{sp}`$, leads to the requirements $`[W_0,s_3]0`$ and $`\{W_3,s_3\}0`$, which implies that $`W_0`$ is off-diagonal in the energy retarded/advanced space (as in the standard localization theory), while $`W_3`$ is diagonal. In other words, the omogeneous and staggered modes, $`W_0`$ and $`W_3`$, respectively, have different structure in the energy space. The energy diagonal $`W_3`$-modes betray the presence, at $`E=0`$, of diffusive poles in the disorder averaged products of retarded and advanced Green’s functions, $`\overline{G_RG_R}`$ and $`\overline{G_AG_A}`$, with $`G_{R,A}=\left(H\pm i0^+\right)^1`$. In the localization theory , only the mixed products $`\overline{G_RG_A}`$ have a singular behavior. This explains why singular corrections to the density of states (which involves connected diagrams with same energy Green’s functions) are present in the two-sublattice model, while they are absent in the localization theory. In a square lattice, the energy diagonal modes have the transparent meaning of density fluctations with wave-vector $`q`$ nearby the nesting vector $`G=(\pi ,\pi ,\mathrm{})`$, see Appendix A. Indeed $`Q_3(q)`$ $``$ $`{\displaystyle \underset{RA}{}}\mathrm{e}^{iqR}\left(\mathrm{\Psi }_{1R}\overline{\mathrm{\Psi }}_{1R}\mathrm{\Psi }_{2R}\overline{\mathrm{\Psi }}_{2R}\right)`$ $`=`$ $`{\displaystyle \underset{RA,B}{}}\mathrm{e}^{i(q+G)R}\mathrm{\Psi }_R\overline{\mathrm{\Psi }}_R=Q(q+G),`$ where $`A`$ and $`B`$ label the two sublattice, and, for $`RA`$, we have taken by definition $`\mathrm{\Psi }_{1R}=\mathrm{\Psi }_R`$ and $`\mathrm{\Psi }_{2R}=\mathrm{\Psi }_{R+a\widehat{x}}`$, being $`a`$ the lattice spacing and $`\widehat{x}`$ the unit vector in the $`x`$-direction. As soon as $`E0`$, nesting is not more important and indeed $`Q_3`$ becomes massive. Finally, because of $`\overline{G_RG_R}`$ and $`\overline{G_AG_A}`$, also the conductance acquires other corrections with respect to standard localization theory. Indeed, these corrections add to give a vanishing $`\beta `$-function for $`\sigma _{xx}`$, as first indicated by Gade and Wegner . We have explained in the Introduction section (see also Appendix E) that this is a consequence of a simple abelian gauge symmetry generated by Eq.(1) at the particle-hole symmetry point $`E=0`$. Similarly to the results for the density of states, this behavior of the $`\beta `$-function is at odds with the standard theory. ## XIV Conclusions In this work, we have derived the effective non-linear $`\sigma `$-model of a disordered electronic system on a generic bipartite lattice. This model, if the hopping matrix elements as well as the disorder only couple one sublattice with the other, shows an interesting behavior close to the band center, i.e. to the particle-hole symmetry point. Namely, the wave-functions are always delocalized at the band center, in any dimension. By a Renormalization Group (RG) analysis, in the framework of an $`ϵ`$-expansion, $`ϵ=d2`$, we have found that the quantum corrections to the conductivity vanish if the chemical potential is exactly at the band center, thus implying a metallic behavior. In two dimensions, in particular, the Kubo conductivity flows to a fixed value by iterating the RG. On the contrary, we have found that the staggered density of states fluctuations, which are controlled by a new parameter in the non-linear $`\sigma `$-model, are singular. This result is reminiscent of what it is found in equivalent one-dimensional models. In fact, models of disordered spinless fermions in one-dimension can be mapped, by a Jordan-Wigner transformation, onto disordered spin chains. In many cases, it is known that, in spite of the presence of disorder, these spin chain models display critical behavior, as shown in great detail by D. Fisher for random Heisenberg antiferromagnets and random transverse-field Ising chains. Indeed, as pointed out by Fisher, the staggered spin fluctuations, in a random antiferromagnetic chain, also display critical behavior in the form of a power law decay, $`(1)^R\overline{S(R)S(0)}R^2`$, where the bar indicates impurity average. Since the staggered spin-density corresponds to the staggered density of the spinless fermions, this result is consistent with the outcome of our analysis, which further suggests that a similar scenario generally holds in such models. Moreover, as in one-dimension, we find that the density of states is strongly modified by the disorder at the band center, and it actually diverges in $`d=2`$. In reality, a random Heisenberg chain, away from the $`XXZ`$ limit, maps onto a spinless fermion random hopping model in the presence of a random nearest-neighbor interaction. However, even in the presence of this additional interaction, the Hamiltonian has still the abelian gauge-like symmetry described in Appendix E, which is at the origin of the delocalization of the band center state. This observation is also compatible with Fisher’s result that the physical behavior does not qualitatively change upon moving away from the XXZ limit towards the isotropic XXX Heisenberg point. Many of the results which we have derived were already known. The existence of delocalized states at the band center of a two-sublattice model was argued already in 1979 by Wegner. The effective non-linear $`\sigma `$-model when the disorder breaks time-reversal invariance, as well as the RG equations, have earlier been derived by Gade, although in a particular two-sublattice model. The extension to disordered systems with time-reversal symmetry was later on argued by Gade and Wegner. Finally, random flux models and disordered Dirac fermion models have recently been the subject of an intensive theoretical study, for their implications to a variety of different physical problems. In spite of that, our analysis has several novelties with respect to earlier studies. First of all, the two-sublattice model which we study is quite general. Secondly, the physical interpretation of the parameters which appear in the non-linear $`\sigma `$-model is quite transparent. Thirdly, the explicit derivation of the RG equations with time-reversal invariance is presented. ## XV Acknowledgments We are grateful to A. Nersesyan, V. Kratsov, E. Tosatti, and Yu Lu for helpful discussions and comments. ## A Specific examples As an example, we consider a tight binding model with only nearest neighbor hopping on a square and honeycomb lattice. In the case of square lattice, the enlarged unit cell is the $`\sqrt{2}\times \sqrt{2}`$ one. The new reciprocal lattice vectors are $`\stackrel{}{G}_1=2\pi (1,1)/a`$, $`\stackrel{}{G}_2=2\pi (1,1)/a`$, and the angle $`\theta _k`$ of Eq.(11) is $$\theta _k=\frac{a}{2}(k_1+k_2)=k_xa.$$ (A1) In the case of the honeycomb lattice, the unit cell contains already two lattice sites. The energy is given by $`ϵ_k^2`$ $`=`$ $`\left[1+2\mathrm{cos}\left({\displaystyle \frac{3}{2}}k_xa\right)\mathrm{cos}\left({\displaystyle \frac{\sqrt{3}}{2}}k_ya\right)\right]^2`$ $`+`$ $`\left[2\mathrm{sin}\left({\displaystyle \frac{3}{2}}k_xa\right)\mathrm{cos}\left({\displaystyle \frac{\sqrt{3}}{2}}k_ya\right)\right]^2,`$ and $`\theta _k=\mathrm{tan}^1\left({\displaystyle \frac{2\mathrm{sin}\left(\frac{3}{2}k_xa\right)\mathrm{cos}\left(\frac{\sqrt{3}}{2}k_ya\right)}{1+2\mathrm{cos}\left(\frac{3}{2}k_xa\right)\mathrm{cos}\left(\frac{\sqrt{3}}{2}k_ya\right)}}\right).`$ The Brillouin zone is still honeycomb, with the $`y`$-axis one of its axes, and side equal to $`4\pi /(3\sqrt{3}a)`$. ## B Ward Identity Let us consider a generic Hamiltonian in the two sublattice representation $$H=\underset{R_1,R_2}{}c_{R_1}^{}H_{R_1,R_2}c_{R_2}^{},$$ (B1) where $`H_{R_1,R_2}`$ is a $`2\times 2`$ Hermitean matrix. The current operator $`\stackrel{}{J}(R)={\displaystyle \underset{R_1,R_2}{}}c_{R_1}^{}\stackrel{}{J}_{R_1,R_2}(R)c_{R_2}^{},`$ can be obtained by the continuity equation, leading to $$\stackrel{}{}\stackrel{}{J}(R)=i\underset{R_1}{}c_R^{}H_{R,R_1}c_{R_1}^{}c_{R_1}^{}H_{R_1,R}c_R^{},$$ (B2) being $`\stackrel{}{}`$ the discrete version of the differential operator. The long-wavelength expression for $`\stackrel{}{J}(R)`$ can be obtained by Fourier transformation, namely, through $`i{\displaystyle \underset{R}{}}\stackrel{}{q}\stackrel{}{J}(R)\mathrm{e}^{i\stackrel{}{q}\stackrel{}{R}}=i{\displaystyle \underset{R,R_1}{}}\left(c_R^{}H_{R,R_1}c_{R_1}^{}c_{R_1}^{}H_{R_1,R}c_R^{}\right)\mathrm{e}^{i\stackrel{}{q}\stackrel{}{R}},`$ and expanding both sides in $`q`$, we get, for the linear term, $`i{\displaystyle \underset{R}{}}\stackrel{}{q}\stackrel{}{J}(R)={\displaystyle \underset{R,R_1}{}}\stackrel{}{q}\stackrel{}{R}\left(c_R^{}H_{R,R_1}c_{R_1}^{}c_{R_1}^{}H_{R_1,R}c_R^{}\right)`$ $`={\displaystyle \underset{R,R_1}{}}\stackrel{}{q}\left(\stackrel{}{R}_1\stackrel{}{R}\right)c_{R_1}^{}H_{R_1,R}c_R^{},`$ hence $$\stackrel{}{J}(R)=i\underset{R_1}{}\left(\stackrel{}{R}_1\stackrel{}{R}\right)c_{R_1}^{}H_{R_1,R}c_R^{}.$$ (B3) Let us define the correlation functions $$\chi _{\mu ,i}(R,R^{};t,t_1,t_2)=T\left[c_{R_1}^{}(t)J_{R_1,R_2}^\mu (R)c_{R_2}^{}(t)c_{R_3}^{}(t_1)J_{R_3,R_4}^i(R^{})c_{R_4}^{}(t_2)\right],$$ (B4) where $`\mu =0,1,2,3`$, $`J_{R_1,R_2}^0(R)=\delta _{RR_1}\delta _{RR_2}`$ are the density matrix elements, and $`J^i`$, for $`i=1,2,3`$, are the matrix element components of the current. By the continuity equation, we find that $`_t\chi _{0,i}+_j\chi _{j,i}`$ $`=`$ $`i{\displaystyle \underset{R_1}{}}\delta (tt_1)Tr\left[G(R_1,R;t_2t)J_{R,R_1}^i(R^{})\right]`$ $`+\delta (tt_2)Tr\left[G(R,R_1;tt_1)J_{R_1,R}^i(R^{})\right].`$ If we integrate both sides by $`{\displaystyle 𝑑t𝑑t_1𝑑t_2},\mathrm{e}^{i(E+\omega )(tt_1)}\mathrm{e}^{iE(t_2t)},`$ at $`\omega =0`$ we find $$_j\chi _{j,i}(R,R^{};E)=i\underset{R_1}{}Tr\left[G(R,R_1;E)J_{R_1,R}^i(R^{})\right]Tr\left[G(R_1,R;E)J_{R,R_1}^i(R^{})\right].$$ (B5) Using once more the continuity equation (B2), we find $`_j_i^{}\chi _{j,i}(R,R^{};E)=Tr\left[G(R,R^{};E)H_{R^{},R}+G(R^{},R;E)H_{R,R^{}}\right]`$ $`+\delta _{RR^{}}{\displaystyle \underset{R_1}{}}Tr\left[G(R,R_1;E)H_{R_1,R}\right]+Tr\left[G(R_1,R;E)H_{R,R_1}\right].`$ By Fourier transform, $`{\displaystyle \underset{RR^{}}{}}_j_i^{}\chi _{j,i}(R,R^{};E)\mathrm{e}^{i\stackrel{}{q}(\stackrel{}{R}\stackrel{}{R^{}})}={\displaystyle \underset{RR^{}}{}}q_iq_j\chi _{j,i}(R,R^{};E)\mathrm{e}^{i\stackrel{}{q}(\stackrel{}{R}\stackrel{}{R^{}})}`$ $`{\displaystyle \underset{RR^{}}{}}\left(1\mathrm{e}^{i\stackrel{}{q}(\stackrel{}{R}\stackrel{}{R^{}})}\right)Tr\left[G(R,R^{};E)H_{R^{},R}+G(R^{},R;E)H_{R,R^{}}\right].`$ At small $`q`$, the above expression is $`{\displaystyle \underset{RR^{}}{}}q_iq_j\chi _{j,i}(R,R^{};E)`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{RR^{}}{}}q_iq_j\left(R_iR_i^{}\right)\left(R_jR_j^{}\right)Tr\left[G(R,R^{};E)H_{R^{},R}+G(R^{},R;E)H_{R,R^{}}\right],`$ leading to $$\underset{RR^{}}{}\chi _{j,i}(R,R^{};E)=\underset{RR^{}}{}\left(R_iR_i^{}\right)\left(R_jR_j^{}\right)Tr\left[G(R,R^{};E)H_{R^{},R}\right].$$ (B6) ## C Longitudinal Modes As discussed in Section VII, the expression of the $`Q`$-matrix which includes also the longitudinal modes is $`Q_P(R)=Q(R)+S(R)`$, where $`Q(R)`$ and $`S(R)`$ have been defined through (101). The free action for these fields contains a local term, Eq.(105), and a non local one, Eq.(105). The latter is $`{\displaystyle \frac{1}{4}}{\displaystyle }dRdR^{}\mathrm{\Gamma }(RR^{})Tr[\mathrm{\Delta }_RQ(R^{})\mathrm{\Delta }_RQ(R^{})^{}`$ (C1) $`+\mathrm{\Delta }_RS(R^{})\mathrm{\Delta }_RS(R^{})^{}+2\mathrm{\Delta }_RQ(R^{})\mathrm{\Delta }_RS(R^{})^{}],`$ (C2) where we have defined the operator $`\mathrm{\Delta }_Rf(R^{})=f(R)f(R^{})={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{n!}}\left(\stackrel{}{R}\stackrel{}{R^{}}\right)^n\stackrel{}{}^nf(R^{}).`$ Let us apply this operator to $`Q(R)`$ and $`S(R)`$, keeping all terms which contains at most two derivatives which act to the transverse matrices $`T`$. We obtain $`\mathrm{\Delta }_RQ(R^{})`$ $``$ $`\stackrel{}{}Q(R^{})\left(\stackrel{}{R}\stackrel{}{R^{}}\right),`$ (C3) $`\mathrm{\Delta }_RS(R^{})`$ $``$ $`\stackrel{~}{T}(R^{})^{}\mathrm{\Delta }_RP(R^{})T(R^{})`$ (C4) $`+`$ $`\left[\left(\stackrel{}{}\stackrel{~}{T}(R^{})^{}\right)P(R)T(R^{})+\stackrel{~}{T}(R^{})^{}P(R)\left(\stackrel{}{}T(R^{})\right)\right]\left(\stackrel{}{R}\stackrel{}{R^{}}\right)`$ (C5) $`+`$ $`{\displaystyle \frac{1}{2}}[\left(\stackrel{}{}^2\stackrel{~}{T}(R^{})^{}\right)P(R)T(R^{})+2\left(\stackrel{}{}\stackrel{~}{T}(R^{})^{}\right)P(R)\left(\stackrel{}{}T(R^{})\right)`$ (C6) $`+`$ $`\stackrel{~}{T}(R^{})^{}P(R)\left(\stackrel{}{}^2T(R^{})\right)](\stackrel{}{R}\stackrel{}{R^{}})^2.`$ (C7) The term which is obtained by (C4) times its hermitean conjugate, together with the local piece (105) give the free action of the longitudinal modes $$S_0[P]=\frac{1}{V}\underset{q}{}\frac{1}{2\omega _q}Tr\left[P_qP_q^{}\right].$$ (C8) The mixed terms, after defining $`\stackrel{}{D}=T\stackrel{}{}T^{}`$, give rise to the coupling between transverse and longitudinal modes $`S[Q,P]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle 𝑑R𝑑R^{}\mathrm{\Gamma }(RR^{})\left(\stackrel{}{R}\stackrel{}{R^{}}\right)Tr\left[\stackrel{}{D}(R^{})\left(P^{}(R^{})P(R)P(R)^{}P(R^{})\right)\right]}`$ (C9) $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle }dRdR^{}\mathrm{\Gamma }(RR^{})(\stackrel{}{R}\stackrel{}{R^{}})^2`$ (C12) $`Tr[\sigma _1\stackrel{}{D}(R^{})\sigma _1P(R)\stackrel{}{D}(R^{})P^{}(R^{})+\sigma _1\stackrel{}{D}(R^{})\sigma _1P(R^{})\stackrel{}{D}(R^{})P^{}(R)`$ $`\stackrel{}{D}(R^{})\stackrel{}{D}(R^{})(P^{}(R)P(R^{})+P^{}(R^{})P(R))]`$ $``$ $`{\displaystyle \frac{1}{4}}{\displaystyle }dRdR^{}\mathrm{\Gamma }(RR^{})(\stackrel{}{R}\stackrel{}{R^{}})Tr[\stackrel{}{}Q(R^{})\mathrm{\Delta }_RP(R^{})^{}+H.c.].`$ (C13) The last term, (C13), gives rise to higher gradient contributions, hence can be neglected. ### 1 Longitudinal propagators Before averaging over the longitudinal modes, we have to evaluate the longitudinal propagators. The matrix $`P=\left(P_{0,0}s_0+P_{0,3}s_3\right)\sigma _0+i\left(P_{3,1}s_1+P_{3,2}s_2\right)\sigma _3,`$ has to satisfy $`cP^tc^t=P`$, and, in addition, all $`P_\alpha =P_\alpha ^{}`$. For $`\alpha =(0,0),(0,3),(3,1)`$, by writing $`P_\alpha =P_\alpha ^{(0)}\tau _0+i\stackrel{}{P}_\alpha \stackrel{}{\tau },`$ we find that $`P_\alpha ^{(0)},\stackrel{}{P}_\alpha e,P_\alpha ^{(0)}=\left(P_\alpha ^{(0)}\right)^t,\stackrel{}{P}_\alpha =\left(\stackrel{}{P}_\alpha \right)^t.`$ For $`\alpha =(3,2)`$, by writing $`P_{3,2}=iP_{3,2}^{(0)}\tau _0+\stackrel{}{P}_{3,2}\stackrel{}{\tau },`$ we must impose $`P_{3,2}^{(0)},\stackrel{}{P}_{3,2}e,P_{3,2}^{(0)}=\left(P_{3,2}^{(0)}\right)^t,\stackrel{}{P}_{3,2}=\left(\stackrel{}{P}_{3,2}\right)^t.`$ If $`P_\alpha ^{(i)}`$ is a symmetric real matrix, its propagator is $$P_{\alpha ,ab}^{(i)}P_{\alpha ,cd}^{(i)}=\frac{G}{2}\left(\delta _{ad}\delta _{bc}+\delta _{ac}\delta _{bd}\right),$$ (C14) while if it is antisymmetric $$P_{\alpha ,ab}^{(i)}P_{\alpha ,cd}^{(i)}=\frac{G}{2}\left(\delta _{ad}\delta _{bc}\delta _{ac}\delta _{bd}\right),$$ (C15) where $`G_k=\omega _k/8V`$, and $`(a,b,c,d)`$ are replica indices. By means of these propagators, we readily find that, if $`M=M_i^{(0)}+i\stackrel{}{M}_i\stackrel{}{\tau }`$, where $`M^{(i)}`$ are matrices in the replica space, then, for $`\alpha =(0,0),(0,3),(3,1)`$, the following results hold $$P_\alpha MP_\alpha =GcM^tc^t+2GTr\left(M^{(0)}\right)=GcM^tc^t+GTr\left(M\right),$$ (C16) while, for $`\alpha =(3,2)`$, $$P_{3,2}MP_{3,2}=GcM^tc^t+2GTr\left(M^{(0)}\right)=GcM^tc^t+GTr\left(M\right).$$ (C17) More generally, $`PMs_i\sigma _jP^{}`$ $`=`$ $`GcM^tc^t\left(s_i\sigma _j+s_3s_i\sigma _js_3+\sigma _3s_1s_i\sigma _js_1\sigma _3\sigma _3s_2s_i\sigma _js_2\sigma _3\right)`$ (C18) $`+`$ $`GTr\left(M\right)\left(s_i\sigma _j+s_3s_i\sigma _js_3+\sigma _3s_1s_i\sigma _js_1\sigma _3+\sigma _3s_2s_i\sigma _js_2\sigma _3\right),`$ (C19) $`PMs_i\sigma _jP`$ $`=`$ $`GcM^tc^t\left(s_i\sigma _j+s_3s_i\sigma _js_3\sigma _3s_1s_i\sigma _js_1\sigma _3+\sigma _3s_2s_i\sigma _js_2\sigma _3\right)`$ (C20) $`+`$ $`GTr\left(M\right)\left(s_i\sigma _j+s_3s_i\sigma _js_3\sigma _3s_1s_i\sigma _js_1\sigma _3\sigma _3s_2s_i\sigma _js_2\sigma _3\right).`$ (C21) For $`j=0,3`$ the above expression simplifies to $`PMs_i\sigma _jP^{}`$ $`=`$ $`2Gc\left(Ms_i\sigma _j\right)^tc^t+2G\sigma _jTr\left(Ms_is_0\right),`$ (C22) $`PMs_i\sigma _jP`$ $`=`$ $`2Gs_3c\left(Ms_i\sigma _j\right)^tc^ts_3+2G\sigma _jTr\left(Ms_is_3\right),`$ (C23) while for $`j=1,2`$ $`PMs_i\sigma _jP^{}`$ $`=`$ $`2Gs_3c\left(Ms_i\sigma _j\right)^tc^ts_3+2G\sigma _jTr\left(Ms_is_3\right),`$ (C24) $`PMs_i\sigma _jP`$ $`=`$ $`2Gc\left(Ms_i\sigma _j\right)^tc^t+2G\sigma _jTr\left(Ms_is_0\right).`$ (C25) ### 2 Averaging $`S[Q,P]`$ We have now all what it is needed to proceed in the averaging over $`P`$. Here we just sketch the calculation, which is quite involved and requires the matrix properties of $`\stackrel{}{D}`$ which are determined in Appendix D. We just remark that (C12) and (C12) do not reproduce the correct stiffness term. Indeed, it is (C9), which contributes at second order, which cancels the additional terms and allows to express everything in terms of the matrix $`Q`$. By means of the previously calculated propagators of the longitudinal modes, we find that $$S[Q,P]_P=\frac{Y}{4d\mathrm{\Sigma }^2}𝑑RTr\left[\stackrel{}{}Q(R)^{}\stackrel{}{}Q(R)\right]+\frac{1}{8\mathrm{\Sigma }^2}\left[Tr\left(Q(R)^{}\stackrel{}{}Q(R)\sigma _3\right)\right]^2,$$ (C26) where $$Y=𝑑R\mathrm{\Gamma }(R)\mathrm{\Gamma }(R)^1R^2.$$ (C27) We notice that the first term is a contribution to the Kubo conductivity and the second to the staggered density of states fluctuations of the diagrams where the current vertices are those proportional to the random hopping. ### 3 Additional terms The last class of corrections which generate new operators is obtained by expanding each Green’s function in (65) at first order, $`G_P=G_0iG_0PG_0`$, leading to the term $$\frac{1}{4}Tr\left(GPGUGPGU\right)_P.$$ (C28) For the sake of clarity, we will analyse this term only in the case of a cubic lattice, where the derivation is more straightforward. We will postpone a discussion about the general case at the end of the section. In the cubic lattice, according to Eq.(33), the electron-$`Q`$ coupling can be brought to a local one also in the diagonal basis. In this basis, $`Q=Q_0+iQ_1\sigma _1`$, $$P(R)=\left(P_{00}s_0+P_{03}s_3\right)\sigma _0+i\left(P_{11}s_1+P_{12}s_2\right)\sigma _1,$$ (C29) the unitary matrix $$T(R)=\mathrm{exp}\left[\frac{W_0(R)}{2}\sigma _0+\frac{W_1(R)}{2}\sigma _1\right],$$ (C30) where $`W_1`$ has the same form as $`W_3`$ defined by Eq.(45), and $`\stackrel{~}{T}=\sigma _3T\sigma _3`$. Moreover, since $`H_{RR^{}}=ϵ_{RR^{}}\sigma _3`$ in the diagonal basis, being $`ϵ_{RR^{}}`$ the Fourier transfom of $`ϵ_k`$, the current operator which appears in the definition of $`U_{RR^{}}`$, Eq.(56), has only matrix elements $`\stackrel{}{J}_k=\stackrel{}{}ϵ_k\sigma _3`$. Once averaged over $`P`$, (C28) gives, among other terms which correct the Kubo conductivity, a new term \[see Eq.(Greenfunctiond)\] $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a,b=1,2}{}}{\displaystyle \underset{h=\pm }{}}{\displaystyle \underset{pkq}{}}\omega _{pk}Tr\left(U_q^{b,h;a,h}𝒢_{p+q}^{a,h}𝒢_p^{b,h}\right)Tr\left(U_q^{a,h;b,h}𝒢_k^{b,h}𝒢_{k+q}^{a,h}\right)`$ $`Tr\left(U^{b,h;a,h}𝒢_{p+q}^{a,h}𝒢_p^{b,h}\right)Tr\left(U^{a,h;b,h}𝒢_k^{b,h}𝒢_{k+q}^{a,h}\right),`$ where the first piece derives from $`P_0`$ and the second from $`P_1`$. The structure in the energy/sublattice indices can be shortly represented by $`{\displaystyle \underset{i=1,\mathrm{},4}{}}\mathrm{\Lambda }_i\mathrm{\Lambda }_i={\displaystyle \frac{1}{4}}\left(\sigma _2\sigma _2+s_3\sigma _1s_3\sigma _1+\sigma _3\sigma _3+s_3s_3\right),`$ so that the above term can be written, at small $`q`$, as $$\frac{1}{2}\underset{i=1,\mathrm{},4}{}\underset{pkq}{}\omega _{pk}Tr\left(U_q𝒢_p\mathrm{\Lambda }_i𝒢_p\right)Tr\left(U_q𝒢_k\mathrm{\Lambda }_i𝒢_k\right)$$ (C31) We remind that $`𝒢_pU_q𝒢_p=i𝒢_p\stackrel{}{D}_q\stackrel{}{J}_p𝒢_p,`$ where $`\stackrel{}{D}_q`$ is the Fourier transform of $`\stackrel{~}{T}(R)\stackrel{}{}\stackrel{~}{T}(R)^1`$. We notice that only the diagonal component in energy of $`\stackrel{}{D}`$ enters. Moreover, for the diagonal matrices $`s_j=s_0,s_3`$, since $`\stackrel{}{D}=\stackrel{}{D}_0\sigma _0+\stackrel{}{D}_1\sigma _1`$, it derives that, for $`i=0,1`$, the following equality holds $`Tr\stackrel{}{D}s_j\sigma _i=Tr\left(\stackrel{}{}\stackrel{~}{W}s_j\sigma _i/2\right)`$. Therefore, just $`W_1\sigma _1`$ contributes. By means of (78), for any $`\mathrm{\Lambda }`$’s we have $`iTr\left(\stackrel{}{D}_{1,q}\sigma _1\stackrel{}{J}_p𝒢(p)\mathrm{\Lambda }_i𝒢(p)\right)=+\stackrel{}{}ϵ_pTr\left(\stackrel{}{D}_{1,q}s_1𝒢(p)s_1\sigma _2\mathrm{\Lambda }_i𝒢(p)\right)`$ (C32) $`=\stackrel{}{}ϵ_pTr\left[\stackrel{}{D}_{1,q}\left(G_3^2\sigma _3\sigma _2\mathrm{\Lambda }_i\sigma _3G_0^2s_3\sigma _2\mathrm{\Lambda }_is_3+G_3G_0s_3[\sigma _3,\sigma _2\mathrm{\Lambda }_i]\right)\right].`$ (C33) Only for $`\mathrm{\Lambda }_1=\sigma _2/4`$ the trace over the $`\sigma `$’s is finite, leading to $$2\left(G_3^2G_0^2\right)\stackrel{}{}ϵ_pTr\left(\stackrel{}{D}_{1,q}\right).$$ (C34) In conclusion, going back to the original sublattice representation, and defining $$\mathrm{\Pi }^{(0)}=\frac{1}{4\pi V^2d}\underset{pk}{}\stackrel{}{}ϵ_k\stackrel{}{}ϵ_p\omega _{pk}Tr\left(\sigma _3G_p^+\sigma _3G_p^+\right)Tr\left(\sigma _3G_k^+\sigma _3G_k^+\right),$$ (C35) we obtain the following explicit expression of (C28) \[notice that $`Q=Q_0\sigma _0+iQ_1\sigma _1Q_0\sigma _0+iQ_3\sigma _3`$ in the sublattice basis\] $$\frac{2\pi }{832\mathrm{\Sigma }^4}\mathrm{\Pi }^{(0)}𝑑RTr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right]Tr\left[Q^{}(R)\stackrel{}{}Q(R)\sigma _3\right].$$ (C36) This term, which derives from the regular current vertices, has the same form as the second term in (C26), which, on the contrary, is due to the random current vertices. Therefore, the coupling constant $`\mathrm{\Pi }`$ which appears in the final expression (108) is the sum of both terms. The symmetry of the operator (C36), which, due to the trace, involves the Nambu, energy and sublattice components $`\tau _0`$, $`s_0`$ and $`\sigma _3`$, respectively, suggests that the prefactor $`\mathrm{\Pi }`$ represents the fluctuations of the staggered density of states, as discussed in section VIII B. In the case of a generic bipartite lattice, we do expect an analogous term to appear, because the staggered-density of states fluctuations still acquire singular contributions, and the operator is not forbidden by the symmetry properties of the $`Q`$-matrix. The reason why we decided to show only the case of cubic lattices is that the distinction between the longitudinal from the transverse modes is a bit ambiguous at large momenta, where the both are in a sense massive. This is not a problem for cubic lattices, where one can show that the large momentum components of the transverse modes do not contribute, hence (C36) exhausts the whole contribution. On the contrary, in other cases, we do have to keep into account the contribution of the small-wavelength transverse modes to recover the full expression. This makes the calculations more involved than in the case of cubic lattices. ## D Explicit derivation of the RG equations In this Appendix we outline the derivation of the RG equations. Before that, it is convenient to list some useful results. ### 1 Averages Besides the propagators, we will also need the explicit expression of particular averages which enter in the derivation of the RG equations. The following results hold: $`B_{ab}P_{bc}B_{cd}^{}`$ $`=`$ $`4D(k)\delta _{ad}Tr\left(P_0\right),`$ (D1) $`B_{ab}P_{bc}B_{cd}`$ $`=`$ $`2D(k)P_{ad}^{},`$ (D2) where $`P=P_0+i\stackrel{}{P}\stackrel{}{\tau }`$ is a quaternion real matrix. In addition, $`A_{ab}P_{bc}A_{cd}`$ $`=`$ $`2D(k)P_{ad}^{}+4D(K)\delta _{ad}Tr\left(P_0\right)`$ (D3) $``$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}P_{ad},`$ (D4) $`C_{ab}P_{bc}C_{cd}`$ $`=`$ $`2D(k)P_{ad}^{}+4D(K)\delta _{ad}Tr\left(P_0\right)`$ (D5) $``$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}P_{ad},`$ (D6) $`A_{ab}P_{bc}C_{cd}`$ $`=`$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}P_{ad}.`$ (D7) For instance, if $`P=\widehat{I}`$, then $`BB^{}`$ $`=`$ $`4D(k)m\widehat{I}`$ (D8) $`BB`$ $`=`$ $`2D(k)\widehat{I}`$ (D9) $`AA`$ $`=`$ $`\left[2D(k)+4D(k)mD(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}\right]\widehat{I}`$ (D10) $`CC`$ $`=`$ $`\left[2D(k)+4D(k)mD(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}\right]\widehat{I}`$ (D11) $`AC`$ $`=`$ $`D(k){\displaystyle \frac{\mathrm{\Pi }}{\sigma _{xx}+\mathrm{\Pi }m}}\widehat{I}.`$ (D12) ### 2 RG equations First of all, we need to know the quaternion structure of the matrix $`D=TT^{}`$. Since $`D=D^{},cD^tc^t=\sigma _1D\sigma _1,`$ we can write $`D=D_0\sigma _0+D_3\sigma _3`$, where in the $`\pm `$-space $$D_0=\left(\begin{array}{cc}A_0& B_0\\ B_0^{}& C_0\end{array}\right),$$ (D13) and $$D_3=i\left(\begin{array}{cc}A_3& B_3\\ B_3^{}& C_3\end{array}\right).$$ (D14) We can write each of the above matrices in quaternion form, $`P=P^{(0)}\tau _0+i\stackrel{}{P}\stackrel{}{\tau }`$, where $`P^{(0)},\stackrel{}{P}e`$, but, in addition, we must impose that $`\begin{array}{cccc}\left(A_0^{(0)}\right)^t=A_0^{(0)}& \stackrel{}{A}_0^t=\stackrel{}{A}_0& \left(C_0^{(0)}\right)^t=C_0^{(0)}& \stackrel{}{C}_0^t=\stackrel{}{C}_0\\ \left(A_3^{(0)}\right)^t=A_3^{(0)}& \stackrel{}{A}_3^t=\stackrel{}{A}_3& \left(C_3^{(0)}\right)^t=C_3^{(0)}& \stackrel{}{C}_3^t=\stackrel{}{C}_3\end{array}`$ By making use of the above properties, and by means of the Eqs.(D2), (D1), (D4), (D6) and (D7), after defining $`L={\displaystyle _{\mathrm{\Lambda }/s}^\mathrm{\Lambda }}{\displaystyle \frac{d^2k}{(2\pi )^2}}D(k),`$ we get the following results $`W_fD_0W_f_f=L_A\mathrm{\Gamma }D_0`$ (D15) $`+2L_B\left(\begin{array}{cc}0& B_0\\ B_0^{}& 0\end{array}\right)2L_A\left(\begin{array}{cc}A_0& 0\\ 0& C_0\end{array}\right),`$ (D20) where $`\mathrm{\Gamma }=\mathrm{\Pi }/(\sigma _{xx}+m\mathrm{\Pi })`$, and $`W_fD_3W_f_f=L_A\mathrm{\Gamma }D_3`$ (D21) $`+2L_Ai\left(\begin{array}{cc}A_3& 0\\ 0& C_3\end{array}\right)2L_Bi\left(\begin{array}{cc}0& B_3\\ B_3^{}& 0\end{array}\right)`$ (D26) $`4L\widehat{I}Tr\left(iA_3^{(0)}+iC_3^{(0)}\right).`$ (D27) For further convenience, when useful, we have labelled $`L_B`$ the propagator of $`B_f`$, and $`L_A`$ the ones of $`A_f`$ and $`C_f`$. The next useful result is $$W_fW_f_f=\left(4L_Bm4L_Am+2L_A+L_A\mathrm{\Gamma }\right)\widehat{I}.$$ (D28) Through (D20) and (D27) we therefore get $`F_1_f=`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^2}}{\displaystyle 𝑑R}2Tr\left[\stackrel{}{D}\sigma _1Q_{sp}W_f\stackrel{}{D}W_fQ_{sp}\sigma _1\right]_f`$ (D30) $`+2Tr\left[\stackrel{}{D}\sigma _1Q_{sp}\stackrel{}{D}Q_{sp}\stackrel{~}{W}_f^2\sigma _1\right]_f`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{32}}{\displaystyle 𝑑R\mathrm{\hspace{0.17em}4}\left(L_BL_A+2L_Am+2L_Bm\right)Tr\left(\begin{array}{cc}0& B_0\\ B_0^{}& 0\end{array}\right)^2}`$ (D33) $`+`$ $`4\left(2L_Am+2L_Bm\right)\left(\begin{array}{cc}iA_3& 0\\ 0& iC_3\end{array}\right)^2`$ (D36) $``$ $`2L\left[Tr\left(D\sigma _3\right)\right]^2`$ (D37) $``$ $`4\left(2L_Am+2L_Bm+2L_A\right)\left(\begin{array}{cc}A_0& 0\\ 0& C_0\end{array}\right)^2`$ (D40) $``$ $`4\left(2L_Am+2L_Bm+L_A+L_B\right)\left(\begin{array}{cc}0& iB_3\\ iB_3^{}& 0\end{array}\right)^2.`$ (D43) The calculation of $`F_2^2_f`$ is more involved, since one needs the average of four $`W_f`$’s. For sake of lengthy, we just quote the final result that such a term cancels (D40) and (D43). We next notice that $`{\displaystyle \frac{1}{\mathrm{\Sigma }^2}}Tr\left[Q^{}Q\right]=2Tr\left[Ds_3\sigma _1Ds_3\sigma _1DD\right]`$ $`=4Tr\left(\begin{array}{cc}0& B_0\\ B_0^{}& 0\end{array}\right)^24Tr\left(\begin{array}{cc}iA_3& 0\\ 0& iC_3\end{array}\right)^2,`$ and that $`Tr\left[Q^{}Q\sigma _3\right]=2\mathrm{\Sigma }^2Tr\left(D\sigma _3\right).`$ Therefore, for $`L_A=L_B`$, $`F_1{\displaystyle \frac{1}{2}}F_2^2_f`$ (D45) $`=4Lm{\displaystyle \frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^2}}{\displaystyle 𝑑RTr\left[Q^{}Q\right]}`$ (D46) $`{\displaystyle \frac{1}{2}}L{\displaystyle \frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^4}}\left[Tr\left(Q^{}Q\sigma _3\right)\right]^2.`$ (D47) In the cases in which the $`A`$ and $`C`$ modes are gaped ($`L_A=0`$, and no $`D_3`$), we obtain the standard result $$F_1\frac{1}{2}F_2^2_f=\left(2Lm+L\right)\frac{2\pi \sigma _{xx}}{32\mathrm{\Sigma }^2}𝑑RTr\left[Q^{}Q\right].$$ (D48) ## E Gade and Wegner’s proof of the vanishing $`\beta `$-function The equations from (45) to (51) imply that $`W_3`$ is not a traceless matrix. Indeed, we can alternatively write $`W`$ as $$W=W^{}+\frac{1}{4m}Tr\left(W_3\right)\sigma _3W^{}+i\varphi \sigma _3,$$ (E1) with $`W^{}=W_0\sigma _0+W_3^{}\sigma _3`$, now being $`W_3^{}`$ a traceless matrix. Since $`\sigma _3`$ commutes with $`W^{}`$, this means that $$T(R)^2=\mathrm{e}^{W(R)}=\mathrm{e}^{i\varphi (R)\sigma _3}\mathrm{e}^{W^{}(R)}\mathrm{e}^{i\varphi (R)\sigma _3}V(R),$$ (E2) which also defines the matrix field $`V(R)`$. By means of this parametrization, the non-linear $`\sigma `$-model (3) can also be written as $`S[T]=S[V,\varphi ]`$ $`=`$ $`{\displaystyle \frac{2\pi \sigma _{xx}}{32}}{\displaystyle 𝑑RTr\left[\stackrel{}{}V(R)^1\stackrel{}{}V(R)\right]}`$ (E3) $`+`$ $`{\displaystyle \frac{m\pi }{2}}\left(\sigma _{xx}+m\mathrm{\Pi }\right){\displaystyle 𝑑R\stackrel{}{}\varphi (R)\stackrel{}{}\varphi (R)}.`$ (E4) Therefore the action of $`V`$ is distinct from that of $`\varphi `$, and the latter, being a phase, is gaussian. This implies that the combination $`\sigma _{xx}+m\mathrm{\Pi }`$ is not renormalized and scales with its bare dimension $`ϵ`$, for any number of replicas. In turns, it means that, in the zero replica limit, it is $`\sigma _{xx}`$ which is not renormalized! This is completely equivalent to the nice proof given by Gade and Wegner that the quantum corrections to the $`\beta `$-function of the conductance of a $`\mathrm{U}(N)/\mathrm{SO}(N)`$ model vanish at all orders in the $`N0`$ limit. The other important result concerns the renormalization of an operator $`T^{2q}=\mathrm{e}^{iq\varphi \sigma _3}V^q.`$ Within RG, $$\mathrm{e}^{iq\varphi \sigma _3}\frac{t}{2}q^2\mathrm{ln}s\left(\frac{t}{c+mt}\frac{1}{m}\right)\mathrm{e}^{iq\varphi \sigma _3}.$$ (E5) The second term, which is singular in the $`m0`$ limit, has to be canceled by the one-loop renormalization of $`V^q`$. Gade and Wegner showed that this cancellation holds for any $`q`$. Furthermore, they argued that, apart from the one-loop correction, the renormalization of $`V^q`$ does not contain any other singular term in the $`m0`$ limit. This, as they pointed out, has very important consequences. In $`2d`$, $`t`$ does not scale, while $`c`$ goes to zero. Therefore the term which dominates the renormalization of $`T^{2q}`$ for $`m=0`$ is just the first term in the right hand side of (E5). This argument implies that the one-loop correction, which we have derived for the density of states ($`q=1`$ case), is sufficient to identify the correct asymptotic behavior. To conclude, let us discuss more in detail the origin of this gaussian field $`\varphi `$. In the Grassmann variable path-integral representation, the action for the particle-hole symmetric model at $`E=\omega =0`$ $`S={\displaystyle \underset{RR^{}}{}}\overline{\mathrm{\Psi }}_RH_{RR^{}}\mathrm{\Psi }_R^{},`$ posseses a simple abelian gauge-like symmetry $$\mathrm{\Psi }\mathrm{e}^{i\varphi \sigma _3}\mathrm{\Psi },$$ (E6) because $`\{\sigma _3,H_{RR^{}}\}=0`$. It is just this symmetry which causes the appearance of the gaussian part of the non-linear $`\sigma `$-model. Notice that this symmetry implies a particle-hole symmetric Hamiltonian, which is invariant under the transformation $`c_{1,R}^{}c_{1,R}^{}`$ but $`c_{2,R}^{}c_{2,R}^{}`$. In fact, $`\{\sigma _3,H_{RR^{}}\}=0`$ also means that $`\{\tau _1\sigma _3,H_{RR^{}}\}=0`$, being $`H_{RR^{}}\tau _0`$. The interesting fact is that (E6) is not a symmetry of the fermion operators. Indeed, under this transformation, $`c_R\mathrm{e}^{i\varphi \sigma _3}c_R`$, but $`\overline{c}_R\mathrm{e}^{i\varphi \sigma _3}\overline{c}_R`$, and not $`\overline{c}_R\mathrm{e}^{i\varphi \sigma _3}\overline{c}_R`$ as we would expect if $`c_R`$ and $`\overline{c}_R`$ had to be identified with the operators $`c_R^{}`$ and $`c_R^{}`$. Finally, we notice that if, besides $`\sigma _3`$, the Hamiltonian commutes with another Pauli matrix (as it can be the case for specifically built particle-hole symmetric models), the above gauge symmetry would be non abelian, hence spoiling all peculiar properties which we have shown to occur. Indeed, in the last case, the system can be mapped into a standard localization model with an additional sublattice index.
warning/0002/astro-ph0002396.html	ar5iv	text	# The spectral variations of the O-type runaway supergiant HD 188209 ## 1 Introduction Runaway O stars have been defined as a group by Blaauw (1961), who introduced the term runaway to describe the space motions of AE Aur and $`\mu `$ Col. Blaauw (1961) has also suggested that such stars were ejected in the breakup of binary systems in supernova explosions by their companions. In later evolutionary stages, the initial secondary appears as a most massive star and transfers matter to the compact companion (the initial primary) making the system appear as a massive X-ray binary (van den Heuvel 1976). Given the possibility of the binary nature of runaway stars, it appears to be an important task to measure the radial velocity (RV) variations of the photospheric lines. Systematic searches for RV variations have been made in order to assess the binary frequency of O stars (e.g. Garmany, Conti & Massey 1980; Stone 1982, Gies 1987). In many cases the amplitude of RV variations is quite large, and the additional presence of a clear periodicity immediately suggests a binary nature for the system. However, there are stars which show more complicated RV curves, and the interpretation of their spectral variability is not straightforward. HD 188209 (O9.5Iab) is one of those objects. Garmany et al. (1980) have concluded from three spectra that this star is probably not a binary, and that the RV variations must be attributed to atmospheric motions. This conclusion was supported by Musaev & Chentsov (1988). However, based on 21 measurements Stone (1982) has concluded that HD 188209 can be considered as a spectroscopic binary with a period 57 days and small semiamplitude. More recently, Fullerton, Gies & Bolton (1996) included HD 188209 in a large sample of stars investigated on the presence of line profile variability (LPV) and found LPVs only in He i 5876 Å. However, they did not flag HD 188209 as a velocity variable (their Table 10). The binarity of many O supergiants has been proposed recently by Thaller (1997). The fact that binaries have a higher incidence and an H$`\alpha `$ emission strength in post-MS stages may indicate that wind interactions are a common source of emission in massive stars. In other words, even in cases where RV measurements are not available, the presence of H$`\alpha `$ emission in a spectrum could be linked with colliding winds. One needs to study orbital phase variations in the H$`\alpha `$ profile in order to be sure that the latter is due to colliding winds instead of some other mechanism. Note that HD 188209 is an X-ray source detected by ROSAT (Berghoefer, Schmitt & Cassinelli 1996). In this paper we focus on the high-resolution spectroscopic data of HD 188209. Our observations can possibly account for the small semi-amplitudes and eccentric orbits of this binary candidate since they have been accumulated at different periods over a long baseline. ## 2 Observations The observations have been carried out in different runs (Table 1) using the the Coudé Echelle Spectrometer (Musaev 1993) at the 1-m telescope of the Special Astrophysical Observatory of the Russian Academy of Science. Most of the spectra have a signal-to-noise ratio S/N$``$100 per resolution element, and an average resolution $`R=`$ 40000 in the wavelength region 4400–7000 Å. Preliminary reduction of the echelle spectra CCD images was made using the dech code (Galazutdinov 1992), which allows the flat-field division, bias/background subtraction, one-dimensional spectrum extraction from two-dimensional images, excision of cosmic-ray features, spectrum addition, correction for diffuse light, etc. Numerous bias, flat-field have been obtained every night. Each image was subject to a bias-frame subtraction and flat-field division using nightly means. Comparison exposures of a Th-Ar lamp were taken for each stellar spectrum. The control measurements of interstellar Ca ii and Na i D lines revealed a small scatter of 0.8 $`\mathrm{km}\mathrm{s}^1`$ (1 $`\sigma `$). However, the 1 $`\sigma `$ dispersion of the velocity of the DIB was 1.5 $`\mathrm{km}\mathrm{s}^1`$. These interstellar lines have been used to align all the spectra in the time series accurately. As an indicator of the overall precision of our measurements we have adopted 1 $`\sigma `$ dispersion 1.5 $`\mathrm{km}\mathrm{s}^1`$ of the velocity of the DIB. All stellar absorption lines exhibited variations about their respective mean velocities at least 2-3 times the dispersion of the DIB velocities. The mean rms obtained from different dispersion curves was at least 0.003 Å. The Coudé Echelle Spectrograph was not a subject to mechanical and/or thermal instabilities. We used a 580$`\times `$530 (pixel size 24 $`\times `$ 18 $`\mu `$m) CCD camera in all runs except 8 and 9. The last run was carried out using the Coudé Echelle Spectrograph (Musaev, 1999) at the 2-m telescope located in Terskol (North Caucasus, Russia). The CCD used in last two runs had a larger matrix (WI 1242$`\times `$1152 pixel with pixel size 22.5 $`\times `$ 22.5 $`\mu `$m) allowing a coverage in a single exposure of the region $``$ 3500–10100 Å with almost the same resolution. ## 3 Photometry The photometry of HD 188209 obtained with Hipparcos (ESA 1997) is presented in Figure 1. The approximate response curve for the Hp<sub>dc</sub> passband (see van Leeuwen et al. 1997) is very extended with a maximum at $``$ 4500 Å. Our observations did not reveal a strong variability in the equivalent widths of the strongest lines in the spectrum of HD 188209, which suggests that the variations detected by Hipparcos in the Hp<sub>dc</sub> band are due to changes in the continuum flux. The mean value of Hp<sub>dc</sub>=5.605 has a standard deviation 0.0155. The level of photometric variability is significant and cannot be ascribed to standard errors in the Hp<sub>dc</sub> magnitude of the order of 0.005 (van Leeuwen et al. 1997). Tests on the periodicity of the photometric variations were performed but no convincing period has been found. However, considering the large gaps in different Hipparcos measurements we cannot definitely rule out short-term photometric periodicity. ## 4 Fundamental Parameters of HD 188209 ### 4.1 Plane–parallel analysis To analyse the spectrum of HD 188209 we first determine the rotational velocity from the width of 11 metal lines of C, N, O, Si, Mg and Ca, adopting a Gaussian instrumental profile with a FWHM of 0.13 Å. The resulting value was a projected rotational velocity of 82.0 $`\pm `$ 8.5 km s<sup>-1</sup>, in good agreement with the value of 87$`\mathrm{km}\mathrm{s}^1`$ reported by Penny (1996), and between those given by Conti and Ebbets (1977) and Howarth et al. (1997) who give 70 and 92$`\mathrm{km}\mathrm{s}^1`$, respectively. The method followed in determining the stellar parameters from the spectrum using NLTE, plane–parallel hydrostatic model atmospheres has been described in detail by Herrero et al. (1992, and references therein). Briefly, we determine, at a fixed helium abundance, the gravity that best fits the different H and He profiles at a given temperature for a set of temperatures. If the abundance is right, the lines in the $`T_{\mathrm{eff}}`$$`\mathrm{log}g`$ diagram will ideally cross at a point, giving the stellar $`T_{\mathrm{eff}}`$ and $`\mathrm{log}g`$. Usually, they form an intersection region, whose central point is taken as giving the stellar parameters, and whose limits give the adopted error. If the lines do not cross at any point, the helium abundance is changed. The helium abundance giving the smaller intersection region for all profiles is the one selected. The center of the intersection region is taken again as that giving the stellar parameters. Recently, McErlean, Lennon & Dufton and Smith & Howarth have shown that different He i lines give different helium abundances in the region of the $`T_{\mathrm{eff}}`$$`\mathrm{log}g`$ diagram occupied by HD 188209. They attribute this to the effect of microturbulence and show that a value of around 10 km s<sup>-1</sup> is appropriate for bringing most of the He i lines into agreement. Thus, we adopt this value for HD 188209 and carry out the analysis in the way described above. With this method we have determined the stellar parameters of HD 188209. We obtained $`T_{\mathrm{eff}}`$= 31 500 K $`\pm _{500}^{1000}`$, $`\mathrm{log}g`$= 3.0 $`\pm 0.1`$ (uncorrected for centrifugal force; because of the low rotational velocity we neglect this small correction here) and $`ϵ`$= 0.12 $`\pm 0.03`$ (the abundance of helium with respect to the total abundance of hydrogen plus helium, by number; the solar abundance is $`ϵ`$= 0.09). The final fits for the H and He lines are shown in Figure 2. The larger errors towards higher temperatures is due to a small difference in the two spectrograms available for He ii 4541 Å. Using the second one, we would have obtained a $`T_{\mathrm{eff}}`$ of 32 000 K, all other parameters remaining the same. For this reason, we have enlarged the error bar in this direction. Remember that the errors given are formal errors, in the sense that they express the uncertainty in the fit using the models described above. As has been show by Herrero the metal opacity in the UV (line blocking) could also affect the value of the parameters determined. However, at the relatively low temperature of HD 188209 the effect would be minor, moving $`T_{\mathrm{eff}}`$ towards higher temperatures within the error box. With the stellar parameters given above we can determine the radius, luminosity and mass of HD 188209 as described in Herrero et al. (1992). For $`M_\mathrm{v}`$= $``$6.0 mag ($`M_{\mathrm{bol}}`$ = $``$9.0 mag) given by Howarth & Prinja (1989) we obtain $`R/R_{}`$= 20.9, $`\mathrm{log}(L/L_{})`$= 5.59 and $`M/M_{}`$= 16.6. The errors are again as in Herrero et al. (1992): $`\pm `$0.06 in $`\mathrm{log}(R/R_{})`$, $`\pm `$0.16 in $`\mathrm{log}(L/L_{})`$ and $`\pm `$0.22 in $`\mathrm{log}(M/M_{})`$. HD 188209 does not formally show the helium discrepancy, as the solar helium abundance is within the error bars. However, it shows the mass discrepancy: the mass derived from the plane–parallel spectroscopic analysis is, even including the error bars, much lower than the one derived from the evolutionary tracks from Schaller et al. . ### 4.2 Unified model analysis After having the parameters from the plane–parallel analysis, we can try to use a spherical, non-hydrostatic model atmosphere in order to improve the already derived parameters and also to obtain the mass-loss rate. In a supergiant like HD 188209 this can have an important impact on the final parameters. Usually, it is also assumed that this will contribute to the reduction in the mass discrepancy. The unified code we use is that recently developed by Santolaya–Rey, Puls & Herrero . The reader will find all the details therein, but for our present purposes we mention that the code uses spherical geometry, with a $`\beta `$–velocity field, and treats the wind and the photosphere in a unified way. It also makes use of the NLTE-Hopf functions. Stark broadening is included in the formal solution, and the model atoms are the same as in the plane–parallel case (slightly adapted for the new program). We begin by estimating the mass-loss rate from the H$`\alpha `$ profile, and then, with this mass-loss rate, we try again to find the best gravity (from the H$`\gamma `$ wings), effective temperature and He abundance (from the He ionization equilibrium). With the new parameters, we again try to fit the H$`\alpha `$ by varying the mass-loss rate, and so on. In the whole process we take the wind terminal velocity from Haser (1995), who gives 1700 km s<sup>-1</sup>. Note that Howarth et al. (1997), give a similar value of 1650 km s<sup>-1</sup>. We have used for the analysis the same lines as in Fig. 2, but have included He i 4471 Å instead of H$`\beta `$(that improved as does H$`\gamma `$) to show the dilution effect (see below). In Fig. 3 we show the fit of the H$`\gamma `$, He i and He ii lines. The adopted parameters are now $`T_{\mathrm{eff}}`$= 31 500 K, $`\mathrm{log}g`$= 3.00 and $`ϵ`$= 0.12, i.e. the same parameters as in the plane–parallel case, even for the gravity (again adopting a microturbulent velocity of 10 km s<sup>-1</sup>). Thus, the unified models do not contribute in this case to changes in the mass discrepancy found with the plane–parallel models (nor in the He abundance). We can also see in Fig. 3 that the He i 4471 line shows the well known dilution effect mentioned by Voels et al. , although the other He lines fit perfectly well. This effect merits an explanation. The fitting of the H$`\alpha `$ profile needed to derive the mass-loss rate cannot be done properly. The profile is highly variable, and we have adopted a qualitative approach: we have simply tried to give upper and lower limits to the mass-loss rate. As an example, we illustrate the procedure in Fig. 4, where we show one of the profiles with various mass-loss rates for the stellar parameters given above. The theoretical mass-loss rate values shown in Fig. 4 are log Ṁ= $``$5.70, $``$5.80 and $``$5.90 (with Ṁ in $`M_{}`$/yr). This situation is similar to that found by Herrero et al. for Cygnus X–1. The profile shown cannot be adopted as either an average or a representative one, as the profile varies a lot. The figure is only for illustrative purposes. An average logarithmic mass-loss rate for HD 188209 would be between $``$6.0 and $``$5.7. These values are also in agreement with the profiles in Fig. 3, where the adopted mass-loss rate was log Ṁ=$``$5.80 (corresponding to 1.58 10<sup>-6</sup> $`M_{}`$/yr). We should point out that other values of the mass-loss rate give worse fits to the H$`\gamma `$ and He i lines, by strongly modifying the line cores. ## 5 Radial velocity variations of the absorption lines The radial velocity variations in stars can have instrumental, internal (atmospheric) and/or external (Keplerian) origin. Instrumental effects in our measurements are minimized due to the high resolution and high S/N of the data presented. The internal accuracy achieved for the wavelength calibrations is of the order 1.5 $`\mathrm{km}\mathrm{s}^1`$ as derived from the scatter of measured radial velocities of interstellar and telluric lines in the spectra. Note that all former studies of HD 188209 (except for five spectra obtained by Fullerton et al. 1996) were based on the photographic spectra. Atmospheric pulsations of early-type stars have been the subject of extensive studies (e.g. Burki 1978; Bohannan & Garmany 1978, Kaufer et al. 1997, Fullerton et al. 1996) and have been reviewed by Baade (1988, 1998). Gies (1987) has compared the velocity distribution and binary frequency among 195 Galactic O-type stars (cluster and association, field and runaway) and found a deficiency of spectroscopic binaries among field stars (and especially among the runaway). The most comprehensive radial velocity studies of OB and BA supergiants to date are those of Fullerton et al. (1996) and Kaufer et al. (1996; 1997). Radial pulsation periods with P$``$5<sup>d</sup> have been predicted for O-type supergiants (Burki 1978; de Jager 1980; Levy et al. 1984). The presence of pulsations or random motions in stellar atmospheres results in complex velocity curves for different spectral lines due to stratification effects (Abt 1957). This is in contrast to Keplerian motions where all the lines vary synchronously with time (Ebbets 1979; Garmany et al. 1980). However, in many cases the amplitude of the RV variations does not exceed 25–30 $`\mathrm{km}\mathrm{s}^1`$, and it is very difficult to distinguish whether these variations are of an internal or an external nature. Additional difficulty comes from the possible presence of non-radial pulsations (NRP). Variable profiles (LPVs) have been detected in many narrow-line supergiants (e.g. Baade 1988; Kaufer et al. 1997; Fullerton et al. 1996) and in some cases they correspond to the radial velocity variations measured in photographic spectra. However, in contrast to broad-line supergiants, clear evidence of NRPs in narrow- and intermediate-line O-type supergiants has not yet been found. Although LPVs are a common occurrence among the O-type stars, some of them (like HD 34656; Fullerton, Gies & Bolton 1991) show a variability which consists of cyclical fluctuations in radial velocity due to pulsations in a fundamental mode. Careful inspection of our time series showed that all absorption lines varied in position over the course of a day. We have detected a few asymmetric profiles of He lines but there were no signatures of moving features. It is however possible that the S/N ratio of our data is not high enough to trace LPVs. Small distortions in profile shape are typically less than 1% of the continuum strength and a minimum S/N ratio required would be at least 300 (Fullerton et al. 1996). ### 5.1 The velocity-excitation relationship We have selected unblended lines by utilizing theoretical synthetic spectra computed for the model atmosphere of HD 188209 and measured RVs by fitting a Gaussian to the line profile. The following groups of lines have been selected: He i (4921, 4471, 4713, 5015, 5047, 5876, 6678 Å), He ii (5411, 4541, 4686 Å), Si iv (4654, 4631 Å), C iii (4650, 4647, 5696 Å), Si iii (4567, 4552, 4574 Å), C iv (5801, 5811 Å), N iii (4514, 4510, 4518, 4523 Å), O ii 4661 Å and Mg ii 4481 Å. The average RVs computed for each of these groups of lines are listed in Table 2. It is normally assumed that lines of different total excitation energy (TEE, ionization energy plus excitation energy of the lower level) form in different layers of the atmosphere (Hutchings 1976). However, it is not clear whether the stratification exists in a dynamically active, pulsating atmosphere. One can assume that the time scale of dynamical processes (pulsations, stochastic motions etc.) is much less than the time necessary for the establishment of radiative equilibrium. Thus, while pulsating, the atmosphere is supposed to pass through a chain of hydrostatic stages. The assumption that TEE correlates with line formation depths can be tested as well. We applied our plane–parallel models to compute formation depths of the line cores of the He i and He ii lines. It appeared that the core of the strong He ii 4686 Å line forms much closer to the surface (at column mass $``$ 0.01 g cm<sup>-2</sup>) than any of the He i lines. However, this was not the case with the two other He ii lines (5411 and 4541 Å) we used in our study. Similar tests have been carried out for other groups of lines but assuming an LTE line formation. These exercises suggest that before combining lines in different groups and computing their average RVs, one has to be sure (of course under the assumption that our plane–parallel models are applicable) that their depths of formation are similar. The velocity–excitation relationship found for many hot supergiants (Hutchings 1976) exists also in HD 188209. In Figs. 5 and 6 we present plots of mean RV versus TEE and standard deviations of the mean RV versus TEE for the different groups of lines computed for all dates (the last two lines in Table 2). These plots exclude a pure Keplerian motion as the only cause of the RV variations. Pulsations and stochastic motions (intrinsic wind variations caused by some hydrodynamical instabilities) in the wind can bring to the RV variations as well. If we suppose that stochastic variations are not important, then the existence of a standard deviation-TEE relationship would suggest that deeper layers in the atmosphere pulsate with smaller amplitudes. The amplitude of pulsations increases when approaching to the surface. ### 5.2 The periodicity of the radial velocity variations The period search was carried out with help of the period package (Dhillon & Privett 1997) of the starlink software. The following strategy was applied when looking for a periodic signal in the RV curves of different groups of lines. Due to the large gaps (especially between runs 5 and 6) in our observations, we first decided to study each of the runs 5, 6 and 7 separately. The clean algorithm (Roberts et al. 1987) was employed to cover a space of loop gains from 0.2 to 0.6 and the number of iterations from 10 to few hundreds. The convergence of the periodograms was achieved for the majority of groups of lines in all three runs. The mean frequency suggested by most of the groups in all three runs is 0.44$`\pm `$0.05 days<sup>-1</sup> (2.24 days). However, this period is very close to the Nyquist frequency (1/(2$`\times `$Smallest Data Interval)) of the data and might be misleading. We have also looked for periodic signals in the combined data of all groups obtained in all runs (Table 2). The maximum and minimum frequencies were set to 100 and 0, respectively. A clean analysis of the time series of the majority of groups revealed a frequency of 0.51$`\pm `$0.1 days<sup>-1</sup> (1.95 days). The gain factor was 0.1 at the first iteration, then was decreased by 15-20 iteration until stabilization. The average RV for each date obtained by averaging the RVs of all the groups revealed a frequency of 0.47$`\pm `$0.12 days<sup>-1</sup> (2.1 days). Again, both frequencies are very close to the Nyquist frequency and we should discard them. We must point out that our periodograms did not show any peaks at frequencies smaller than 0.4 days<sup>-1</sup>. The next strongest peak which appeared in our periodograms was near 0.156$`\pm `$0.15 days<sup>-1</sup> (6.4 days). Clearly this period is not affected by sampling. We have also analysed the RV data using the Lomb-Scargle method (Lomb 1976, Scargle 1982) which allows to compute statistical probability of peaks in periodograms. To ensure reliable significance values, the minimum number of permutations was set 100. The probability that the period is not equal to 6.4 days was always less than 30 $`\%`$. The peak at 6.4 days appears in all periodograms but given its significance value, we cannot definitely rule out its non-physical nature. In Figs. 7 and 8 we show the clean and the lomb-scargle periodograms and the fitting of a sin curve to folded data, respectively. ## 6 Variability of H$`\alpha `$ All hot supergiants have variable H$`\alpha `$ profiles in their spectra (Rosendahl 1973). The shape of the H$`\alpha `$ may vary from P Cyg to inverse P Cyg, double-peaked, pure absorption and/or emission (Ebbets 1982) with typical time-scales of the order of days. The nature of this variability is not yet understood. The existence of variable asymmetric outflows/infalls of matter and some corotating structures related to surface inhomogenities and possible magnetic fields have been proposed for BA-type (Kaufer et al. 1996) and O-type (Fullerton et al. 1996; Kaper et al. 1997) supergiants. In addition, there have been detailed studies of the rotating giant loop in $`\beta `$ Orionis (Israelian, Chentsov & Musaev 1997) and the corotating spiral structures in HD 64760 and HD 93521 (Howarth et al. 1998; Fullerton et al. 1997). It is of course very difficult to distinguish binary systems from single stars without understanding the nature of H$`\alpha `$ variability. As Thaller (1997) suggests, the H$`\alpha `$ can suffer some peculiar variability due to the colliding winds in a binary system. The time evolution of H$`\alpha `$ profiles in three different runs is shown in Figure 9. The average H$`\alpha `$ profile consists of three components, a central emission accompanied by blue and red absorptions. We have not observed a single H$`\alpha `$ profile without a central reversal. The emission is not always centered exactly on the rest wavelength but is varying. It may approach the continuum level, go above it and decrease rapidly in strength. Apparently the time-scale of the H$`\alpha `$ variability is at least one day. The 5<sup>th</sup> run has been divided into two parts (runs 5a & 5b) with four successive nights in each. The H$`\alpha `$ variability is observed over a wide range from about $``$400 to 200 $`\mathrm{km}\mathrm{s}^1`$. We have already seen in Section 4.2 that our spherical unified models can account for the central reversal (or at least set upper and lower limits of the mass-loss). Thus, we know that the central emission forms in the expanding envelope and accounts for the filling-in effect observed in H$`\alpha `$, H$`\beta `$ and some other lines. The filling-in effect observed in H$`\beta `$ correlated perfectly with the strength of the H$`\alpha `$ central reversal. The H$`\alpha `$ wings originate deep in the atmosphere, whereas the central reversal comes from the thin layers of the envelope. One cannot use the term “underlying photospheric absorption line” since the central reversal is not emitted by a detached layer far from the photosphere. It is important to stress that we deal with a $`single`$ line formed in a $`unified`$ model atmosphere. A variable amount of incipient emission can be due to density (or radius) variations in the outer atmosphere. However, these variations are not expected to produce an asymmetry as long as we deal with spherically symmetric mass-loss. Our observations indicate that the velocity of the central emission varies as well. We $`do`$ expect RV variations in the central emission (the upper atmosphere) since we know that the RVs of all absorption lines vary. The amplitude of the RV variations of Mg ii 4481 Å can reach 30$`\mathrm{km}\mathrm{s}^1`$ (Table 2) and 40$`\mathrm{km}\mathrm{s}^1`$ in H$`\beta `$. Thus, it is not unusual for the amplitude of RV variations of the H$`\alpha `$ central emission to reach 50$`\mathrm{km}\mathrm{s}^1`$. A period analysis of the RV curves of the central emission resulted in the detection of quasi-periodic variability with a frequency 0.42$`\pm `$0.14 d<sup>-1</sup> (2.35 days). It turns out that the RV curves of all groups of lines vary in phase. However, the RV curve of the H$`\alpha `$ was shifted half phase relative to all groups. We have measured the RVs of the blue and red absorption components of the H$`\alpha `$ and plotted them against the residual intensity of the central emission reversal. This plot (Fig. 10) shows a correlation with a large scatter due to the RV variations of the central reversal. This kind of correlation can be expected when the central reversal is moving up and down relative to a local continuum. We found similar correlations between the RV of the central emission and the residual intensities of the red and blue absorptions. Apparently all the changes observed in the H$`\alpha `$ wings at velocities $`v`$ 100$`\mathrm{km}\mathrm{s}^1`$ and $`v`$ +100$`\mathrm{km}\mathrm{s}^1`$ are due to the variations of the central reversal. We do not anticipate such large RV variations deep in the atmosphere where these wings are formed. The conclusion is that the overall shape of the H$`\alpha `$ is determined by the central emission. We have performed a period search of the integrated equivalent width (EW) of the H$`\alpha `$ data set and found a maximum in power at frequency 0.22 day<sup>-1</sup>. Figures 11 and 12 show the phase diagram for the period 4.41 days and a grey-scale representation of the phase spectrum, respectively. A phase spectrum represents a two-dimensional case of the clean algorithm where each velocity bin of the H$`\alpha `$ is treated as a time series of the H$`\alpha `$ intensity. ## 7 Discussion Our target belongs to the group of stars for which the existence of a compact companion has been proposed in the literature. The task of disproving or confirming the binary nature of the system can be tackled only if sufficiently accurate analysed observational data are available. In this paper we used state-of-the-art models of atmospheres to determine the fundamental parameters of HD 188209. To establish the presence of a possible companion we have studied the RVs of absorption lines by combining them in different groups. Fourier analysis based on the iterative clean algorithm was used to search for periodic variability. Unfortunately the time coverage of runs 1–4 and 8–9 was too sparse to set constraints on their time-dependent behaviour. For this reason we first analysed a few runs separately and then utilized the clean algorithm to search for periods in a whole data set. The highest peak in the Fourier power spectrum was centered near the frequency 0.156 day<sup>-1</sup> (6.4 days). The 6.4 days period can be due to the binary nature of the system if one assumes very small and unlikely values for the mass ratio (q$``$0.1). Taking the values derived in this article ($`M/M_{}`$= 16.6, $`R/R_{}`$=20.9) and assuming q=0.1, we obtain for the Roche radius and for the major semi-axis of the binary orbit 15 $`R_{}`$ and 25 $`R_{}`$, respectively. This simple estimate shows that an O supergiant can hardly fit within the orbit because its Roche radius would be less than the stellar radius. Even if it would fit, the tides in such a tight binary would be very strong making the star to speed up quickly until the rotation period matches the orbit. Even if we assume that the system is very young, it’s hard to explain that the putative orbital period is 2 times shorter than the rotational period (13 days). The second difficulty with the binary interpretation comes from the variability of H$`\alpha `$. A TVS (temporal variance spectrum, showing the extent and distribution of statistically significant profile variability) has been computed recently (Baade 1998b, Kaper et al. 1998) for 15 spectroscopic binaries and it was found that all they show a characteristic double-peaked profile. This is due to two H$`\alpha `$ absorption/emission profiles moving in a composite spectra. In our case the H$`\alpha `$ profile is splitted because of the central emission coming from the lower wind. The last argument comes from the clear relations between excitation energies and radial velocity amplitude and excitation energy and mean radial velocities and from the model atmosphere atmosphere calculations. The latter is a good discriminant between internal variations (pulsations & wind instabilities) and Keplerian motions. In a binary system one would expect all lines to have the same amplitude independent on their TEE. It has been known for a long time (Abt 1957) that the quasi-periodicity in hot supergiants might be ascribed to radial pulsations. A simple relation (Burki 1978; de Jager 1980) can be used to estimate the period of radial pulsation, $$\mathrm{log}P_{\mathrm{fund}}=10.930.5\mathrm{log}(M/M_{})0.38M_{\mathrm{bol}}3\mathrm{log}T_{\mathrm{eff}}$$ (1) Using the values of parameters obtained in Section 4 we arrive at $`P_{\mathrm{fund}}`$=1.75 d. Note that the form of the relation (1) depends on the stellar evolutionary models and the input parameters; both are subject to large errors. In particular, note that we found no large differences in the parameters determined with plane–parallel and unified model atmospheres. Nevertheless, Levy et al. (1984) have pointed out that periods a factor of 1.5 longer than the corresponding periods of the radial pulsations can be ascribed to non-radial pulsations. This means that a factor of two difference between the evolutionary and the spectroscopic masses can easily result in the mis-identification of the pulsating mode. Another difficulty has been pointed by the referee of the article. A more sophisticated approach shows (Unno et al. 1979) that f-mode pulsation (which is the lowest-frequency mode supported by radial pulsation) periods are about 10 times larger than the one suggested by a period-luminosity relation. In any case, the theoretical period of 1.75 days is very close to the Nyquist frequency of our data which means that we have a little chance to identify it in our data set even if it exists. Our data not allow to distinguish between pulsations and stochastic variations of the stellar wind. It is also quite possible that we have a combination of both effects. Note that the projected rotational period of this star ($``$ 13 d) is much longer than any of the quasi-periods found in this paper (but of course our runs do not cover a whole rotation cycle). The surface features (if any) will always be visible on the projected disc of the star independently of the inclination angle. Thus, any periods due to the rotation of these features must correspond directly to the rotation period. We do not find any peaks in the power spectra at $``$13 d and this leads us to discard rotational modulation as a possible explanation of the RV variations reported here. The quality and the sampling of our data do not allow a careful study of the line asymmetries, moving components (if later exists) and/or long-term spectroscopic variability to be made. It is quite possible that the non-sinusoidal character of the RV curve for 6.4 days period (Fig 8) is caused by some disturbances due to the NRPs and/or moving features in the profiles plus any stochastic instabilities of the wind. New monitoring with much higher S/N may allow NRPs, multimode pulsations and clearly separate a sinusoidal curve of the radial pulsations to be revealed. However, we found convincing evidence that the atmospheric motions cannot be ascribed solely to Keplerian motions and probably are not of a binary origin. ## Acknowledgments We thank D. Baade, O. Pols and Pablo Rodriguez for their useful comments. A. G. thanks the Canadian Astronomical Society for the travel grant to SAO, and I. Bikmaev for helpful discussions. We wish to thank the anonymous referee for his careful reading of the manuscript and several constructive suggestions.
warning/0002/astro-ph0002340.html	ar5iv	text	# The Toronto Red-Sequence Cluster Survey: First Results ## 1. The TRCS #### The Toronto Red-Sequence Cluster Survey (TRCS) is a major new observational effort designed to identify and characterize a large sample of galaxy clusters to redshifts as high as $`z1.4`$. When completed, the TRCS will be the largest imaging survey ever completed on 4m telescopes, and will provide a large and homogeneous sample of galaxy clusters for detailed follow-up study. The basic survey is envisioned as 100 deg<sup>2</sup> of 2 filter ($`R`$ and $`z^{}`$) imaging, to a depth which is $``$2 mag past $`M^{}`$ at $`z=1`$ in both filters. The design of the survey is based on a new method for identifying galaxy clusters (Gladders & Yee 2000a) developed specifically for the TRCS. In brief, this method searches for clustering in the 5-D space of: x-y positions, $`Rz^{}`$ color, $`z^{}`$ mag, and morphology in the form of a concentration index. The x-y positions provide the surface density enhancement. A color slice in the color-mag plane provides separation in $`z`$ space via the red sequence of early-type galaxies in clusters (Figure 1) and increases the S/N of density enhancements. Morphology allows us to key onto early-type galaxies, the primary population in cluster centers. ### 1.1. Scientific Goals #### The TRCS is being driven by two major scientific goals. The first is based on the theoretical prediction that the evolution of the mass-spectrum of galaxy clusters with redshift, $`N(M,z)`$, should be a strong function of two cosmological parameters, $`\mathrm{\Omega }_m`$ and $`\sigma _8`$ (Figure 2). The goal is to use the clusters identified in the survey to measure $`N(M,z)`$ directly from the survey data. Redshift can be estimated from the color of red sequence (e.g., López-Cruz & Yee 2000), and the mass of each cluster can be estimated from its richness, as measured by the parameter $`B_{gc}`$ (e.g., Yee & López-Cruz 1999). The second major scientific goal is a study of the cluster galaxy populations, which can be done using the TRCS for the first time with a complete sample. The definition of a complete, or volume limited, sample is derived from extensive simulations of the survey selection functions. ### 1.2. Survey Completeness and Selection Functions #### Any detailed understanding of the cosmological or galaxy evolution results deriving from the TRCS requires a good understanding of the survey selection functions. Specifically, we wish to know how well the cluster-finding algorithm finds clusters of various sorts (as described by various parameters). To this end, we have constructed a number of cluster and field simulations (Gladders & Yee 2000b) to directly test the algorithm. A large suite of possible clusters have been tested; the parameters describing the clusters are given in Table 1. The results of this process demonstrate that the TRCS should be complete for all reasonable clusters of Abell Richness Class $``$ 1 clusters ($`\sigma _v750`$ km s<sup>-1</sup>) to at least $`z=1.1`$. ## 2. Some First Results The first run for the TRCS occurred at CFHT in May, 1999. A total of 21 pointings were acquired with the CFH12K camera, with each pointing covering 0.272 deg<sup>2</sup>. The bulk of the images have seeing better than 0$`.^{\prime \prime }`$7, with some as good 0$`.^{\prime \prime }`$5. At the time of writing, the total TRCS dataset consists of about 35 deg<sup>2</sup> of data. However, the results presented below are based on only the first 6 deg<sup>2</sup>. Figure 3 shows a rich ($`B_{gc}2000`$), compact cluster at photometric redshift of $`z0.95`$ (left panel). The cluster appears to be embedded in a large filamentary structure ($``$10 h<sup>-1</sup> Mpc long) traced out by red galaxies (center panel). The excess is undetectable without a color cut (right panel). The efficacy of color information in isolating high-$`z`$ structures is obvious. #### Figure 4 shows the cores of several other rich clusters, one of which appears to have a gravitational arc. The success of the TRCS in finding rich, high-z cluster candidates in the few degrees searched so far implies that the total survey will contain several hundred $`z0.8`$ cluster candidates, a preliminary result which is supportive of a low-density, high normalization cosmological model. ## 3. Secondary Projects The TRCS dataset is also well suited for a number of other detailed studies. For example, preliminary work has already revealed a significant population of extremely red ($`Rz^{}3.5`$) point sources. Such objects are likely L and T dwarfs, with some contamination by $`z5.5`$ QSOs. Other studies are possible using the survey data, e.g. studies of cluster lensing (strong and weak), cosmic shear, halo structure, low surface brightness galaxies and early-type galaxy correlations. ## References Gladders, M.D., & Yee, H.K.C. 2000a, in prep. Gladders, M.D., & Yee, H.K.C. 2000b, in prep. Kodama, T., & Arimoto, N. 1997, 320, 41 López-Cruz, O., & Yee, H.K.C. 2000, (to be submitted to ApJ) Yee, H.K.C., & López-Cruz, O. 1999, AJ, 117, 1985
warning/0002/math0002136.html	ar5iv	text	# On Krammer’s Representation of the Braid Group ## 1 Introduction In \[BW\], Joan Birman and Hans Wenzl constructed a two-parameter family of algebras related to braid groups and the Kauffman knot polynomial, and analyzed its semisimple structure. The same algebras were discovered simultaneously and independently by Jun Murakami in \[M\]. Since the braid group maps homomorphically into this algebra, the representations on the simple summands of the algebra give irreducible representations of the braid group. We will focus on one particular irreducible representation (for each $`n`$) out of this collection, and show that it is identical to an independently discovered representation. This is parallel to a similar identification which has long been known between the Burau representation and a certain summand of the Hecke algebra. Daan Krammer, in \[K\], defined a representation of the braid group using its interpretation as the automorphism group of the punctured disk. Using geometric arguments, he constructed a skein relation between items called forks. From the skein relation, he was able to list a set of algebraic equations describing the action of braid group generators on these forks, some of which form a basis for an invariant module. These equations thus describe matrix entries for the representation. Using techniques related to the solution of the word problem in the braid groups in \[BKL\], Krammer was then able to prove that this representation was faithful for $`n=4`$, thus showing ###### Theorem. (Krammer) $`B_4`$ is linear. This discovery revolutionized the study of braid group linearity, since for a long time it seemed that the Burau representation was the best candidate to give a faithful representation of the braid group. Nevertheless, Krammer’s construction was not widely publicized until Stephen Bigelow expanded Krammer’s result using topological methods: ###### Theorem. (Bigelow) Krammer’s representation is faithful for all $`n`$, thus the braid groups are linear. This finally answered an important question in braid theory which had stood since the introduction of the Burau representation in 1935. The main result of this paper gives a connection between Krammer’s representation and the Birman-Murakami-Wenzl (BMW) algebra: ###### Main Theorem. Krammer’s representation of the braid group $`B_n`$ is identical to the $`(n2)\times 1`$ irreducible representation of the BMW algebra. My proof uses purely algebraic methods, and can be deduced solely from the equations given in \[K\] for the action of the braid group, and the equations given in \[BW\] for relations within the algebra. One immediate consequence of this theorem is ###### Corollary 1. Krammer’s representation is irreducible. And another follows quickly from Bigelow’s result: ###### Corollary 2. The regular representation of the BMW algebra is faithful. Vaughan Jones also discovered this main result, simultaneously and independently. However, his methods are different from mine, and involve a somewhat deeper understanding of the algebraic structure. The representation which Krammer spelled out explicitly was originally discovered by Ruth Lawrence in \[L\] as a representation of the Hecke algebra, using homology of configuration spaces similar to those used by Bigelow. This connection relates to a topological interpretation of the Jones polynomial which is still in the process of being explored. With the addition of Krammer’s presentation of the representation, we can also make connections to other algebras and knot polynomials. ### 1.1 Acknowledgements This paper will be part of my Ph.D. thesis. I’d like to thank my advisor, Joan Birman, for her invaluable assistance. Besides her everyday commentary and guidance, she brought to my attention many of the structures and concepts used in this work, and had several ideas for ways to try to connect them. I would also like to thank Stephen Bigelow, who helped me understand his proof and Lawrence’s contribution, through both correspondence and a visit to Columbia. I would like to thank Daan Krammer and Vaughan Jones for correspondence concerning their contributions in this area, and Justin Roberts, who shared with me his notes and insights from attending relevant talks given by both Bigelow and Jones. ## 2 Background and Definitions The braid group $`B_n`$ has many definitions and interpretations. Perhaps the easiest to see is a pictorial one. A braid is a diagram consisting of two horizontal bars, one at the top and one at the bottom of the figure, with $`n`$ nodes on each bar (usually drawn equally spaced), and $`n`$ strands, always running strictly downward, connecting the upper and lower nodes. This figure represents an isotopy class of embeddings of the strands in 3-space, so the strands are allowed (in fact, required) to cross over and under each other rather than intersect, and the directions of these crossings are marked in a conventional way on the diagram. Among all braids, those which are isotopic in $`\times I`$ are identified. Braids form a group. The identity element is the braid with no crossings, multiplication is concatenation (draw one diagram above the other, and erase the center bar), and the standard set of generators (known as Artin generators) consists of braids $`\sigma _i`$ ($`1i<n`$), where the only crossing is that of the $`i`$th strand under the next strand. See Figure 1(a). The relations in the braid group are easily described: $$\sigma _i\sigma _j=\sigma _j\sigma _i\text{ for }\|ij\|2$$ (1) $$\sigma _i\sigma _{i+1}\sigma _i=\sigma _{i+1}\sigma _i\sigma _{i+1}.$$ (2) One of the long-standing questions about the braid group has been whether it is linear, that is, whether there exists any faithful representation into a matrix group. A common method of constructing representations of the braid group is to map the braid group homomorphically into a finite-dimensional algebra, and use the algebra’s regular representation. One such algebra that was used is the Hecke algebra, a deformation of the complex symmetric algebra $`S_n`$, and another is the Birman-Murakami-Wenzl algebra, similarly a deformation of the Brauer algebra. The BMW algebra $`C_n(l,m)`$ can be defined on invertible generators $`G_i`$ ($`1i<n`$), which satisfy the braid relations described above, and non-invertible elements $`E_i`$ defined via the formula $$G_i+G_i^1=m(1+E_i).$$ (3) The additional relations are: $`E_iE_{i\pm 1}E_i=E_i`$ (4) $`G_{i\pm 1}G_iE_{i\pm 1}=E_iG_{i\pm 1}G_i=E_iE_{i\pm 1}`$ (5) $`G_{i\pm 1}E_iG_{i\pm 1}=G_i^1E_{i\pm 1}G_i^1`$ (6) $`G_{i\pm 1}E_iE_{i\pm 1}=G_i^1E_{i\pm 1}`$ (7) $`E_{i\pm 1}E_iG_{i\pm 1}=E_{i\pm 1}G_i^1`$ (8) $`G_iE_i=E_iG_i=l^1E_i`$ (9) $`E_iG_{i\pm 1}E_i=lE_i`$ (10) $`E_i^2=(m^1(l+l^1)1)E_i`$ (11) $`G_i^2=m(G_i+l^1E_i)1.`$ (12) See \[BW\]. Some of these relations can be deduced from others; however, I am not concerned here with a minimal presentation. For those readers who are, one can be found in \[W1\]. In addition, it can be deduced that $`E_i`$ commutes with both $`E_j`$ and $`G_j`$, if $`\|ij\|2`$. As in the case of the braid group, we can associate several elements of the algebra to pictures. The generators $`G_i`$ can be identified with the same braid diagrams (Figure 1(a)) described above for $`\sigma _i`$. The elements $`E_i`$ can be identified with similar diagrams where the $`i`$th node on the top is joined to the next node on top, and similarly on the bottom (with all other nodes connected vertically as with $`G_i`$). See Figure 1(b). Since the “strands” do not run strictly downward, this is not a braid element. However, we can still “multiply” these diagrams by each other and by braids, to again get isotopy classes of embeddings in 3-space. The reader may notice that each of the above relations in the algebra which does not involve addition actually represents an isotopy relation of these braid-like elements. Consequently, every monomial in the algebra can be uniquely represented by such a picture, up to isotopy. The irreducible representations of an algebra such as this one can be identified using its Bratelli diagram, which encodes the decomposition of each semisimple algebra $`C_n`$ into its simple summands. The regular representation on each summand (which, as a vector space, is an invariant subspace) is then an irreducible representation of the algebra and, by extension, the braid group which maps into it. The connections in the diagram between levels encode the inductions or restrictions of the representations to the next or previous level (by adding or ignoring the final generator). The Bratelli diagram of the Hecke algebra is the same as the Bratelli diagram of the complex group algebra of the symmetric group. So each simple module, thus each irreducible representation, of $`H_n`$ is indexed by a Young diagram with $`n`$ boxes, and connections from a particular module to each succeeding level are to the Young diagrams obtained by adjoining one box to the Young diagram in question. Define notation as follows: let $`V_{n,\lambda }`$ denote the module or representation in the $`n`$th level of the Bratelli diagram, labeled by the Young diagram $`\lambda .`$ The Hecke algebra is a quotient of the BMW algebra (obtained by setting $`E_i=0`$). Consequently, it appears as a direct summand of $`C_n`$, and its Bratelli diagram is contained in that of the larger algebra. The complete structure of the Bratelli diagram of $`C_n`$, as is explained in \[W1\], is the same as the Brauer algebra $`D_n`$: ###### Theorem. (Wenzl) (a) $`D_n`$ is semisimple. (b) The simple components of $`D_n`$ are labeled by the set of Young diagrams with $`n2k`$ boxes ($`k^+`$). (c) If $`V_{n,\lambda }`$ is a simple $`D_n`$ module it decomposes as a $`D_{n1}`$ module into a direct sum of simple $`D_{n1}`$ modules $`V_{n1,\mu }`$, where $`\mu `$ ranges over all Young diagrams obtained by removing or (if $`\lambda `$ contains fewer than $`n`$ boxes) adding a box to $`\lambda `$. This means that the Bratelli diagram can be easily constructed using an inductive method. The $`n=1`$ level is a single module, indexed by the Young diagram consisting of a single box. The $`n=2`$ level has three modules, one indexed by the empty Young diagram, and the other two indexed by the two Young diagrams of two boxes. All three are connected to the module on the previous level. Starting at level $`n=3`$, we can construct the levels using reflections. Reflecting the $`n2`$ level of the diagram, including its connections to the $`n1`$ level, across a line drawn through the $`n1`$ level, gives the portion of the $`n`$ level indexed by Young diagrams with less than $`n`$ boxes. This portion of the algebra, following \[BW\], is given the notation $`H_n^{}`$. The portion indexed by Young diagrams with $`n`$ boxes is then constructed as usual from the modules at the $`n1`$ level with $`n1`$ boxes in their Young diagrams. This portion of the algebra is given the notation $`H_n`$, and is isomorphic to the Hecke algebra. Consequently, all monomials in the algebra which contain a $`E_i`$ factor are located in $`H_n^{}`$. The Bratelli diagram of $`C_n`$ up to $`n=4`$ is shown in Figure 2. As with any Bratelli diagram, the dimension of each module (and its representation) is the sum of the dimensions of the representations it restricts to, which is thus the number of paths from the top of the diagram. For $`H_n`$, these dimensions can also be found from the hook length formula. ## 3 Dimension Argument Since we are claiming that the Krammer representation, which has dimension $`\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$, is an irreducible representation of $`C_n`$, it would be good to know that a representation of the appropriate size exists: ###### Theorem 1. In the $`n`$th level of the Bratelli diagram of $`C_n`$, the representations labeled by rectangular Young diagrams of shapes $`(n2)\times 1`$ and $`1\times (n2)`$ have dimension $`\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$. ###### Proof. Due to the symmetric and dual nature of the Young diagrams, both arguments will be similar, and I only need discuss one. So without loss of generality, let $`\lambda `$ denote the $`(n2)\times 1`$ Young diagram, consisting of 1 row with $`n2`$ boxes. We can use the Bratelli diagram (see Figure 3) to count the dimension of each of the representations, as it will be the sum of the dimensions of the representations on the previous level that this one is connected to. Since $`\lambda `$ has fewer than $`n`$ boxes, this module is part of $`H_n^{}`$, and using the inductive construction of the Bratelli diagram, we can see that all the connections leading into this representation come from connections leading out of the representation $`V_{n2,\lambda }H_{n2}`$ (and, in fact, go to/from the same representations on the $`n1`$ level). These connections are of two types: those connecting $`V_{n2,\lambda }`$ to modules in $`H_{n1}`$, and those connecting it to modules in $`H_{n1}^{}`$. The connections to $`H_{n1}`$ are easy to see, as they result from the standard Young’s lattice. There is one representation with a rectangular Young diagram of shape $`(n1)\times 1`$, of dimension 1, and one representation whose Young diagram has a single box in the second row. By the hook length formula, this representation has dimension $`n2`$. Whatever connections exist between $`V_{n2,\lambda }`$ and $`H_{n1}^{}`$ again come from reflections of connections from $`V_{n2,\lambda }`$ to $`C_{n3}`$. Of these, there is only one: $`V_{n2,\lambda }H_{n2}`$ only connects to $`H_{n3}`$, and since $`\lambda `$ has shape $`(n2)\times 1`$, it only connects to the representation labeled by the Young diagram $`\mu `$ of shape $`(n3)\times 1`$. So our third and last downward connection from $`V_{n2,\lambda }`$ (and by reflection, upward connection from $`V_{n,\lambda }`$) is to $`V_{n1,\mu }`$. Notice that this is precisely the $`n1`$ version of the representation we are investigating! Thus by induction, we may assume that this representation has dimension $`\left(\genfrac{}{}{0pt}{}{n1}{2}\right)`$. (One may verify that the claimed formula holds for a sufficiently low base case.) The dimension of the representation $`V_{n,\lambda }`$ is therefore $$\left(\genfrac{}{}{0pt}{}{n1}{2}\right)+(n2)+1=\left(\genfrac{}{}{0pt}{}{n}{2}\right).$$ ## 4 Explicit form of the representation In Jones \[J\], it is shown that for any braid index $`n`$, there exists a 1-dimensional irreducible representation of the Hecke algebra. More specifically, the Hecke algebra has a 1-dimensional invariant subspace which is preserved under multiplication by any of the Hecke algebra (or braid) generators. It is similarly shown in \[BW\] and \[W1\] that $`C_n`$, which contains a subalgebra isomorphic to the Hecke algebra, also contains a 1-dimensional irreducible representation of $`C_n`$ for any index $`n`$. This is easy to see on the Bratelli diagram as the representation $`V_{n,\lambda }`$ (with $`\lambda `$ as in the proof of Theorem 2). Thus we again have a 1-dimensional invariant subspace which is preserved under multiplication by any of the braid generators. Note that this subspace does not consist of the same algebra elements as in the Hecke algebra. The quadratic relation is different if we consider the full Birman-Murakami-Wenzl algebra, so either subspace is not preserved if we use the multiplication rule corresponding to the wrong algebra. Pick an algebra element which is in this 1-dimensional subspace (and thus generates it). We will call it $`v`$. Since the actions of the braid group on this vector (by left multiplication) result in a 1-dimensional representation, all such multiplications are scalar multiplications by the same factor, which we will call $`\kappa `$. To prove that the representation of the braid group $`B_n`$ given explicitly by Krammer is a representation of $`C_n`$, we will make use of the vector $`v`$ which corresponds to the braid index $`n2`$. Thus, $`vV_{n2,\lambda }H_{n2}`$ can be expressed as a linear combination of monomials in $`C_n`$ using only the positive generators $`G_1,\mathrm{},G_{n3}`$. An explicit computation of $`v`$ in these terms is possible, but not required for the computations in this paper. An inductive construction can be found in \[W1\]. Importantly, $`v`$ has the property that $$G_iv=\kappa vi<n2.$$ Now we will begin constructing a representation of $`C_n`$ which we will later show is equivalent to Krammer’s representation. A basis for the invariant subspace is the following vectors: for $`1i<jn`$, let $$\begin{array}{cc}\hfill T_{ij}& =\left(\underset{k=i}{\overset{j2}{}}G_k\right)\left(\underset{l=j1}{\overset{n1}{}}E_l\right)v\hfill \\ & =G_i\mathrm{}G_{j2}E_{j1}E_j\mathrm{}E_{n1}v.\hfill \end{array}$$ (13) Pictorially, each vector $`T_{ij}`$ corresponds to the braid-like diagram where the last two nodes on the bottom bar are connected, and the $`i`$th and $`j`$th nodes on the upper bar are connected (under all other strands), and the remaining nodes connect from top to bottom without crossing each other. See Figure 4. (Pictures much like these are used in Jones’ proof.) We will occasionally ignore the restriction $`i<j`$ in order to write more general statements. Now we will see how these vectors behave under a left action by the braid group. Traditionally, Artin braid generators $`\sigma _i`$ are mapped into $`C_n`$ under the map $`\sigma _iG_i`$; however, that mapping turns out not to work for our purposes. We will instead rescale the braid group, and use the mapping $`\sigma _iG_i/\kappa `$. There are four different types of multiplication to consider (and three of them have subtypes which depend on the ordering of the indices): Type A: $`\sigma _iT_{i,i+1}`$ Type B: $`\sigma _iT_{jk}`$, with $`\{i,i+1\}\{j,k\}=\mathrm{}`$. Type C: $`\sigma _iT_{i+1,j}`$ Type D: $`\sigma _iT_{ij}`$ For each of these actions, we will map $`\sigma _iG_i/\kappa `$ as described above, and expand $`T_{ij}`$ as defined in equation (13), and simplify the resulting expression according to the relations given in equations (1)-(12). Most of the steps are straightforward; here are lemmas for those that are less so: ###### Lemma 1. $`E_iG_{i+1}=E_iE_{i+1}G_i^1`$ ###### Proof. Both expressions are simplifications of $`E_iE_{i+1}E_iG_{i+1}`$. ∎ ###### Lemma 2. $`E_{i2}G_iE_{i1}E_i=E_{i2}E_{i1}G_{i2}E_i`$ ###### Proof. Both expressions can simplify to $`E_{i2}G_{i1}^1E_i`$. ∎ Now to the left action: A. $$\sigma _iT_{i,i+1}=\kappa ^1G_i(E_iE_{i+1}\mathrm{}E_{n1}v)=\kappa ^1l^1E_i\mathrm{}E_{n1}v=\kappa ^1l^1T_{i,i+1}$$ B. Because of the form of $`T_{jk}`$, the calculation will depend on the order of the indices $`i,j,k`$. If $`i+1<j<k`$, then $`G_i`$ commutes past $`T_{jk}`$ to multiply by $`v`$: $$\begin{array}{c}\sigma _iT_{jk}=\kappa ^1G_i(G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}v)\hfill \\ \hfill =\kappa ^1(G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1})G_iv\\ \hfill =\kappa ^1(G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1})(\kappa v)=T_{jk}\end{array}$$ (14) If $`j<i<i+1<k`$, then $`G_i`$ commutes as far as it can, but it is stopped when it gets to $`G_{i1}G_i`$. However, the three-term braid relation lets it transform into a $`G_{i1}`$ and continue commuting to the right: $$\begin{array}{c}\sigma _iT_{jk}=\kappa ^1G_i(G_j\mathrm{}G_i\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}v)\hfill \\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}G_iG_{i1}G_iG_{i+1}\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}G_{i1}G_iG_{i1}G_{i+1}\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}G_{i1}G_iG_{i+1}\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}(G_{i1}v)\\ \hfill =\kappa ^1G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}(\kappa v)=T_{jk}\end{array}$$ (15) Similarly, if $`j<k<i`$, then $`G_i`$ commutes as far as it can, and then we use Lemma 2 to commute it the rest of the way, at the expense of lowering the index: $$\begin{array}{c}\sigma _iT_{jk}=\kappa ^1G_i(G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_i\mathrm{}E_{n1}v)\hfill \\ \hfill =\kappa ^1G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{i2}G_iE_{i1}E_i\mathrm{}E_{n1}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{i2}E_{i1}G_{i2}E_i\mathrm{}E_{n1}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}G_{i2}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{k2}E_{k1}\mathrm{}E_{n1}(\kappa v)=T_{jk}\end{array}$$ (16) C. Again, the exact calculations will depend on the order of the indices. If $`i+1<j`$, the multiplication trivially gives us new indices for $`T`$. $$\sigma _iT_{i+1,j}=\kappa ^1G_i(G_{i+1}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}v)=\kappa ^1T_{ij}$$ If $`j<i`$, then the multiplication will behave much as in the latter parts of Case B, but with the obstruction this time located at the junction of the $`G`$ terms and the $`E`$ terms of $`T_{j,i+1}`$. $$\begin{array}{c}\sigma _iT_{j,i+1}=\kappa ^1G_i(G_j\mathrm{}G_{i1}E_i\mathrm{}E_{n1}v)\hfill \\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}G_iG_{i1}E_i\mathrm{}E_{n1}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}E_{i1}E_i\mathrm{}E_{n1}v=\kappa ^1T_{ji}\end{array}$$ (17) D. This time, not only are the calculations dependent on the order of the indices, but the results come out differently as well. If $`i+1<j`$, then the first term in $`T_{ij}`$ is $`G_i`$, so we need to use the quadratic relation in $`C_n`$. $$\begin{array}{c}\sigma _iT_{ij}=\kappa ^1G_i(G_i\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}v)\hfill \\ \hfill =\kappa ^1(mG_i+ml^1E_i1)G_{i+1}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}v\\ \hfill =\kappa ^1(mT_{ij}T_{i+1,j}+ml^1E_iG_{i+1}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}v)\end{array}$$ (18) Two of the three summands we get are recognizable as $`T`$ vectors, but the other begins with an $`E`$ followed by a product of $`G`$s. For this, we need to apply Lemma 1. $$\begin{array}{c}E_iG_{i+1}G_{i+2}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}v\hfill \\ \hfill =E_iE_{i+1}G_i^1G_{i+2}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}v\\ \hfill =E_iE_{i+1}G_{i+2}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}(G_i^1v)\\ \hfill =E_iE_{i+1}G_{i+2}\mathrm{}G_{j2}E_{j1}\mathrm{}E_{n1}(\kappa ^1v)\end{array}$$ (19) It’s still not in the standard form for a $`T`$ vector, but notice what one application of Lemma 1 accomplished. We still have a word made up of $`E`$s followed by $`G`$s followed by $`E`$s, and our indices have not changed (they’re still increasing by one each time). We’ve just changed the first $`G`$ that appears into an $`E`$ of the same index (and gotten a $`\kappa ^1`$ scalar out of the multiplication). Repeated applications of Lemma 1 will have the same effect, and we can continue until all our $`G`$ terms have been transformed into $`E`$ terms. This gives us the vector $`E_i\mathrm{}E_{n1}=T_{i,i+1}`$, multiplied by one $`\kappa ^1`$ for each $`G`$ term we had before the first application of the lemma, thus $`\kappa ^{(j2i)}`$. Conclusion: $$\sigma _iT_{ij}=m\kappa ^1T_{ij}\kappa ^1T_{i+1,j}+ml^1\kappa ^{ij+1}T_{i,i+1}$$ If $`j<i`$, then things come out differently: $$\begin{array}{c}\sigma _iT_{ji}=\kappa ^1G_i(G_j\mathrm{}G_{i2}E_{i1}\mathrm{}E_{n1}v)\hfill \\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}G_iE_{i1}E_i\mathrm{}E_{n1}v=\kappa ^1G_j\mathrm{}G_{i2}G_{i1}^1E_i\mathrm{}E_{n1}v\\ \hfill =\kappa ^1G_j\mathrm{}G_{i2}(m+mE_{i1}G_{i1})E_i\mathrm{}E_{n1}v\\ \hfill =\kappa ^1(mT_{ji}T_{j,i+1}+mG_j\mathrm{}G_{i2}E_i\mathrm{}E_{n1}v)\end{array}$$ (20) The last of these terms simplifies by commuting each of the $`G`$ factors past all the $`E`$ factors. Each $`G`$ that reaches $`v`$ becomes multiplication by $`\kappa `$, and what remains is $`E_i\mathrm{}E_{n1}v=T_{i,i+1}`$. So we have $`\sigma _iT_{ji}`$ $`=\kappa ^1(mT_{ji}T_{j,i+1}+m\kappa ^{i1j}T_{i,i+1})`$ $`=m\kappa ^1T_{ji}\kappa ^1T_{j,i+1}+m\kappa ^{ij2}T_{i,i+1}`$ It is now clear that these vectors form an invariant subspace, of dimension $`\left(\genfrac{}{}{0pt}{}{n}{2}\right)`$. As we have described an action on this space by the generators of the braid group, we have a representation of the braid group. This representation has two (complex) parameters, $`m`$ and $`l`$. (Recall that $`\kappa `$ is the eigenvalue of the specific vector $`v`$, so is not a parameter in the same sense. However, it can be adjusted or reset by rescaling, as we did above.) A final detail is to locate this invariant subspace in $`C_n`$. ###### Theorem 2. The invariant subspace described above is $`V_{n,\lambda }`$. ###### Proof. From \[W2\], Prop. (1.2), it is clear that $`vV_{n2,\lambda }`$ implies $`T_{n1,n}=E_{n1}vV_{n,\lambda }`$. The cited proposition, applied to the present situation, states that if $`pV_{n2,\lambda }`$ is a minimal idempotent, then $`pE_{n1}`$ is a minimal idempotent of $`V_{n,\lambda }`$. For the present purposes, it is not necessary to deal with minimal idempotents, but notice that since $`V_{n2,\lambda }`$ is 1-dimensional, the minimal idempotent $`p`$ is a scalar multiple of $`v`$. Therefore, the cited proposition guarantees that a certain scalar multiple of $`vE_{n1}`$ is a minimal idempotent of $`V_{n,\lambda }`$, which is a stronger result than we require. (Notice also that since $`v`$ is written using only the generators $`G_1,\mathrm{},G_{n3}`$, it commutes with $`v`$.) ∎ ## 5 The Krammer representation We will now recall Krammer’s representation from \[K\], written in terms of actions on a module with basis $`v_{ij}`$, with $`1i,jn`$ and $`ij`$. Like the representation above, it has two (complex) parameters, $`q`$ and $`t`$. (Stephen Bigelow has found a fascinating interpretation of these parameters from a topological perspective.) Krammer’s presentation is longer than that described here, because he included formulae for multiplication by band generators (a larger set than Artin generators), and because I am taking the liberty to combine formulae when the order of the indices does not matter. $`\sigma _iv_{i,i+1}=tq^2v_{i,i+1}`$ $`\sigma _iv_{jk}=v_{jk}`$ $`\text{for }\{i,i+1\}\{j,k\}=\mathrm{}`$ $`\sigma _iv_{i+1,j}=v_{ij}`$ $`\text{for }ji,i+1`$ $`\sigma _iv_{ij}=tq(q1)v_{i,i+1}+(1q)v_{ij}+qv_{i+1,j}`$ $`\text{if }i+1<j`$ $`\sigma _iv_{ji}=(1q)v_{ji}+qv_{j,i+1}+q(q1)v_{i,i+1}`$ $`\text{if }j<i`$ To show that this is the same as the representation of $`C_n`$ constructed above, we will rescale the basis slightly and set a correspondence between parameters: ###### Theorem 3. Under the identifications $`q=\kappa ^2`$, $`m=\kappa (1q)`$, $`l^1=\kappa tq^2`$, $`v_{ij}=\kappa ^{i+j}T_{ij}`$, the two actions described above are identical. ###### Proof. A. $$\sigma _iv_{i,i+1}=\kappa ^{2i+1}\sigma _iT_{i,i+1}=\kappa ^{2i}l^1T_{i,i+1}=\kappa ^{2i+1}tq^2T_{i,i+1}=tq^2v_{i,i+1}$$ B. $$\sigma _iv_{jk}=\kappa ^{j+k}\sigma _iT_{jk}=\kappa ^{j+k}T_{jk}=v_{jk}$$ C. $$\sigma _iv_{i+1,j}=\kappa ^{i+j+1}\sigma _iT_{i+1,j}=\kappa ^{i+j}T_{ij}=v_{ij}$$ D. If $`i+1<j`$, $$\begin{array}{c}\sigma _iv_{ij}=\kappa ^{i+j}\sigma _iT_{ij}=m\kappa ^{i+j1}T_{ij}\kappa ^{i+j1}T_{i+1,j}+ml^1\kappa ^{2i+1)}T_{i,i+1}\hfill \\ \hfill =(1q)\kappa ^{i+j}T_{ij}\kappa ^{i+j1}T_{i+1,j}+\kappa ^2(1q)tq^2\kappa ^{2i+1)}T_{i,i+1}\\ \hfill =(1q)v_{ij}\kappa ^2v_{i+1,j}+\kappa ^2(1q)tq^2v_{i,i+1}=(1q)v_{ij}+qv_{i+1,j}+tq(q1)v_{i,i+1}\end{array}$$ (21) If $`j<i`$, $$\begin{array}{c}\sigma _iv_{ji}=\kappa ^{i+j}\sigma _iT_{ji}=m\kappa ^{i+j1}T_{ji}\kappa ^{i+j1}T_{j,i+1}+m\kappa ^{2i2}T_{i,i+1}\hfill \\ \hfill =(1q)\kappa ^{i+j}T_{ji}\kappa ^{i+j1}T_{j,i+1}+(1q)\kappa ^{2i1}T_{i,i+1}\\ \hfill =(1q)v_{ji}\kappa ^2v_{j,i+1}+\kappa ^2(1q)v_{i,i+1}=(1q)v_{ji}+qv_{j,i+1}+q(q1)v_{i,i+1}\end{array}$$ (22) ###### Remark. The computations in the proof above will actually work under any additional rescaling of $`v_{ij}`$ with respect to $`T_{ij}`$, namely using the identification $`v_{ij}=\kappa ^{i+j+k}T_{ij}`$, for any value of $`k`$. For simplicity in the proof, I set $`k=0`$, but if these identifications are to be used for any explicit calculations, I recommend instead using $`k=n+1`$, so that the values of the exponent range symmetrically from $`(n2)`$ to $`(n2)`$.
warning/0002/quant-ph0002029.html	ar5iv	text	# Quantum entanglement using trapped atomic spins ## A Possible trap parameters
warning/0002/cond-mat0002467.html	ar5iv	text	# Impurity spin magnetization of thin Fe doped Au films \[ ## Abstract In order to probe the influence of the surface-induced anisotropy on the impurity spin magnetization, we measure the anomalous Hall effect in thin AuFe films at magnetic fields up to $`15\mathrm{T}`$. The observed suppression of the anomalous Hall resistivity at low fields as well as the appearance of a minimum in the differential Hall resistivity at higher fields can be explained by our theoretical model which takes into account the influence of a polycrystalline film structure on the surface-induced anisotropy. Our results imply that the apparent discrepancy between different experimental results for the size effects in dilute magnetic alloys can be linked to a different microstructure of the samples. \] In very dilute magnetic alloys finite size effects may occur for sample dimensions comparable to the size of the Kondo screening cloud . While some experiments revealed a considerable decrease of the logarithmic Kondo anomaly in the resistivity of thin films and narrow wires already at the $`\mu \mathrm{m}`$ scale, other experiments showed an almost constant Kondo anomaly for wire widths down to $`40\mathrm{nm}`$. Recent theoretical calculations indicated the complex, dynamical nature of the screening cloud, implying that the simple picture of a static, spherically symmetric screening cloud is not correct . Zawadowski et al. linked the size dependent Kondo scattering to a surface-induced anisotropy of the magnetic impurity spins. Due to the interaction of the impurities with the conduction electrons, which suffer from the spin-orbit scattering by the non-magnetic host atoms, the impurity spins tend to be aligned parallel to the sample boundaries . The surface-induced anisotropy should remain active for more concentrated spin glass alloys. Finite size effects have indeed been observed in the resistivity of spin glasses with reduced dimensions . Finite size effects should also be observable in the impurity spin magnetization. Earlier magnetization experiments on single film and on multilayered spin glasses had linked a depression of the freezing temperature $`T_f`$ to a lower critical dimensionality for the spin glass transition. Here, we show that measuring the magnetic field dependence of the anomalous Hall effect in thin films of a AuFe spin glass provides a powerful method to probe the influence of the surface-induced anisotropy on the Fe spin magnetization. At low magnetic fields, the magnetization signal related to the spin glass freezing is strongly suppressed when compared to the bulk behavior. The appearance of an extra magnetization signal at higher fields can be linked to a re-orientation of impurity spins which are blocked by the surface-induced anisotropy at lower fields. The low field magnetization is further suppressed when moving the Fe doping towards the film surface. We have developed a theoretical model which takes into account the polycrystalline film structure. The presence of additional internal surfaces at the grain boundaries introduces a particular spatial dependence of the spin anisotropy which is essential for understanding the Hall effect data. Our model provides an explanation for the apparent discrepancies between the different experiments which have probed the size dependence of the Kondo and spin glass resistivity. Thin films of AuFe alloys have been prepared by codeposition on oxidized silicon wafers of Au ($`99.9999\%`$ purity) and Fe ($`99.99\%`$ purity) in an ultra high vacuum molecular beam epitaxy growth chamber. A multiterminal sample geometry is obtained by depositing the films through a contact mask. Apart from $`30\mathrm{nm}`$ thick Fe doped Au films and pure Au reference films, we have also prepared Au/AuFe/Au and AuFe/Au/AuFe trilayers. For the $`30\mathrm{nm}`$ thick trilayer samples the central layer has a thickness of $`15\mathrm{nm}`$, while the two outer layers have a thickness of $`7.5\mathrm{nm}`$. The Hall resistance as well as the longitudinal resistance are measured with an ac resistance bridge in perpendicular magnetic fields up to $`15\mathrm{T}`$. Adding Fe impurity spins to pure Au causes the appearance of a non-linear anomalous Hall component. Figure 1(a) shows typical Hall effect data at $`T=4.2\mathrm{K}`$ which have been obtained by McAlister and Hurd for polycrystalline pure bulk Au as well as for a $`0.98\mathrm{at}.\%`$ AuFe bulk spin glass. The anomalous Hall resistivity varies proportional to the Fe spin magnetization and its temperature dependence at low fields reveals a sharp peak near the freezing temperature $`T_f`$. In agreement with direct magnetization measurements, the peak is washed out at higher magnetic fields . As illustrated in Fig. 1(b), the anomalous Hall effect is much weaker in a thin $`2\mathrm{at}.\%`$ AuFe spin glass film when compared to the bulk material. Although the Fe con- centration is about twice as large as in Fig. 1(a), deviations from the classical linear Hall resistivity of the pure Au film are very small. In order to make the non-linear anomalous Hall component more clearly visible, we also plot in Fig. 1(b) the differential Hall resistivity $`\mathrm{d}\rho _H/\mathrm{d}B`$. The initial variation of the differential Hall resistivity at low fields reflects the destruction of the spin glass state due to the alignment of the impurity spins. The magnetic response resulting from this alignment has been reduced by an order of magnitude when compared to the bulk material. As discussed in detail below, we can attribute the reduced amplitude of $`\mathrm{d}\rho _H/\mathrm{d}B`$ at low fields as well as the appearance of a minimum at higher fields to the surface-induced anisotropy. The strength of the anisotropy of an impurity spin is predicted to increase when the impurity approaches the sample surface . Consequently, the anisotropy effects should be less pronounced for a Au/AuFe/Au trilayer when compared to a AuFe/Au/AuFe trilayer. This is confirmed in Fig. 1(c) for AuFe layers with an Fe concentration of $`3.5\mathrm{at}.\%`$. Although both the total thickness and the Fe content are the same, the non-linear behavior at low fields is clearly enhanced when moving the Fe impurities away from the surface. The reduced amplitude of $`\mathrm{d}\rho _H/\mathrm{d}B`$ for the trilayers when compared to a single AuFe film (see Fig. 1(b)) results from a shortcircuiting by the pure Au layer with a much lower resistivity. In order to calculate the non-linear impurity magnetization in polycrystalline AuFe films, we extended the theory of the surface-induced anisotropy to the case of thin films consisting of small grains. Following Ref. , the Hamiltonian $`_{\mathrm{an}}`$, which describes the surface-induced anisotropy for an impurity spin in an isolated brick-shaped grain, can be written as $`_{\mathrm{an}}=𝒜{\displaystyle \underset{\alpha ,\beta }{}}_{\alpha \beta }S_\alpha S_\beta (\alpha ,\beta =x,y,z).`$ (1) The $`S_\alpha `$ are the operators for the components of the impurity spin $`𝐒`$. In contrast to a semi-infinite sample , the strength of the anisotropy has to be described in terms of a matrix with elements $`𝒜_{\alpha \beta }`$. Moreover, the calculated matrix elements $`_{\alpha \beta }`$ are complicated functions of the dimensions ($`a_x`$, $`a_y`$, $`a_z`$) of the grain and of the impurity position. The material dependent constant $`𝒜`$ should range between 0.01 and $`1\mathrm{eV}`$ for dilute AuFe alloys . Taking the $`z`$-axis parallel to the magnetic field, the magnetization of an Fe spin is given by $`S_z={\displaystyle \frac{1}{2\mu _\mathrm{B}𝒵}}{\displaystyle \underset{k=1}{\overset{5}{}}}\mathrm{exp}\left({\displaystyle \frac{_k}{k_\mathrm{B}T}}\right){\displaystyle \frac{\mathrm{d}_k}{\mathrm{d}B}},`$ (2) where $`\mu _\mathrm{B}`$ is the Bohr magneton, and $`𝒵=_{k=1}^5\mathrm{exp}(_k/k_\mathrm{B}T`$). The index $`k=1,\mathrm{},5`$ labels the roots $`_k`$ of the secular equation $`\|_{S_z^{}S_z}\delta _{S_z^{}S_z}\|=0(S_z,S_z^{}=2,1,0,1,2)`$ (3) with the Hamiltonian $`=2\mu _\mathrm{B}S_zB+_{\mathrm{an}}`$. In Fig. 2 the calculated differential magnetization $`\left[\mathrm{d}S_z/\mathrm{d}B\right]_{\mathrm{gr}}`$ (the square brackets with subscript ‘gr’ correspond to an average over impurity positions within a grain) is shown as a function of the magnetic field for different lateral dimensions of the grains. The height $`a_z`$ of the grains is fixed at $`30\mathrm{nm}`$, i.e., the thickness of the AuFe film shown in Fig. 1(b). In the limiting case of a single-crystal film ($`a_x,a_y\mathrm{}`$), the impurity spin states at $`B=0`$ are known to be the eigenstates of $`S_z`$, the energy eigenvalues being proportional to $`S_z^2`$. At low temperatures, only the state with $`S_z=0`$ will be populated. Hence, the impurity spin does not respond to a weak magnetic field. With increasing $`B`$ the energy level with $`S_z=1`$ becomes lower than that with $`S_z=0`$. This gives rise to the first peak in the field dependence of $`\left[\mathrm{d}S_z/\mathrm{d}B\right]_{\mathrm{gr}}`$. The second peak appears when the state with $`S_z=2`$ becomes the ground state. The theoretical result for the single-crystal film in Fig. 2 clearly fails to describe our experimental results, since the calculated impurity spin magnetization remains zero at low magnetic fields (all spins are blocked parallel to the surface). In Ref. it was already shown that in narrow wires the presence of differently oriented surfaces leads to a rather intricate behavior of the impurity-spin anisotropy when compared to the case of a film. In brick-shaped grains, the competing influence of mutually perpendicular surfaces on the magnetic anisotropy gives rise to a partial cancellation of the anisotropy effect, resulting in a shift of the peaks in $`\left[\mathrm{d}S_z/\mathrm{d}B\right]_{\mathrm{gr}}`$ towards lower fields. Moreover, there exist specific locations within the grains where the magnetic anisotropy energy becomes negligibly small. At these locations the bulk behavior of the impurity spins is restored. Hence, a limited number of spins contribute to the spin glass freezing and are sensitive to very weak magnetic fields. In the inset of Fig. 2 we illustrate the large difference in the predicted response of spins located near the corners of the grains when compared to spins located more towards the center of the grains. For grains with lateral sizes larger than (but still comparable to) the height, our theoretical model predicts that a minimum in $`\left[\mathrm{d}S_z/\mathrm{d}B\right]_{\mathrm{gr}}`$ can appear at relatively small magnetic fields. For those grains the initial part of the theoretical curves is similar to the measured differential Hall resistivity shown in Fig. 1(b) and in Fig. 1(c). At higher temperatures ($`T>10\mathrm{K}`$) the anomalous Hall effect signal gradually weakens. A consistent description of our anomalous Hall data requires to assume a large value $`𝒜=0.12\mathrm{eV}`$ in Eq. (1). This implies that a larger fraction of the alignment of the Fe spins by the magnetic field occurs at extremely high fields which are not accessible in the experiment. Figure 3(a) shows the calculated field dependence of an average $`\left[\mathrm{d}S_z/\mathrm{d}B\right]_{\mathrm{cs}}`$ over the positions of the impurities located within a cross-section of a grain perpendicular to the magnetic field. The different curves correspond to different distances $`\mathrm{\Delta }z`$ from the cross section to the top (or bottom) of the grain. The theoretical curves can be compared to the experimental field dependence of $`\mathrm{d}\rho _H/\mathrm{d}B`$ for the Au/AuFe/Au and the AuFe/Au/AuFe trilayer shown in Fig. 1(c). The low field variation is much weaker for the AuFe/Au/AuFe trilayer. This is in agreement with our theoretical result that the surface-induced anisotropy becomes stronger for Fe spins which are closer to the top (or bottom) of a grain. The inset of Fig. 3(b) shows a scanning tunneling microscopy (STM) image of the surface of the Au/AuFe/Au sample, revealing a rather broad distribution of grain sizes. In order to quantitatively fit the experimental data we therefore calculate the magnetization for an ensemble of grains of different lateral size $`a`$, distributed lognormally with a statistical median $`\overline{a}`$ and a standard deviation $`\sigma `$. In Fig. 3(b), the measured differential Hall resistivity $`\mathrm{d}\rho _H/\mathrm{d}B`$ for the Au/AuFe/Au trilayer is compared to the field dependence of $`\left[\mathrm{d}S_z/\mathrm{d}B\right]_{\mathrm{gr}}`$, calculated for an ensemble of Au/AuFe/Au grains with $`\overline{a}=80\mathrm{nm}`$ and $`\sigma =1.5`$. We note that the value of $`80\mathrm{nm}`$ is larger than the typical grain size inferred from the STM images ($`20\mathrm{nm}`$). Our theoretical approach with isolated grains is strictly valid only when elastic defect scattering at the grain boundaries fully destroys the anisotropy which is caused by the conduction-electron-mediated interaction of an Fe spin with host atoms from neighboring grains. Using a larger grain size for the calculations allows to take into account that a fraction of the electrons still moves ballistically between adjacent grains. We conclude that our theoretical model for the surface-induced anisotropy of magnetic impurities in small metallic grains allows a quantitative description of the anomalous Hall resistivity in polycrystalline AuFe spin glass films. The small impurity spin magnetization at low magnetic field also provides an explanation for the small amplitude of the spin glass resistivity in thin AuFe films . It is clear that the surface-induced anisotropy strongly affects the spin glass freezing, implying that the analysis of previous measurements, which investigated the reduction of the freezing temperature $`T_f`$ in thin films and multilayers , should be revised. At this point it is not clear in how far this reduction, which also occurs for the thin film samples discussed in this Letter, is caused by an intrinsic finite size effect or can be totally accounted for by the surface-induced anisotropy. For Kondo alloys with a small impurity content ($`100\mathrm{ppm}`$) we are no longer able to measure the anomalous Hall effect. We expect the surface-induced anisotropy to still produce a strong reduction of the impurity spin magnetization, resulting in a suppression of the Kondo resistivity. Since the surface-induced blocking of a magnetic impurity spin is sensitive to the specific polycrystalline sample structure, the apparent discrepancy between different experimental results for the size dependence of the Kondo resistivity can be linked to a different microstructure of the samples. The collaboration between the universities of Antwerpen and Leuven has been supported by the Fund for Scientific Research – Flanders (Belgium) as well as by the Belgian Inter-University Attraction Poles research program on Reduced Dimensionality Systems (IUAP No. 4/10). Additional support has been obtained in Leuven from the Flemish Concerted Action (GOA) research program and in Antwerpen from the Scientific Fund (BOF) of the Universiteit Antwerpen.
warning/0002/cond-mat0002019.html	ar5iv	text	# Vortex Lattice Depinning vs. Vortex Lattice Melting: a pinning-based explanation of the equilibrium magnetization jump ## Abstract In this communication we argue that the Vortex Lattice Melting scenario fails to explain several key experimental results published in the literature. From a careful analysis of these results we conclude that the Flux Line Lattice (FLL) does not melt along a material- and sample-dependent boundary $`H_j(T)`$ but the opposite, it de-couples from the superconducting matrix becoming more ordered. When the FLL depinning is sharp, the difference between the equilibrium magnetization $`M_{eq}(T,H)`$ of the pinned and unpinned FLL leads to the observed step-like change $`\mathrm{\Delta }M_{eq}(T,H)`$. We demonstrate that the experimentally obtained $`\mathrm{\Delta }M_{eq}(T,H)`$ can be well accounted for by a variation of the pinning efficiency of vortices along the $`H_j(T)`$ boundary. and thanks: On leave from Instituto de Fisica, Unicamp, 13083-970 Campinas, Sao Paulo, Brasil. Supported by the Deutsche Forschungsgemeinschaft under DFG IK 24/B1-1, Project H, and by CAPES proc. No. 077/99. The phenomenological description of superconductors is based on the knowledge of their magnetic field-temperature $`(HT)`$ phase diagram. The equilibrium behavior of conventional type-II superconductors in an applied magnetic field is well known: At fields below the lower critical field $`H_{c1}(T)`$, the superconductor is in the Meissner-Ochsenfeld phase in which surface currents screen the magnetic field from the interior of the sample. Above $`H_{c1}(T)`$ the field penetrates the superconductor in the form of a lattice of vortices, the so-called Abrikosov vortex lattice. This flux-line-lattice (FLL) persists up to the upper critical field $`H_{c2}(T)`$ where superconductivity vanishes in the bulk of the sample. In high-temperature superconductors (HTS), however, due to strong thermal fluctuations, a first-order phase transition of the FLL to a liquid-like state, the “melting” of the FLL, has been predicted to occur well below $`H_{c2}(T)`$. Since then, an enormous amount of experimental and theoretical work has been done trying to find this, or other more sophisticated transitions and to extend the original theoretical treatment. At the beginning of these research activities it was claimed that the melting transition of the FLL in HTS manifests itself at the damping peak of vibrating HTS in a magnetic field. However, this interpretation was controversial . We know nowadays that the damping peaks in vibrating superconductors attributed to the melting transition can be explained quantitatively assuming thermally activated depinning and the diffusive motion of the FLL under a small perturbation generated by the vibration of the sample . To find the true experimental evidence for the melting transition of the FLL is by no means simple: Because the FLL interacts with the superconducting matrix through pinning centers (atomic lattice defects, surface barriers, etc.) every property of the FLL one measures will be influenced by the pinning and, therefore, no direct and straightforward proof of the melting phase transition can be achieved. Several years after the above cited first experimental attempt, a striking jump in the equilibrium magnetization in Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8</sub> (Bi2212) single crystals has been measured using a SQUID as well as sensitive micro-Hall sensors. This magnetization jump, which was interpreted as a first-order transition of the FLL, lies at more than one order of magnitude lower fields than the thermally activated depinning line measured with vibrating superconductors. These interesting resultsare important because if the melting transition would be of the first order, a discontinuous change in the equilibrium magnetization $`M_{eq}(T,H)`$ at the transition is expected. The step-like increase of $`M_{eq}(T,H)`$ (a decrease in absolute value) with increasing magnetic field and temperature has been also observed in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> (Y123) single crystals. Apart from the clearly defined magnetization jump, another important fact was revealed by the experiments. It has been found that the line $`H_j(T)`$ along which $`M_{eq}(T,H)`$ jumps, and the temperature dependence of the field $`H_{\mathrm{SMP}}(T)`$ where an anomalous maximum in the width of the magnetization hysteresis loop takes place (the so-called “second magnetization peak” (SMP)), define a unique boundary on the $`HT`$ plane for a given sample, demonstrating their intimate relationship. In order to account for this behavior, a second-order phase transition associated with the increase of the critical current density $`j_c(T,H)`$ was suggested as the origin of SMP based on a pinning-induced disordering of the FLL. Recently, a similar SMP occurring at a temperature dependent field $`H_{\mathrm{SMP}}(T)H_{c2}(T)`$ was also measured in conventional Nb superconducting films. The studies of the SMP performed on Nb films, Bi2212 single crystals, as well as on a non-cuprate isotropic single crystalline Ba<sub>0.63</sub>K<sub>0.37</sub>BiO<sub>3</sub> thick film provide a clear evidence that the SMP is not related to a critical current enhancement, but originates from a thermomagnetic instability effect and/or a non-uniform current distribution, leading to the “hollow” in $`M(H)`$ at $`H<H_{\mathrm{SMP}}(T)`$. In agreement with the models, the SMP vanishes in all these superconductors when the lateral sample size becomes less than $`100\mu `$m . Because $`100\mu `$m is much larger than all relevant vortex-pinning-related characteristic lengths, the strong influence of the sample geometry on the SMP cannot be explained by a change of the pinning efficiency of vortices. Besides, the results suggest that the interaction between vortices starts to dominate that between vortices and the matrix at $`H>H_{\mathrm{SMP}}(T)`$. We note that the results obtained in Bi2212 crystals are actually in good agreement with the second-order diffraction small-angle neutron scattering (SANS) experiments which revealed a well-ordered FLL at $`H>H_{\mathrm{SMP}}(T)`$ . Moreover, the formation of a more ordered FLL with increasing temperature due to thermal depinning has been found in the above mentioned SANS experiments at $`H>H_{\mathrm{SMP}}(T)`$ and for intermediate temperatures. A similar result was obtained by means of Lorentz microscopy in Bi2212 thin films near the low field - high temperature portion of the “irreversibility line”. The high resolution SANS measurements, recently reported for Y123 crystals also revealed a well defined FLL up to a field of 4 T and at low temperatures, i.e. above the $`H_{\mathrm{SMP}}(T)`$-line measured in similar crystals. Based on this experimental evidence and the intimate relationship between SMP- and the magnetization-jump-lines, we propose that the jump in $`M_{eq}(T,H)`$ results from a magnetic-field- and temperature-driven FLL depinning transition to a more ordered state of the FLL, effectively de-coupled from the atomic lattice. There exist already experimental as well as theoretical works that show that the interaction of vortices with pinning centers increases $`\|M_{eq}(T,H)\|`$, indicating clearly that pinning influences the thermodynamic, equilibrium properties of superconductors. If the FLL depinning is sharp, then one expects a step-like change of the equilibrium magnetization $`\|\mathrm{\Delta }M_{eq}(T,H)\|=\|M_{eq}^{dis}(T,H)M_{eq}(T,H)\|`$ along the $`H_j(T)`$ boundary. Here $`\|M_{eq}^{dis}(T,H)\|>\|M_{eq}(T,H)\|`$ is the absolute equilibrium magnetization in the presence of the quenched disorder which is measured below $`H_j(T)`$. We note further that the sharpness of the vortex depinning onset, irrespectively of the underlying mechanism, manifests itself as a sudden increase in the electrical resistivity at $`H_j(T)`$, below which the vortex behavior is irreversible. In what follows we present a phenomenological approach to describe the magnetization jump observed at the depinning transition. The equilibrium magnetization of an ordered, unpinned FLL in the London regime and neglecting fluctuations is given by the equation $$M_{eq}=\frac{\varphi _0}{2(4\pi \lambda )^2}\mathrm{ln}(\eta H_{c2}/H),$$ (1) where $`\lambda (T)\lambda _{ab}(T)`$ is the in-plane London penetration depth, $`\varphi _0`$ is the flux quantum, and $`\eta `$ is a parameter analogous to the Abrikosov ratio $`\beta _A=<\|\mathrm{\Psi }\|^4>/<\|\mathrm{\Psi }^2\|>^2(\mathrm{\Psi }`$ being the superconducting order parameter) that depends on the FLL structure. We assume now that the $`\mathrm{\Delta }M_{eq}`$ results from a change of the parameter $`\eta `$ due to a change in the vortex arrangement triggered by the interaction of the FLL with pinning centers. Hence, the magnetization jump can be written as $$\mathrm{\Delta }M_{eq}=\frac{\varphi _0}{2(4\pi \lambda )^2}\mathrm{ln}(\eta ^{dis}/\eta _0),$$ (2) where the parameter $`\eta ^{dis}\eta _0`$ is related to the strength of the quenched disorder, and $`\eta _0`$ applies for the FLL above the $`H_j(T)`$. Therefore, it is reasonable to assume that $`\eta ^{dis}`$ is proportional to the critical current density $`j_c(T,H)`$ which measures the vortex pinning strength. In fact, the correlation between $`\mathrm{\Delta }M_{eq}`$ and the pinning of the FLL has been observed experimentally. We emphasise that although the relationship between $`\eta `$ and $`\beta _A`$ is unknown, one expects an increase of the parameter $`\eta `$ in the vortex liquid state compared to that of the vortex solid, similar to the results obtained for $`\beta _A`$ . We show in Fig. 1 the jumps of the induction $`\mathrm{\Delta }B(T,H)`$ and of the magnetization $`4\pi \mathrm{\Delta }M(T,H)`$ measured along the $`H_j(T)`$ boundaries in Bi2212 (a) and Y123 (b) single crystals, respectively. The difference in the behavior of $`\mathrm{\Delta }M_{eq}(T,H)`$ in Bi2212 and Y123 crystals (see Fig. 1) can be easily understood noting that $`H_j(T)<H^{}=\varphi _0/\lambda ^2`$ for Bi2212, whereas $`H_j(T)H^{}`$ in the case of Y123. At fields $`H<H^{}`$ the FLL shear modulus exponentially decreases with field as $`c_{66}(ϵ_0/\lambda ^2)(H\lambda ^2/\varphi _0)^{1/4}\mathrm{exp}(\varphi _0/H\lambda ^2)`$, whereas at $`H>H^{}`$ it is proportional to the field $`c_{66}(ϵ_0/4\varphi _0)H`$ . The exponential decrease of $`c_{66}`$ with decreasing $`H_j`$ (increasing $`T_j`$) in the case of Bi2212, implies an enhancement of the interaction between vortices and the quenched disorder (pinning centers) which leads to an increase of $`\mathrm{\Delta }M_{eq}`$ increasing $`T_j`$ (or decreasing $`H_j`$), see Fig. 1(a). As temperature tends to the critical one $`T_c`$, the pinning of vortices vanishes. Therefore, above a certain temperature, $`\mathrm{\Delta }M_{eq}(T,H)`$ decreases with temperature and tends to zero. In the case of Y123, however, and due to a relatively weak field dependence of $`c_{66}`$ along the $`H_j(T)`$ line, the reduction of the vortex pinning efficiency with temperature is the dominant effect which explains the monotonous decrease of $`\mathrm{\Delta }M_{eq}(T,H)`$ with temperature, see Fig. 1(b). We stress that the vanishing of the magnetization jump $`(\mathrm{\Delta }M_{eq}(T,H))`$ at $`T_090`$K, i.e. approximately 3 K below the superconducting transition temperature $`T_c=92.9`$K (similar result was obtained in another untwinned Y123 crystal), can be also explained naturally by the effect of thermal fluctuations which smear out the pinning potential. On the other hand, the theory based on the FLL melting hypothesis predicts $`\mathrm{\Delta }M_{eq}\varphi _0/\lambda ^2(T)`$ (dotted line in Fig. 1(b)) which implies the vanishing of $`\mathrm{\Delta }M_{eq}`$ at $`T_c`$. While the FLL-melting theory requires different approaches in order to explain $`\mathrm{\Delta }M_{eq}(T,H)`$ in Y123 and Bi2212, our analysis can equally well be applied to both weakly (Y123) and strongly (Bi2212) anisotropic superconductors. In order to use Eq. (2) to calculate the magnetization jump we need the relationship between $`\eta ^{dis}`$ and the critical current density $`j_c(T,H)`$ at the boundary $`H_j(T)`$, which is unknown at present. Nevertheless and as an illustration of our ideas we present here a simple fitting approach. The solid lines in Fig. 1(a,b) were obtained from Eq. (2) with $`\eta ^{dis}=\eta _0[1+aj_c(T,H)]=\eta _0[1+aj_{c0}(1T/T_0)^n/(H^\alpha +H_0)]`$ calculated at the boundary $`H_j(T)`$, where $`T_0`$ corresponds to the temperature at which $`j_c=0`$, and $`a,j_{c0},H_0`$ are model-dependent constants, $`n`$ and $`\alpha `$ are pinning related exponents. Note, that the here used $`j_c(T,H)`$ is a rather general expression which reflects the well-known experimental fact that the critical current density generally decreases with temperature and increasing field. In our fits (solid lines in Fig. 1) we have set the exponents $`n=1,\alpha =1.8(1)`$, and used $`\lambda (T)=250(1T/T_c)^{1/3}`$nm ($`140(1T/T_c)^{1/3}`$ nm) for Bi2212 (Y123); other fitting parameters close to those used give also satisfactory fits. Furthermore, within the here proposed physical picture we expect a reduction and ultimately the vanishing of $`\mathrm{\Delta }M_{eq}(T,H)`$ if by applying external driving forces one de-couples the FLL from the matrix. Such effect was observed in Bi2212 crystals, indeed. Also, within our picture we expect that the FLL remains in a more ordered, metastable state if the sample is field cooled as compared to the zero-field-cooled state. Several published results have indicated such a behavior, both directly (see, e.g. Ref. ) and indirectly. Among them, the pioneer work which demonstrates that the magnetization jump at the irreversibility line can be much larger in the zero-field-cooled sample compared to the field-cooled one, providing a clear evidence that $`\mathrm{\Delta }M_{eq}(T,H)`$ is essentially related to the competition between vortex-vortex and vortex-pinning centers interactions. Supporting the above ideas, the magnetic-field-driven transition from a disorder-dominated vortex state to a moving well-ordered FLL was observed in NbSe<sub>2</sub> low-$`T_c`$ layered superconductor . Finally, we would like to point out that the equilibrium magnetization jump at a first-order depinning transition would imply the use of the Clausius-Clapeyron relation. However, we are not aware of any prediction on the entropy change at the depinning transition which we could use for comparison. Therefore, we believe that it has little sense to comment here on the use of the Clausius-Clapeyron equation at this stage of our study. To summarise, based essentially on experimental facts we propose the magnetic-field and temperature-driven vortex-lattice-ordering transition as an alternative to the FLL melting scenario in high-$`T_c`$ superconductors. Simple arguments allowed us to account for the equilibrium magnetization jump associated with the FLL depinning transition. This transition awaits for a rigorous theoretical treatment. We acknowledge illuminating discussions with G. Blatter, E. H. Brandt, A.M. Campbell, G. Carneiro, F. de la Cruz, M. Däumling, R. Doyle, E.M. Forgan, D. Fuchs, A. Gurevich, B. Horovitz, M. Konczykowski, A. H. MacDonald, A. Schilling, T. Tamegai, V. Vinokur, E. Zeldov, and M. Ziese. This work was supported by the German-Israeli-Foundation for Scientific Research and Development and the Deutsche Forschungsgemeinschaft.
warning/0002/cond-mat0002419.html	ar5iv	text	# Evidence of a structural anomaly at 14 K in polymerised CsC60 ## Abstract We report the results of a high-resolution synchrotron X-ray powder diffraction study of polymerised CsC<sub>60</sub> in the temperature range 4 to 40 K. Its crystal structure is monoclinic (space group $`I`$2/$`m`$), isostructural with RbC<sub>60</sub>. Below 14 K, a spontaneous thermal contraction is observed along both the polymer chain axis, $`a`$ and the interchain separation along , $`d_1`$. This structural anomaly could trigger the occurrence of the spin-singlet ground state, observed by NMR at the same temperature. PACS.61.48+c Fullerenes and fullerene-related materials. PACS.61.10Nz Single crystal and powder diffraction. Introduction. \- Alkali fullerides with stoichiometry AC<sub>60</sub> ($`A`$= K, Rb, Cs) undergo a structural transition from a high-temperature monomer to a polymer phase in the vicinity of 350 K . Polymerisation occurs by a \[2+2\] cycloaddition mechanism, leading to the formation of one-dimensional C-C bridged C$`{}_{}{}^{}{}_{60}{}^{}`$ chains . The AC<sub>60</sub> polymers exhibit interesting structural properties and form a variety of conducting and magnetic phases. A metal-insulator transition below 50 K is accompanied by the stabilisation of a magnetic state in RbC<sub>60</sub> and CsC<sub>60</sub>, whereas KC<sub>60</sub> remains metallic to low temperatures . The nature of the magnetic transition has remained controversial, as both quasi-one dimensional electronic instabilities and three-dimensional magnetic ordering have been proposed. Initial X-ray powder diffraction studies have described the structure of the polymers as orthorhombic (space group $`Pmnn`$) with the orientation of the C$`{}_{}{}^{}{}_{60}{}^{}`$ chains about the short axis $`a`$ described by $`\mu `$=45$`\pm `$5, where $`\mu `$ is the angle between the cycloaddition planes and the $`c`$ axis. However, recent single crystal X-ray diffraction and diffuse scattering studies on KC<sub>60</sub> and RbC<sub>60</sub> have revealed that they are not isostructural. They adopt different relative chain orientations (Figure 1) and their crystal structures are described by distinct space groups, $`Pmnn`$ (orthorhombic) and $`I`$2/$`m`$ (body-centred monoclinic), respectively. The difference in their respective physical properties can thus be attributed to the distinct relative chain orientations. Moreover the angle $`\mu `$ was determined as 51 for KC<sub>60</sub> and 47 for RbC<sub>60</sub>, the difference also affecting the electronic band structure . CsC<sub>60</sub> exhibits overall similar physical properties to RbC<sub>60</sub>. However, recent NMR measurements have detected the appearance of a spin-singlet ground state below $`T_S`$= 13.8 K which coexists with the long range ordered magnetic state . It was proposed that the development of the non-magnetic phase may be correlated either with a structural change, tentatively ascribed to the occurrence of a spin-Peierls transition or with an electronic instability. We note that earlier $`\mu `$SR data had described the low-temperature state of CsC<sub>60</sub> as comprising of co-existing static magnetic and fluctuating paramagnetic domains . In this paper, we report the crystallographic characterisation of the CsC<sub>60</sub> polymer by synchrotron X-ray powder diffraction in the temperature range 4-40 K. This reveals that CsC<sub>60</sub> is isostructural with RbC<sub>60</sub>, adopting a body-centred monoclinic structure (space group $`I`$2/$`m`$). In addition, our results show that a spontaneous strain appears along both the polymer chain axis and the interchain direction below $`T_S`$, providing the signature of magnetoelastic coupling. Experimental. \- The CsC<sub>60</sub> sample was prepared by solid state reaction of stoichiometric quantities of Cs metal and C<sub>60</sub> in sealed quartz tubes at 800 K for four weeks with intermittent shaking. High-resolution synchrotron X-ray powder diffraction measurements were performed on a CsC<sub>60</sub> sample sealed in a 0.5-mm diameter glass capillary. Data were collected in continuous scanning mode using nine Ge(111) analyser crystals on the BM16 beamline at the European Synchrotron Radiation Facility (ESRF), Grenoble, France in the temperature range 4-40 K ($`\lambda `$= 0.79972 Å). Data were rebinned in the 2$`\theta `$ range 5\- 45 to a step of 0.01. Analysis of the diffraction data was performed with the GSAS suite of powder diffraction programmes . In the course of the Rietveld refinements, the characteristic modulated background was fitted to a 15th-order Chebyshev polynomial. It arises from diffuse scattering, probably due to the static rotational disorder of residual C<sub>60</sub> molecules and/or to the different orientational domains arising from the symmetry lowering at the monomer- polymer transition. The initial refinements included the scale factor, the background coefficients, the lattice constants, the zero point and the peak width parameters. The final refinements incorporated the isotropic temperature factors, the positional parameters of the bridging C atoms and the occupation number of Cs. Results. \- We first performed detailed Rietveld refinements of the diffraction profile of CsC<sub>60</sub> at 20 K. We used three starting structural models, derived from those reported before for the AC<sub>60</sub> polymers. In all cases, the C$`{}_{}{}^{}{}_{60}{}^{}`$ ions are located at the (0,0,0) and ($`\frac{1}{2}`$,$`\frac{1}{2}`$,$`\frac{1}{2}`$) positions, while the Cs<sup>+</sup> ions at the (0,0,$`\frac{1}{2}`$) and ($`\frac{1}{2}`$,$`\frac{1}{2}`$,0) positions in the unit cell. However, in each case a different polymer chain orientation ordering is adopted, as illustrated in Figure 1. For space group $`Pmnn`$, the orientation of the chains at the origin and the centre of the unit cell are $`\mu `$ and -$`\mu `$, respectively, while for $`I`$2/$`m`$, they are identical ($`\mu `$). The third model is described by space group $`Immm`$ and involves a disordered arrangement of chains with $`\mu `$ or -$`\mu `$ ($``$0 or 90) orientations. In the course of the refinements, we monitored the resulting quality-of-fit factors ($`R_{wp}`$) as a function of the orientation of the chains for all three space groups. Figure 2 presents the evolution of $`R_{wp}`$ with the rotation angle $`\mu `$ which was varied between 40 and 50. The refinements were stable throughout this $`\mu `$-range with a minimum in all cases at $`\mu `$46. However, the deepest minimum is clearly obtained for the monoclinic $`I`$2/$`m`$ space group, in a similar fashion to RbC<sub>60</sub> . Subsequent refinements concentrated on this structural model and resulted in excellent quality fits ($`R_{wp}`$= 3.97%, $`R_{exp}`$= 1.44%). The final results are plotted in Figure 3. The resulting lattice parameters at 20 K are $`a`$= 9.0968(3) Å, $`b`$= 10.1895(3) Å, $`c`$= 14.1351(4) Å and $`\alpha `$= 89.820(6). The intra- and inter-molecular C<sub>1</sub>-C<sub>1</sub> bond distances also refined to 1.72(1) Å and 1.60(1) Å, respectively, where C<sub>1</sub> is the bridging C atom between two C<sub>60</sub> molecules. Following the successful determination of the structure of the CsC<sub>60</sub> polymer at 20 K, we attempted Rietveld refinements of the datasets at other temperatures. In the course of these refinements, the positional and thermal atomic parameters and the peak width functions were kept fixed. The resulting variation of the monoclinic lattice constants in the temperature range 4-40 K is strongly anisotropic (Figure 4). While the lattice constants, $`b`$ and $`c`$ vary smoothly with temperature, a clear anomaly occurs below $``$14 K along the polymer chain axis $`a`$. A spontaneous contraction of the $`a`$ constant, $`\mathrm{\Delta }a/a`$ 1.5$`\times `$10<sup>-4</sup> is evident between 14 and 4 K. At the same time, the monoclinic angle $`\alpha `$ also increases by roughly the same proportion, leading to an overall contraction in the volume of the unit cell of $`\mathrm{\Delta }V/V`$ 2$`\times `$10<sup>-4</sup>. The observed anomaly is also mirrored in the interfullerene chain separations. In space group $`I`$2/$`m`$, there are two types of intermolecular environments, $`d_1`$ along the and $`d_2`$ along the \[1$`\overline{1}`$1\] directions of the unit cell. These are not affected equally below 14 K (Figure 5) with a sharp decrease in interchain separation occurring only for $`d_1`$ along ($`\mathrm{\Delta }d_1/d_1`$ 1.0$`\times `$10<sup>-4</sup>). Discussion. \- The temperature-dependent synchrotron X-ray data have shown that the low-temperature structure of CsC<sub>60</sub> is identical to that of RbC<sub>60</sub> with similar orientational angles $`\mu `$. Thus both polymers have comparable electronic structures and show the appearance of a low-temperature magnetic insulating phase, in contrast to the non-isostructural KC<sub>60</sub> whose metallic state is robust and shows no low-temperature instabilities. The difference in size between K<sup>+</sup>, Rb<sup>+</sup> and Cs<sup>+</sup> also leads to increased nearest-neighbour separations between the polymer chains (succesively by $``$ 0.11% and 0.12%) and increased quasi-one-dimensional character of the electronic structure, as we progress from KC<sub>60</sub> to RbC<sub>60</sub> to CsC<sub>60</sub>. In the related family of fulleride polymers, Na<sub>2</sub>Rb<sub>1-x</sub>Cs<sub>x</sub>C<sub>60</sub> (0$`x`$ 1), the electronic properties were also shown to be very sensitive to such subtle size effects . The details of the electronic structure of the RbC<sub>60</sub> and CsC<sub>60</sub> polymers remain controversial. While many experiments reveal one-dimensional characteristics , electronic band structure calculations predict three-dimensional energy bands . In addition, recent <sup>13</sup>C magic angle spinning NMR experiments on CsC<sub>60</sub> have shown that the conduction electron density is concentrated along the equator of the C$`{}_{}{}^{}{}_{60}{}^{}`$ ions and away from the C-C bridging bonds . This indicates that the band structure is dominated by strong transverse coupling between the polymer chains rather than a strong one-dimensional coupling along them. Nonetheless, <sup>13</sup>C and <sup>133</sup>Cs NMR measurements also revealed a transition to a non-magnetic (spin-singlet) ground state at $`T_S`$= 13.8 K, ascribed to a structural or electronic instability . The former was interpreted as a spin-Peierls transition, consistent with the presence of strong 1D features in the electronic description of the CsC<sub>60</sub> polymer. A highly significant result of the present diffraction experiments is the observation of a spontaneous thermal contraction along both the chain axis, $`a`$ and the interchain separation along , $`d_1`$ below a transition temperature of about 14 K, which coincides with the spin-singlet state transition temperature, $`T_S`$ and unambiguously points to the structural origin of the instability. This result is certainly reminiscent of the situation encountered for the linear chain compound, CuGeO<sub>3</sub> in which the magnetic transition to a spin-Peierls state is accompanied by shifts of the Cu<sup>2+</sup> ions along the chain direction and oxygen displacements perpendicular to the chains ; the magnetoelastic coupling was evident in diffraction experiments as a spontaneous thermal contraction along the axis $`b`$, perpendicular to the chain direction, of comparable magnitude to that observed for CsC<sub>60</sub>. However, no evidence of any superstructure peaks is established in CsC<sub>60</sub>, while the quality of the powder diffraction data is not such as to allow us to determine whether the spontaneous strain is accompanied by a structural phase transition. In this respect, it is also interesting to note that the $`b`$ and $`c`$ axes parameters of CsC<sub>60</sub> and the interchain distance along \[1$`\overline{1}`$1\], $`d_2`$ show no anomalies at $`T_S`$. The current structural evidence is certainly consistent with the occurrence of a structural transition in CsC<sub>60</sub>. However, while a spin-Peierls scenario is possible, it is equally likely that the observed structural changes may arise from a CDW (commensurate or incommensurate) transition. An appealing simple interpretation of the structural anomaly at 14 K could be that a soft phonon mode along the polymer chain may be responsible for the lattice contraction along $`a`$ which then induces a local structural rearrangement through the variation of the interchain distance $`d_1`$. In this respect, it would be of interest to perform high-resolution structural measurements at elevated pressures. The intermolecular distance $`d_1`$ should vary the most, while contraction along $`a`$ should be limited because of the rigidity of the polymer chain. Band structure calculations should then give valuable insight on the influence of the $`d_1`$ variation on the nature of the non-magnetic phase, especially as it has been found experimentally that application of pressure gradually supresses the magnetic order in favour of the spin-singlet state in both CsC<sub>60</sub> and RbC<sub>60</sub> . Conclusions. \- In conclusion, high-resolution synchrotron X-ray powder diffraction has established that the structure of CsC<sub>60</sub> is monoclinic (space group $`I`$2/$`m`$), isostructural with that of RbC<sub>60</sub>. The similarity in the orientational ordering of the polymer chains in these systems explains their comparable electronic and conducting properties. A structural anomaly is observed below $``$14 K, coincident with the appearance of a non-magnetic state. The observed structural changes provide the signature of magnetoelastic coupling between the $`C_{60}^{}`$ localised spins and the phonon degrees-of-freedom. \*** We thank the European Union (TMR Research Network ‘FULPROP’, ERBFMRXCT970155) and the NEDO Frontier Carbon Technology Program for financial support and the ESRF for provision of synchrotron X-ray beam time.
warning/0002/nucl-th0002002.html	ar5iv	text	# QCD sum rules for 𝐽/𝜓 in the nuclear medium: calculation of the Wilson coefficients of gluon operators up to dimension 6 ## 1 Introduction Identifying higher twist operators and calculating their corresponding Wilson coefficients are very important in several aspects. First, these provide a systematic building blocks to analyze the available data from deep inelastic scattering(DIS) at lower $`Q^2`$ region. Second, these contributions are essential for a realistic generalization of QCD sum rule methods to finite baryon density. The twist-4 operators were classified and its anomalous dimensions for some of the operators were calculated by S. Gottlieb. However, the operators were over-determined and the independent set of twist-4 operators appearing in the DIS were first identified by Jaffe and Soldate, where the twist-4 operators are of four quark type and quark gluon mixed type. Among the uses of these results, the available DIS data were analyzed to determine the nucleon matrix elements of the twist-4 operators. These estimates were then used in generalizing the QCD sum rule approaches for light vector mesons to finite density. In the correlation function of two heavy vector currents, only gluon operators contribute in the operator product expansion(OPE). This is so because in the heavy quark system, all the heavy quark condensates are generated via gluonic condensates. In dimension 6, there are scalar operators, twist-2 and twist-4 operators. For the scalar gluonic operators at dimension 6, there are two independent operators. In ref., the two were identified and the corresponding Wilson coefficient were calculated . For twist-2 gluon operator, the calculation for the leading order(LO) Wilson coefficient is simple and its matrix element is just the second moment of the gluon structure function. The twist-4 dimension 6 gluon operators are more involved. In this work, we have identified the three independent local gluon operators and calculated their corresponding LO Wilson coefficients in the correlation function between two heavy vector currents. This result is new and complimentary to a previous work on gluon twist-4 operators, where they start from certain diagrams and identify the three independent twist-4 gluon structure functions. Together with the previous calculation of the Wilson coefficients for the dimension 6 scalar gluon operators by Nikolaev and Radyushkin, our result completes the list of all the Wilson coefficients of dimension 6 gluon operators in the correlation function between heavy vector currents. As an application, we will use our result in QCD sum rule approach to investigate the property of $`J/\psi `$ in nuclear matter. This is particularly interesting because the on-going discussion of $`J/\psi `$ suppression in RHIC as a possible signal for quark gluon plasma, inevitably requires a detailed knowledge of the changes of $`J/\psi `$ properties in “normal” nuclear matter. Furthermore, the large charm quark mass $`m_c\mathrm{\Lambda }_{QCD}`$ provides a natural renormalization point for which a perturbative QCD expansion is partly possible. In fact, studies have shown that the multi-gluon exchange between a $`c\overline{c}`$ pair and nucleons might induce a bound $`c\overline{c}`$ state with even light nuclei. In such analysis, the low energy multi-gluon potential was modeled either from the effective theory obtained in the infinitely large $`m_Q`$ limit or from extrapolating the high energy scattering via pomeron exchange to lower energy. Although both approaches, gave similar binding for the $`J/\psi `$ in nuclear matter, it is not clear how reliable these results are unless one systematically calculates the corrections. In order to confirm this findings in an alternative but a systematic approach and to establish a basis for further studies, we have previously applied the QCD sum rules to heavy quark system in nuclear medium and calculated the mass of $`J/\psi `$ and $`\eta _c`$ in nuclear medium. This was the generalization of the sum rule method for the light vector mesons in medium to the heavy quark system. It was found that the mass of $`J/\psi `$ ($`\eta _c`$) would reduce by 7 MeV (5 MeV), which is indeed consistent with previous results based on completely different methods. However, it was not possible to reliably estimate the uncertainties of the result, because the contribution from the operator product expansion was truncated at the leading dimension 4 operators. To overcome the limitations , we will here make use of our calculation to include the complete dimension 6 contributions. Unfortunately, at present, there is no data to identify the nucleon expectation value of any of the dimension 6 operators(scalar, twist-2, twist-4) nor is there any lattice result. Nevertheless, in this article, we will estimate the nucleon matrix elements of these operators and then use QCD sum rule approach to study the reliability of our previous result on the mass shift of $`J/\psi `$ in nuclear matter. If in the future, the matrix elements are better known, we will be able to study the non trivial momentum dependence of the $`J/\psi `$ in nuclear medium, as has been done for the light vector mesons. The momentum dependence is especially interesting because, there are inelastic channels opening when the charmonium system is moving with respect to the medium. For example, when the charmonium is moving with sufficient velocity, it will have enough energy to interact with a nucleon to produce a D meson and a charmed nucleon. This effect would be very important in relation to $`J/\psi `$ suppression in RHIC. In section 2 we characterize all gluon operator up to dimension 6 and obtain identities to be used to reduce the operators to an independent set. In section 3, we calculate the Wilson coefficients for the independent set of gluon operators and show current conservation. In section 4, we calculate the moments and perform a moment sum rules analysis to calculate the $`J/\psi `$ mass in nuclear matter. We conclude with some discussions. The appendix includes some detailed calculation of the Wilson coefficients. ## 2 Operators In the operator product expansion of heavy quark system, only gluonic condensates are relevant. This is so because all the heavy quark condensates can be related to the gluon condensates via heavy quark expansion . This is also true in nuclear medium since there are no valence charm quarks to leading order in density and any interaction with the medium is gluonic. Let us start by categorizing gluonic operators up to dimension 6, which does not vanish in nuclear matter. For dimension 4 operators, the scalar and twist-2 gluon operators contribute , $`g^2G_{\mu \nu }^aG_{\mu \nu }^a,g^2G_{\mu \alpha }^aG_{\nu \alpha }^a.`$ (1) For dimension 6 operators, one can think of generating a number of gluon operators constructed from three gluon fields $`G_{\mu \nu }^a`$ or two gluon fields with two covariant derivatives. However, for scalar operators, there are only two independent scalar operators. They are, $`g^3f^{abc}G_{\mu \nu }^aG_{\mu \alpha }^bG_{\nu \alpha }^c,g^2G_{\mu \alpha }^aG_{\nu \alpha ;\nu \mu }^a.`$ (2) The second operator can also be written in terms of four quark operator using the equation of motion $`G_{\mu \nu ;\nu }^a=g\overline{q}\gamma _\mu {\displaystyle \frac{\lambda ^a}{2}}q=gj_\mu ^a,`$ (3) As for the spin 2 operators in dimension 6, which are also called dimension 6 twist 4 (twist=dimension-spin) operator, we can first categorize possible operators as follows. First, there is again one operator with three gluon field strength tensor. Then, assuming the free symmetric and traceless indices to be $`\mu `$ and $`\nu `$, depending on whether or not the free index goes into the covariant derivative, there are 6 more operators possible. Starting with the three gluon operator, the 7 are, $`g^3f^{abc}G_{\mu \kappa }^aG_{\nu \lambda }^bG_{\kappa \lambda }^cg^3fG_{\mu \nu }^3`$ $`g^2G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a`$ $`g^2G_{\mu \kappa }^aG_{\nu \lambda ;\kappa \lambda }^a,g^2G_{\mu \kappa }^aG_{\nu \kappa ;\lambda \lambda }^a,g^2G_{\mu \kappa }^aG_{\nu \lambda ;\lambda \kappa }^a`$ $`g^2G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a,g^2G_{\kappa \lambda }^aG_{\mu \kappa ;\lambda \nu }^a.`$ (4) However, they are not independent. Using the identities in appendix A, one can show that there are three independent spin-2 operators at dimension 6. The set we will use are, $`g^2G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a,g^2G_{\mu \kappa }^aG_{\nu \lambda ;\lambda \kappa }^a,g^2G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a.`$ (5) It is interesting to note that in reference starting from certain diagrams with two, three and four gluon exchange in the t-channel, they were able to derive three twist-4 gluon distribution amplitudes, from which one can calculate the nucleon matrix element of the three independent operators. As for the spin 4 dimension 6 operator, it is just the twist-2 gluon operator. ## 3 Polarization (OPE) Having established the independent gluon operators in dimension 6, we will calculate their LO Wilson coefficients in the correlation function between two vector currents made of heavy quarks, $`j_\mu =\overline{h}\gamma _\mu h`$. $`\mathrm{\Pi }_{\mu \nu }(q)`$ $`=`$ $`i{\displaystyle d^4xe^{iqx}\mathrm{T}\{j_\mu (x)j_\nu (0)\}_\rho }`$ (6) $`=`$ $`i{\displaystyle }d^4xe^{iqx}\mathrm{Tr}[\gamma _\mu S(x,0)\gamma _\nu S(0,x))]_\rho `$ $`=`$ $`i{\displaystyle \frac{d^4k}{(2\pi )^4}\mathrm{Tr}[\gamma _\mu S(k+q)\gamma _\nu \stackrel{~}{S}(k)]_\rho },`$ where $`_\rho `$ represents the expectation value at finite nuclear density $`\rho `$. The fourier transforms are defined by $`iS(p)={\displaystyle d^4xe^{ipx}iS(x,0)}`$ $`i\stackrel{~}{S}(p)={\displaystyle d^4xe^{ipx}iS(0,x)}.`$ (7) $`S(x,0)`$ is the heavy quark propagator in the presence of external gluon field . To calculate the Wilson coefficients, we obtain the quark propagator in the presence of external gauge fields in the Fock-Schwinger gauge , $`x_\mu A_\mu ^a(x)=0.`$ (8) In appendix C we list the momentum space representation of the quark propagator multiplying the gauge invariant gluon operators. We will use these quark propagators in eq.(6) and extract the Wilson coefficients by collecting the appropriate tensor and gluon structures. In general the polarization tensor in eq.(6) will have two invariant functions. They can be divided into the longitudinal and transverse part of the external three momentum, which are the following when the nuclear matter is at rest. $`\mathrm{\Pi }_L={\displaystyle \frac{1}{k^2}}\mathrm{\Pi }_{00},\mathrm{\Pi }_T={\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{q^2}}\mathrm{\Pi }_\mu ^\mu +{\displaystyle \frac{1}{k^2}}\mathrm{\Pi }_{00}),`$ (9) where $`q=(\omega ,k)`$ is the external momentum. In the vacuum or in the limit when $`k0`$ they become the same so that there is only one invariant function $`\mathrm{\Pi }_L(\omega ^2)=\mathrm{\Pi }_T(\omega ^2)={\displaystyle \frac{1}{3\omega ^2}}g^{\mu \nu }\mathrm{\Pi }_{\mu \nu }\mathrm{\Pi }(\omega ^2).`$ (10) So in this work, we will construct the sum rule for $`\mathrm{\Pi }(\omega ^2)`$ at $`k0`$ and nuclear matter at rest. However, in the calculation of the OPE, we will start from the general expression of eq.(6) at nonzero value of $`k`$ and calculate each tensor structure separately. This way, current conservation will be a nontrivial check of our calculation and the generalization to $`k0`$ will be straightforward. In the following subsections, we will summarize the OPE for operators of dimension 4 and 6. The new results of the present work are the OPE for dimension 6 spin 2 and dimension 6 spin 4 operators. ### 3.1 scalar contributions Here we summarize the OPE of scalar operators of dimension 4 and dimension 6 used in the vacuum sum rule. $`\mathrm{\Pi }_{\mu \nu }^{\mathrm{scalar}}(q)`$ $`=`$ $`(q_\mu q_\nu g_{\mu \nu }q^2)[C^{pert.}+C_{G^2}^0g^2G_{\alpha \beta }^aG_{\alpha \beta }^a`$ (11) $`+C_{G^3}^0g^3f^{abc}G_{\alpha \beta }^aG_{\beta \lambda }^bG_{\lambda \alpha }^c+C_{j^2}^0g^4j_\kappa ^aj_\kappa ^a],`$ where $`C_{G^2}^0={\displaystyle \frac{1}{48\pi ^2\left(Q^2\right)^2}}\left[1+3J_22J_3\right],`$ $`C_{G^3}^0={\displaystyle \frac{1}{72\pi ^2\left(Q^2\right)^3}}\left[{\displaystyle \frac{2}{15}}{\displaystyle \frac{1}{10}}y+4J_2{\displaystyle \frac{31}{3}}J_3+{\displaystyle \frac{43}{5}}J_4{\displaystyle \frac{12}{5}}J_5\right],`$ $`C_{j^2}^0={\displaystyle \frac{1}{36\pi ^2\left(Q^2\right)^3}}\left[{\displaystyle \frac{41}{45}}{\displaystyle \frac{3}{5}}y+\left({\displaystyle \frac{2}{3}}+{\displaystyle \frac{1}{3}}y\right)J_1J_2{\displaystyle \frac{4}{9}}J_3{\displaystyle \frac{26}{15}}J_4+{\displaystyle \frac{8}{5}}J_5\right],`$ (12) and $`J_N(y)={\displaystyle _0^1}{\displaystyle \frac{dx}{\left[1+x(1x)y\right]^N}},`$ (13) with $`y=Q^2/m^2`$ and $`m`$ being the heavy quark mass. This will give the following contribution to $`\mathrm{\Pi }(\omega ^2)`$ defined in eq.(10). $`\mathrm{\Pi }(\omega ^2)`$ $`=`$ $`C^{pert.}+C_{G^2}^0g^2G_{\alpha \beta }^aG_{\alpha \beta }^a+C_{G^3}^0g^3f^{abc}G_{\alpha \beta }^aG_{\beta \lambda }^bG_{\lambda \alpha }^c`$ (14) $`+C_{j^2}^0g^4j_\kappa ^aj_\kappa ^a.`$ For later convenience, we will use $`g^3f^{abc}G_{\mu \nu }^aG_{\nu \lambda }^bG_{\lambda \mu }^c=\frac{1}{2}g^2G_{\mu \nu }^aG_{\mu \nu ;\alpha \alpha }^a+g^4j_\mu ^aj_\mu ^a`$ and rewrite the OPE as follows, $`\mathrm{\Pi }(\omega ^2)`$ $`=`$ $`C^{pert.}+C_{G^2}{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \nu }^aG_{\mu \nu }^a+C_{GD^2G}{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \nu }^aG_{\mu \nu ;\kappa \kappa }^a`$ (15) $`+C_{j^2}{\displaystyle \frac{2}{3}}{\displaystyle \frac{\alpha _s}{\pi }}G_{\alpha \kappa }^aG_{\alpha \lambda ;\lambda \kappa }^a,`$ where $`C_{G^2}`$ $`=`$ $`{\displaystyle \frac{1}{12(Q^2)^2}}[1+3J_22J_3],`$ $`C_{GD^2G}`$ $`=`$ $`{\displaystyle \frac{1}{1080(Q^2)^3}}[43y+120J_2310J_3+258J_472J_5],`$ $`C_{j^2}`$ $`=`$ $`{\displaystyle \frac{1}{1080(Q^2)^3}}[(43y+120J_2310J_3+258J_472J_5)`$ (16) $`180+120y(120+60y)J_1300J_2+1320J_3720J_4].`$ Note $`Q^2=\omega ^2`$ in our kinematical limit. ### 3.2 dimension 4 spin 2 operator This part has been calculated recently . Here we summarize the results. $`\mathrm{\Pi }_{\mu \nu }^{4,2}(q)`$ $`=`$ $`{\displaystyle \frac{1}{Q^2}}[_{\mu \nu }^2+{\displaystyle \frac{1}{Q^2}}(q_\rho q_\mu _{\rho \nu }^2+q_\rho q_\nu _{\rho \mu }^2)`$ (17) $`+g_{\mu \nu }{\displaystyle \frac{q_\rho q_\sigma }{Q^2}}𝒥_{\rho \sigma }^2+{\displaystyle \frac{q_\mu q_\nu q_\rho q_\sigma }{Q^4}}(_{\rho \sigma }^2+𝒥_{\rho \sigma }^2)],`$ where $`_{\mu \nu }^2`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}G_{\sigma \mu }^aG_{\sigma \nu }^a\left[{\displaystyle \frac{1}{2}}+(1{\displaystyle \frac{1}{3}}y)J_1{\displaystyle \frac{3}{2}}J_2\right]`$ $`𝒥_{\mu \nu }^2`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}G_{\sigma \mu }^aG_{\sigma \nu }^a\left[{\displaystyle \frac{7}{6}}+(1+{\displaystyle \frac{1}{3}}y)J_1{\displaystyle \frac{1}{2}}J_2+{\displaystyle \frac{2}{3}}J_3\right].`$ (18) ### 3.3 contribution from dimension 6 and spin 4 The diagrams describing interactions with the gluonic field in dimension 6 are shown in Fig 1. The calculation for this part is straightforward. We substitute into eq.(6) parts of the quark propagator containing $`D`$’s and $`G`$’s (summarized in appendix C) so that when these from the two propagators are combined, yield terms proportional to $`GDDG`$. In Fig.1, these come from Fig.1(c) to Fig.1(g). Then we extract the traceless and symmetric spin-4 part of the operator $`GDDG`$. The final result for the dimension 6 spin 4 part of the operators yields, $`\mathrm{\Pi }_{\mu \nu }^{6,4}(q)`$ $`=`$ $`{\displaystyle \frac{q_\kappa q_\lambda }{(Q^2)^3}}[I_{\kappa \lambda \mu \nu }^4+{\displaystyle \frac{1}{Q^2}}(q_\rho q_\mu I_{\kappa \lambda \rho \nu }^4+q_\rho q_\nu I_{\kappa \lambda \rho \mu }^4)`$ (19) $`+g_{\mu \nu }{\displaystyle \frac{q_\rho q_\sigma }{Q^2}}J_{\kappa \lambda \rho \sigma }^4+{\displaystyle \frac{q_\mu q_\nu q_\rho q_\sigma }{Q^4}}(I_{\kappa \lambda \rho \sigma }^4+J_{\kappa \lambda \rho \sigma }^4)],`$ where $`I_{\mu \nu \rho \sigma }^4`$ $`=`$ $`\left[{\displaystyle \frac{266}{45}}+\left({\displaystyle \frac{20}{3}}+{\displaystyle \frac{22}{15}}y\right)J_1+{\displaystyle \frac{138}{5}}J_2{\displaystyle \frac{916}{45}}J_3+{\displaystyle \frac{16}{3}}J_4\right]`$ $`\times {\displaystyle \frac{\alpha _s}{\pi }}G_{\rho \kappa }^aG_{\sigma \kappa ;\mu \nu }^a`$ $`J_{\mu \nu \rho \sigma }^4`$ $`=`$ $`\left[{\displaystyle \frac{362}{45}}\left(4+{\displaystyle \frac{22}{15}}y\right)J_1{\displaystyle \frac{94}{15}}J_2{\displaystyle \frac{44}{45}}J_3+{\displaystyle \frac{16}{3}}J_4{\displaystyle \frac{32}{15}}J_5\right]`$ (20) $`\times {\displaystyle \frac{\alpha _s}{\pi }}G_{\rho \kappa }^aG_{\sigma \kappa ;\mu \nu }^a.`$ ### 3.4 contributions from dimension 6 and spin 2 The calculation for this part goes in a similar way as in the previous subsection. However, this time all the graphs in Fig.1 contribute. Moreover, we have to include all possible combination of $`G`$ and $`D`$’s that makes up dimension 6. We then extract the spin 2 part of each operator and use the identities given in the appendix to reduce the operators into the independent three in eq.(5). An example needed to extract the spin 2 part from a general operator is given in appendix D. Finally we find, $`\mathrm{\Pi }_{\mu \nu }^{6,2}(q)`$ $`=`$ $`{\displaystyle \frac{1}{(Q^2)^2}}[I_{\mu \nu }^2+{\displaystyle \frac{1}{Q^2}}(q_\rho q_\mu I_{\rho \nu }^2+q_\rho q_\nu I_{\rho \mu }^2)`$ (21) $`+g_{\mu \nu }{\displaystyle \frac{q_\rho q_\sigma }{Q^2}}J_{\rho \sigma }^2+{\displaystyle \frac{q_\mu q_\nu q_\rho q_\sigma }{Q^4}}(I_{\rho \sigma }^2+J_{\rho \sigma }^2)],`$ where $`I_{\mu \nu }^2`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a\left[{\displaystyle \frac{31}{240}}{\displaystyle \frac{1}{60}}y+\left({\displaystyle \frac{13}{24}}+{\displaystyle \frac{1}{48}}y\right)J_1{\displaystyle \frac{115}{48}}J_2+{\displaystyle \frac{21}{8}}J_3{\displaystyle \frac{9}{10}}J_4\right]`$ $`+{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\nu \lambda ;\lambda \kappa }^a\left[{\displaystyle \frac{739}{720}}+{\displaystyle \frac{2}{15}}y+\left({\displaystyle \frac{9}{8}}+{\displaystyle \frac{3}{16}}y\right)J_1+{\displaystyle \frac{133}{48}}J_2+{\displaystyle \frac{1}{72}}J_3{\displaystyle \frac{19}{30}}J_4\right]`$ $`+{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a\left[{\displaystyle \frac{293}{240}}{\displaystyle \frac{3}{10}}y+\left({\displaystyle \frac{55}{24}}+{\displaystyle \frac{1}{16}}y\right)J_1{\displaystyle \frac{131}{16}}J_2+{\displaystyle \frac{145}{24}}J_3{\displaystyle \frac{41}{30}}J_4\right]`$ $`J_{\mu \nu }^2`$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a\left[{\displaystyle \frac{103}{240}}+{\displaystyle \frac{1}{60}}y+\left({\displaystyle \frac{5}{24}}{\displaystyle \frac{7}{48}}y\right)J_1{\displaystyle \frac{59}{48}}J_2+{\displaystyle \frac{31}{24}}J_3{\displaystyle \frac{11}{10}}J_4+{\displaystyle \frac{2}{5}}J_5\right]`$ (22) $`+{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\nu \lambda ;\lambda \kappa }^a\left[{\displaystyle \frac{71}{240}}{\displaystyle \frac{2}{15}}y+\left({\displaystyle \frac{1}{8}}+{\displaystyle \frac{1}{48}}y\right)J_1+{\displaystyle \frac{61}{48}}J_2{\displaystyle \frac{61}{24}}J_3+{\displaystyle \frac{29}{30}}J_4+{\displaystyle \frac{2}{15}}J_5\right]`$ $`+{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a\left[{\displaystyle \frac{29}{240}}+{\displaystyle \frac{3}{10}}y\left({\displaystyle \frac{1}{24}}+{\displaystyle \frac{7}{16}}y\right)J_1+{\displaystyle \frac{31}{48}}J_2{\displaystyle \frac{23}{24}}J_3+{\displaystyle \frac{11}{30}}J_4{\displaystyle \frac{2}{15}}J_5\right].`$ There are few checks we can perform to verify our calculation. First, although the polarization function in eq.(6) has only two invariant parts, we performed the calculation directly leaving the indices $`\mu ,\nu `$ free. This means that we have obtained and calculated the 6 different possible tensor structure, given in appendix B, separately. Therefore, showing current conservation, $`q^\mu \mathrm{\Pi }_{\mu \nu }=q^\nu \mathrm{\Pi }_{\mu \nu }=0`$ is a non-trivial check. Another indirect check is that the final result is regular at $`Q^2=0`$. That is, making a small $`Q^2`$ expansion one can show that all the Wilson coefficients in eq.(3.2), eq.(3.3) and eq.(3.4) are regular. Eq.(11), eq.(17), eq.(19) and eq.(21) form the complete OPE up to dimension 6 operators. Since QCD is renormalizable, there are no other power correction up to this dimension and the Wilson coefficients will be an asymptotic expansion in $`1/\mathrm{log}Q^2`$. In the vacuum, only scalar operators will contribute. However, in medium or when the expectation value is with respect to a nucleon state, the tensor operators will also contribute. As an application of our result, we will apply our OPE to the QCD sum rule analysis for the $`J/\psi `$ in nuclear matter. ### 3.5 sum of tensor parts As we discussed before we will take the nuclear matter to be at rest. Moreover, we will use linear density approximation to evaluate the matrix elements, $`_\rho =_0+{\displaystyle \frac{\rho }{2m_N}}N\|\|N`$ (23) where $`_0`$ is the vacuum expectation value, $`\rho `$ is the nuclear density, $`m_N`$ the nucleon mass, and the nucleon state $`\|N`$ is normalized as $`N\|N=(2\pi )^32\omega _N\delta ^3(0)`$. Then, the OPE from dimension 4 and dimension 6 operators in eq.(17), eq.(19) and eq.(21) give the following contribution to $`\mathrm{\Pi }`$. $`\mathrm{\Pi }(\omega ^2)`$ $`=`$ $`\mathrm{\Pi }_{4,2}(\omega ^2)+\mathrm{\Pi }_{6,2}(\omega ^2)+\mathrm{\Pi }_{6,4}(\omega ^2)`$ (24) $`=`$ $`{\displaystyle \frac{\rho }{2m_N}}\left[C_{G_2}G_2+(C_XX+C_YY+C_ZZ)+C_{G_4}G_4\right],`$ where $`C_{G_2}`$ $`=`$ $`{\displaystyle \frac{m_N^2}{12\left(Q^2\right)^2}}\left[9\left(12+2y\right)J_1+9J_26J_3\right],`$ $`C_X`$ $`=`$ $`{\displaystyle \frac{m_N^2}{240\left(Q^2\right)^3}}\left[852y+\left(70+25y\right)J_1+365J_2390J_3+252J_472J_5\right],`$ $`C_Y`$ $`=`$ $`{\displaystyle \frac{m_N^2}{720\left(Q^2\right)^3}}\left[25+48y+\left(27045y\right)J_11185J_2+1370J_3408J_472J_5\right],`$ $`C_Z`$ $`=`$ $`{\displaystyle \frac{m_N^2}{240\left(Q^2\right)^3}}\left[9536y+\left(130+75y\right)J_1+375J_2190J_3+16J_4+24J_5\right],`$ $`C_{G_4}`$ $`=`$ $`{\displaystyle \frac{m_N^4}{108\left(Q^2\right)^3}}\left[205\left(210+33y\right)J_1+99J_2262J_3+240J_472J_5\right],`$ (25) where again in our kinematical limit $`Q^2=\omega ^2`$ and the scalar parts of the matrix elements in nuclear matter to leading density come from the following nucleon expectation values, $`N\|{\displaystyle \frac{\alpha _s}{\pi }}G_{\sigma \mu }^aG_{\sigma \nu }^a\|N`$ $`=`$ $`G_2(p_\mu p_\nu {\displaystyle \frac{1}{4}}m_N^2g_{\mu \nu }),`$ $`N\|{\displaystyle \frac{\alpha _s}{\pi }}G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a\|N`$ $`=`$ $`X(p_\mu p_\nu {\displaystyle \frac{1}{4}}m_N^2g_{\mu \nu }),`$ $`N\|{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\nu \lambda ;\lambda \kappa }^a\|N`$ $`=`$ $`Y(p_\mu p_\nu {\displaystyle \frac{1}{4}}m_N^2g_{\mu \nu }),`$ $`N\|{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a\|N`$ $`=`$ $`Z(p_\mu p_\nu {\displaystyle \frac{1}{4}}m_N^2g_{\mu \nu }),`$ $`N\|{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \kappa }^aG_{\nu \kappa ;\alpha \beta }^a\|N`$ $`=`$ $`G_4[p_\mu p_\nu p_\alpha p_\beta +{\displaystyle \frac{m_N^4}{48}}(g_{\mu \nu }g_{\alpha \beta }+g_{\mu \alpha }g_{\nu \beta }+g_{\mu \beta }g_{\nu \alpha })`$ $`{\displaystyle \frac{1}{8}}m_N^2(p_\mu p_\nu g_{\alpha \beta }+p_\mu p_\alpha g_{\nu \beta }+p_\mu p_\beta g_{\alpha \nu }`$ $`+p_\nu p_\alpha g_{\mu \beta }+p_\nu p_\beta g_{\mu \alpha }+p_\alpha p_\beta g_{\mu \nu })].`$ Note that here we chose the nucleon four momentum to be $`p=(m_N,0,0,0)`$. We discuss the magnitudes of these nucleon matrix elements in section 4. ## 4 Moment sum rule for $`J/\psi `$ at rest ### 4.1 moments As an application of our result, we will use our result to refine our previous work to calculate the mass shift of the $`J/\psi `$ in nuclear medium using the moment sum rule . The moments of the polarization function is defined as, $`M_n(Q_0^2)={\displaystyle \frac{1}{n!}}\left({\displaystyle \frac{d}{dQ^2}}\right)^n\mathrm{\Pi }(Q^2)\|_{Q^2=Q_0^2}.`$ (27) where in our kinematics, $`Q^2=\omega ^2`$. Direct evaluation of these moments using the OPE gives, up to dimension 4 , $`M_n(\xi )=A_n^V(\xi )\left[1+a_n(\xi )\alpha _s+b_n(\xi )\varphi _b^4+c_n(\xi )\varphi _c^4\right],`$ (28) where $`\varphi _b^4`$ $`=`$ $`{\displaystyle \frac{4\pi ^2}{9}}{\displaystyle \frac{\frac{\alpha _s}{\pi }G^2}{(4m_c^2)^2}}`$ $`\varphi _c^4`$ $`=`$ $`{\displaystyle \frac{2\pi ^2}{3}}{\displaystyle \frac{G_2}{(4m_c^2)^2}}m_N\rho _N`$ (29) and $`A_n^V(\xi )`$ $`=`$ $`{\displaystyle \frac{3}{4\pi ^2}}{\displaystyle \frac{2^n(n+1)(n1)!}{(2n+3)!!}}(4m^2)^n(1+\xi )_2^nF_1(n,{\displaystyle \frac{1}{2}},n+{\displaystyle \frac{5}{2}};\rho )`$ $`b_n(\xi )`$ $`=`$ $`{\displaystyle \frac{n(n+1)(n+2)(n+3)}{2n+5}}(1+\xi )^2{\displaystyle \frac{{}_{2}{}^{}F_{1}^{}(n+2,\frac{1}{2},n+\frac{7}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}}`$ $`c_n(\xi )`$ $`=`$ $`b_n(\xi ){\displaystyle \frac{4n(n+1)}{3(2n+5)(1+\xi )^2}}{\displaystyle \frac{{}_{2}{}^{}F_{1}^{}(n+1,\frac{3}{2},n+\frac{7}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}}`$ (30) with $`\xi =\frac{Q_0^2}{4m_c^2}`$ and $`\rho =\frac{\xi }{1+\xi }`$. The factors multiplying the condensate in $`\varphi _b^4`$ and $`\varphi _c^4`$ in eq.(4.1) are defined such that at $`n\mathrm{}`$ $`c_nb_n`$. Moreover, at this limit, one notes that the contribution from dimension 4 is proportional to $`\varphi _b^4+\varphi _c^4`$, which in our kinematical limit originates from the following operator form. $`\varphi _b^4+\varphi _c^4={\displaystyle \frac{16\pi ^2}{9}}{\displaystyle \frac{\frac{\alpha }{\pi }G_{k0}^aG_{k0}^a}{(4m_c^2)^2}},`$ (31) where the operator does not have the trace part of the $`00`$ index. That is, the leading mass shift is proportional to the color electric field squared $`E^2`$ and the effect coming from the magnetic field squared $`B^2`$ disappears. This is consistent with the non-relativistic picture at the infinite quark mass limit, where the Zeeman effect is higher order in $`\alpha _s`$ compared to the Stark effect. The dimension 6 operators contribute to the moment as follows, $`\mathrm{\Delta }M_n^6(\xi )`$ $`=`$ $`A_n^V(\xi )[s_n(\xi )\varphi _s^6+t_n(\xi )\varphi _t^6+x_n(\xi )\varphi _x^6`$ (32) $`+y_n(\xi )\varphi _y^6+z_n(\xi )\varphi _z^6+g_{4n}(\xi )\varphi _{g_4}^6],`$ where $`\varphi _s^6`$ $`=`$ $`{\displaystyle \frac{4\pi ^2}{31080}}{\displaystyle \frac{\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu ;\kappa \kappa }^a}{(4m_c^2)^3}}`$ $`\varphi _t^6`$ $`=`$ $`\left({\displaystyle \frac{2}{3}}\right){\displaystyle \frac{4\pi ^2}{31080}}{\displaystyle \frac{\frac{\alpha _s}{\pi }G_{\alpha \kappa }^aG_{\alpha \lambda ;\lambda \kappa }^a}{(4m_c^2)^3}},`$ (33) come from the scalar operators and $`\varphi _x^6`$ $`=`$ $`\left({\displaystyle \frac{9m_N^2}{2}}\right){\displaystyle \frac{4\pi ^2}{31080}}{\displaystyle \frac{X}{(4m_c^2)^3}}{\displaystyle \frac{\rho }{2m_N}}`$ $`\varphi _y^6`$ $`=`$ $`\left({\displaystyle \frac{3m_N^2}{2}}\right){\displaystyle \frac{4\pi ^2}{31080}}{\displaystyle \frac{Y}{(4m_c^2)^3}}{\displaystyle \frac{\rho }{2m_N}}`$ $`\varphi _z^6`$ $`=`$ $`\left({\displaystyle \frac{3m_N^2}{2}}\right){\displaystyle \frac{4\pi ^2}{31080}}{\displaystyle \frac{Z}{(4m_c^2)^3}}{\displaystyle \frac{\rho }{2m_N}}`$ $`\varphi _{g_4}^6`$ $`=`$ $`(10m_N^4){\displaystyle \frac{4\pi ^2}{31080}}{\displaystyle \frac{G_4}{(4m_c^2)^3}}{\displaystyle \frac{\rho }{2m_N}},`$ (34) are from tensor operators. Note here again that we have put in the prefactors in eq.(33) and eq.(34) such that in the large $`n`$ limit, all the Wilson coefficients become the same. Specifically, the Wilson coefficients are, $`s_n(\xi )`$ $`=`$ $`\sigma _n(\xi )f_n(0,0;120,310,258,72)`$ $`t_n(\xi )`$ $`=`$ $`s_n(\xi )+\sigma _n(\xi )f_n(120,60;300,1320,720)`$ $`x_n(\xi )`$ $`=`$ $`s_n(\xi )+\sigma _n(\xi )f_n(70,25;245,80,6)`$ $`y_n(\xi )`$ $`=`$ $`s_n(\xi )+\sigma _n(\xi )f_n(270,45;1305,1680,666)`$ $`z_n(\xi )`$ $`=`$ $`s_n(\xi )+\sigma _n(\xi )f_n(390,225;1245,880,306)`$ $`g_{4n}(\xi )`$ $`=`$ $`s_n(\xi )+\sigma _n(\xi )f_n(210,33;21,48,18),`$ (35) where $`\sigma _n(\xi )=\left({\displaystyle \frac{8}{1+\xi }}\right){\displaystyle \frac{n(n+2)}{(2n+5)(2n+7)}}`$ (36) and $`f_n(c_1,c_2;a_2,a_3,\mathrm{},a_k)`$ (44) $``$ $`\left(\begin{array}{c}c_1,c_2,a_2,a_3,\mathrm{},a_k\end{array}\right)\left(\begin{array}{cc}(n+3)\frac{{}_{2}{}^{}F_{1}^{}(n+1,\frac{1}{2},n+\frac{9}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}& \\ 2(2n+7)\frac{{}_{2}{}^{}F_{1}^{}(n+1,\frac{1}{2},n+\frac{7}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}& \\ (n+3)(n+4)\frac{{}_{2}{}^{}F_{1}^{}(n+1,\frac{1}{2},n+\frac{9}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}& \\ \frac{(n+3)(n+4)(n+5)}{2!}\frac{{}_{2}{}^{}F_{1}^{}(n+1,\frac{3}{2},n+\frac{9}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}& \\ \mathrm{}& \\ \frac{(n+3)(n+4)\mathrm{}(n+k+2)}{(k1)!}\frac{{}_{2}{}^{}F_{1}^{}(n+1,k+\frac{3}{2},n+\frac{9}{2};\rho )}{{}_{2}{}^{}F_{1}^{}(n,\frac{1}{2},n+\frac{5}{2};\rho )}& \end{array}\right).`$ These functions can be read off from the polarizations which are expressed in terms of linear combinations of the $`J_N(y)`$’s. See appendix G. ### 4.2 estimation of matrix elements Here, we summarize the parameters and matrix elements appearing in our sum rule. As in ref., we will choose the normalization point to be $`Q^2=4m_c^2`$; i.e. $`\xi =1`$. Hence, we will use $`\alpha _s(8m_c^2)=0.21,m_c=1.24\mathrm{GeV}.`$ (45) For the matrix elements, we will use linear density approximation in eq.(23) with $`m_N=0.93`$ GeV and the nuclear matter density to be $`\rho _0=0.17/\mathrm{fm}^3`$. #### 4.2.1 scalar operators There are one scalar operator in dimension 4 and two in dimension 6. 1. $`\frac{\alpha _s}{\pi }G^2_\rho `$ The density dependence of the scalar gluon condensate is obtained from the trace anomaly relation in the chiral limit which to leading order in $`\alpha _s`$ is, $`\theta _\mu ^\mu =\frac{9\alpha _s}{8\pi }G_{\mu \nu }^aG_{\mu \nu }^a`$. Taking the nucleon expectation value, we find, $`{\displaystyle \frac{\alpha _s}{\pi }}G^2_\rho `$ $`=`$ $`{\displaystyle \frac{\alpha _s}{\pi }}G^2_0{\displaystyle \frac{8}{9}}m_N^0\rho `$ (46) $``$ $`(0.35\mathrm{GeV})^4\left(1a_1{\displaystyle \frac{\rho }{\rho _0}}\right).`$ We used $`m_N^0=750`$ MeV for the nucleon mass in the chiral limit to estimate $`a_1={\displaystyle \frac{8m_N^0\rho _0}{9\frac{\alpha _s}{\pi }G^2_0}}0.058.`$ (47) 2. $`\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \alpha ;\alpha \nu }^a_\rho =(\frac{1}{4\pi ^2})g^4j^2_\rho `$ Where we have used the equation of motion in eq.(3) to rewrite the gluon operator in terms of quark operators. This is a four quark operator. Hence to estimate the nuclear matter expectation value, we use ground state saturation hypothesis, where one can factor out the four quark operator in terms of independent two quark operator and take their ground state expectation value. This is a generalization of vacuum dominance hypothesis(VDH) in the vacuum. Hence, $`j^2_\rho `$ $`=`$ $`{\displaystyle \frac{N_c^21}{24^2N_c}}\{\overline{q}q_\rho ^2\mathrm{Tr}[\gamma _\mu \gamma _\mu ]`$ (48) $`+\overline{q}q_\rho \overline{q}\gamma _\kappa q_\rho \mathrm{Tr}[\gamma _\mu \gamma ^\kappa \gamma _\mu +\gamma ^\kappa \gamma _\mu \gamma _\mu ]`$ $`+\overline{q}\gamma _\kappa q_\rho \overline{q}\gamma _\lambda q_\rho \mathrm{Tr}[\gamma ^\kappa \gamma _\mu \gamma ^\lambda \gamma _\mu ]\}`$ $`=`$ $`({\displaystyle \frac{1}{3}})\left[g_{\mu \mu }\overline{q}q_\rho ^2+2\overline{q}\gamma _\mu q_\rho \overline{q}\gamma _\mu q_\rho g_{\mu \mu }\overline{q}\gamma _\lambda q_\rho \overline{q}\gamma _\lambda q_\rho \right]`$ $`=`$ $`({\displaystyle \frac{1}{3}})\left[4\overline{q}q_\rho ^22\overline{q}\gamma _\lambda q_\rho \overline{q}\gamma _\lambda q_\rho \right]`$ $``$ $`{\displaystyle \frac{4}{3}}\left\{\overline{q}q_0^2+{\displaystyle \frac{2\sigma _N\overline{q}q_0}{m_u+m_d}}\rho \right\}`$ $`=`$ $`{\displaystyle \frac{4}{3}}\overline{q}q_0^2\left(1a_2{\displaystyle \frac{\rho }{\rho _0}}\right)(0.24\mathrm{GeV})^6\left(1a_2{\displaystyle \frac{\rho }{\rho _0}}\right),`$ with $`a_2={\displaystyle \frac{2\sigma _N\rho _0}{(m_u+m_d)\overline{q}q_0}}0.69.`$ (49) In getting this, we first note, $`\overline{q}\gamma _\mu q_\rho =\overline{q}\overline{)}uq_\rho u_\mu =q^{}q_\rho u_\mu \rho g_{\mu 0},`$ (50) where we have introduced the four vector $`u_\mu `$ to represent the nuclear matter, which in our case is $`u=(1,0,0,0).`$ We used this in eq.(48) and neglected terms proportional to $`\rho ^2`$ . Other values taken in eq.(48) are, $`\overline{q}q_0=(0.23\mathrm{GeV})^3`$ and $`\frac{2\sigma _N}{m_u+m_d}=0.09/0.014`$. Finally, we have, $`{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \nu }^aG_{\mu \alpha ;\alpha \nu }^a_\rho `$ $`=`$ $`\left({\displaystyle \frac{1}{4\pi ^2}}\right)g^4j^2_\rho `$ (51) $`=`$ $`\left({\displaystyle \frac{1}{4\pi ^2}}\right)g^4j^2_0\left(1a_2{\displaystyle \frac{\rho }{\rho _0}}\right)`$ $`=`$ $`(0.23\mathrm{GeV})^6\left(1a_2{\displaystyle \frac{\rho }{\rho _0}}\right),`$ It should be noted that here we have used a larger value of $`\alpha _s=0.7`$ compared to the perturbative value in eq.(45). 3. $`\frac{\alpha _s}{\pi }G_{\mu \nu }^aG_{\mu \nu ;\alpha \alpha }^a_\rho `$ Using the identities in the appendix, one can rewrite the operator in terms of the three gluon operator and the four quark operator. $`\frac{1}{2\pi ^2}g^3G^3g^4j^2_\rho `$. For the four quark operator, we use the previous result. For the density dependence of the three gluon operator, we assume that the following ratio is constant to linear order in density. $`R={\displaystyle \frac{g^3G^3_0^{2/3}}{g^2G^2_0}}={\displaystyle \frac{g^3G^3_\rho ^{2/3}}{g^2G^2_\rho }},`$ (52) This, together with the previous results on the density dependence of the gluon condensate and the four quark operator, we find, $`{\displaystyle \frac{\alpha _s}{\pi }}G_{\mu \nu }^aG_{\mu \nu ;\alpha \alpha }^a_\rho `$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}g^3G^3g^4j^2_0`$ (53) $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}g^3G^3_0(1{\displaystyle \frac{3}{2}}a_1{\displaystyle \frac{\rho }{\rho _0}}){\displaystyle \frac{1}{2\pi ^2}}g^4j^2_0(1a_2{\displaystyle \frac{\rho }{\rho _0}})`$ $`=`$ $`(0.60\mathrm{GeV})^6\left(1a_3{\displaystyle \frac{\rho }{\rho _0}}\right)`$ where $`a_3={\displaystyle \frac{\frac{3}{2}a_1g^3G^3_0a_2g^4j^2_0}{g^3G^3g^4j^2_0}}0.154`$ (54) #### 4.2.2 twist-2 operators For the twist-2 gluonic operators, we have $`G_{2n}=(i)^{2n2}{\displaystyle \frac{\alpha _s}{\pi }}A_{2n}^G,`$ (55) where $`A_n^G(\mu ^2)=2{\displaystyle _0^1}𝑑xx^{n1}G(x,\mu ^2).`$ (56) We take $`A_2^G(8m_c^2)0.9`$ and $`A_4^G(8m_c^2)0.02`$ . #### 4.2.3 twist-4 operators: dimension 6 Here, there are three independent operators given in eq.(3.5) and contributing to $`\varphi _x^6,\varphi _y^6,\varphi _z^6`$. 1. $`N\|\frac{\alpha _s}{\pi }G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a\|N`$ This operator can be considered to be the second moment of the gluon condensate. For this operator we assume $`N\|{\displaystyle \frac{\alpha _s}{\pi }}G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a\|N`$ $`=`$ $`(p_\mu p_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }m_N^2)\times (){\displaystyle \frac{A_4^G}{A_2^G}}N\|{\displaystyle \frac{\alpha _s}{\pi }}G^2\|N`$ (57) $`=`$ $`(p_\mu p_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }m_N^2)\times {\displaystyle \frac{0.02}{0.9}}\left(+2m_N{\displaystyle \frac{8}{9}}m_N^0\right)`$ $`=`$ $`(p_\mu p_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }m_N^2)\times ({\displaystyle \frac{0.32}{8.1}})m_Nm_N^0.`$ Therefore, $`X=({\displaystyle \frac{0.32}{8.1}})m_Nm_N^0.`$ (58) 2. $`N\|G_{\mu \kappa }G_{\nu \lambda ;\lambda \kappa }\|N=(\frac{1}{4\pi ^2})g^4j_\mu j_\nu _\rho `$ Here, we have used the equation motion again to change this operator to four quark operator. However, this operator is traceless and symmetric. Consider the nuclear matter expectation value of this operator, it becomes, $`𝒮.𝒯.g^4j_\mu ^aj_\nu ^a_\rho =g^4j_\mu ^aj_\nu ^a_\rho \frac{1}{4}g_{\mu \nu }g^4j_\lambda ^aj_\lambda ^a_\rho `$. Assume the ground state saturation hypothesis again, $`j_\mu ^aj_\nu ^a_\rho `$ $`=`$ $`{\displaystyle \frac{N_c^21}{24^2N_c}}\{\overline{q}q_\rho ^2\mathrm{Tr}[\gamma _\mu \gamma _\nu ]`$ $`+\overline{q}q_\rho \overline{q}\gamma _\kappa q_\rho \mathrm{Tr}[\gamma _\mu \gamma ^\kappa \gamma _\nu +\gamma ^\kappa \gamma _\mu \gamma _\nu ]`$ $`+\overline{q}\gamma _\kappa q_\rho \overline{q}\gamma _\lambda q_\rho \mathrm{Tr}[\gamma ^\kappa \gamma _\mu \gamma ^\lambda \gamma _\nu ]\}`$ $`=`$ $`({\displaystyle \frac{1}{3}})\left[g_{\mu \nu }\overline{q}q_\rho ^2+2\overline{q}\gamma _\mu q_\rho \overline{q}\gamma _\nu q_\rho g_{\mu \nu }\overline{q}\gamma _\lambda q_\rho \overline{q}\gamma _\lambda q_\rho \right]`$ $`j_\nu ^aj_\nu ^a_\rho `$ $`=`$ $`({\displaystyle \frac{1}{3}})\left[4\overline{q}q_\rho ^22\overline{q}\gamma _\lambda q_\rho \overline{q}\gamma _\lambda q_\rho \right].`$ (59) So $`𝒮.𝒯.j_\mu ^aj_\nu ^a_\rho `$ $`=`$ $`{\displaystyle \frac{2}{3}}\overline{q}\gamma _\mu q_\rho \overline{q}\gamma _\nu q_\rho +{\displaystyle \frac{1}{6}}g_{\mu \nu }\overline{q}\gamma _\lambda q_\rho \overline{q}\gamma _\lambda q_\rho .`$ (60) Using eq.(50), $`𝒮.𝒯.j_\mu ^aj_\nu ^a_\rho `$ $`=`$ $`{\displaystyle \frac{1}{6}}q^{}q_\rho ^2(u^2g_{\mu \nu }4u_\mu u_\nu )\rho ^2,`$ (61) Hence, $`Y=0`$ (62) 3. $`N\|\frac{\alpha _s}{\pi }G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a\|N`$ Here we use the equation of motion once. Then one finds that the operator has similar structure to a twist-4 quark gluon mixed operator appearing in the twist-4 contribution in deep inelastic scattering. $`N\|G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a\|N`$ $`=`$ $`N\|\overline{q}\gamma _\kappa [D_\nu ,G_{\mu \kappa }]q\|N`$ (63) $``$ $`N\|i\overline{q}\{D_\mu ,^{}F_{\nu \lambda }\}\gamma ^\lambda \gamma _5q\|N`$ $`=`$ $`{\displaystyle \frac{3}{2}}N\|i\overline{u}\gamma _\kappa \gamma _5\{D_\nu ,^{}F_{\mu \kappa }\}u\|N`$ $`=`$ $`(p_\mu p_\nu {\displaystyle \frac{1}{4}}g_{\mu \nu }m_N^2){\displaystyle \frac{3}{2}}K_u^g.`$ Here, $`K_u^g`$ is a parameter introduced in ref. to fit the available deep inelastic scattering data. The fit gives, $`0.3\mathrm{GeV}^2<K_u^g<0.2\mathrm{GeV}^2`$. With this one finds, $`Z(0.097\mathrm{GeV})^2.`$ (64) In Table 1, we summarize values of the re-scaled condensates in Eq. (29),(33) and (34) in the vacuum and at normal nuclear matter density. ### 4.3 numerical analysis The polarization function $`\mathrm{\Pi }`$ satisfies the following energy dispersion relation. $`\mathrm{\Pi }(\omega ^2)={\displaystyle \frac{1}{\pi }}{\displaystyle 𝑑s\frac{\mathrm{Im}\mathrm{\Pi }(s)}{s\omega ^2}}.`$ (65) In the vacuum, the spectral density (Im$`\mathrm{\Pi }(s)`$) consists of the $`J/\psi `$ pole, the contribution from the excited states $`\psi ^{}`$ and the $`D\overline{D}`$ continuum. The contribution from the $`J/\psi `$ and the low-lying resonances can be approximated by delta functions inside the dispersion integral. This is so because the $`J/\psi `$ mass lies below the continuum threshold and is dominated by electromagnetic decays. This is also true in nuclear matter for a $`J/\psi `$ at rest. The inelastic channels opening due to the scattering of the $`J/\psi +N\mathrm{\Lambda }_c(2.28)+\overline{D}(1.87)`$ is forbidden by kinematics. All other processes are OZI rule violating. Hence, the delta function approximation for the $`J/\psi `$ is also good in the nuclear medium. $`\mathrm{Im}\mathrm{\Pi }(s)=f\delta (sm_{J/\psi }^2)+f^{}\delta (sm_\psi ^{}^2)+..`$ (66) The second term just represents a generic excited state contribution. Substituting eq.(66) into the dispersion relation, taking the moments and taking the ratio between the $`(n1)`$th and the $`n`$th moment, one finds, $`{\displaystyle \frac{M_{n1}(Q^2)}{M_n(Q^2)}}=(m_{J/\psi }^2+Q^2)\times \left({\displaystyle \frac{1+\delta _{n1}}{1+\delta _n}}\right)`$ (67) where $`\delta _n={\displaystyle \frac{f^{}}{f}}\left({\displaystyle \frac{m_{J/\psi }^2+Q^2}{m_\psi ^{}^2+Q^2}}\right)^{n+1}.`$ (68) Substituting the vacuum values for the excited states, one finds that the ratio $`\frac{1+\delta _{n1}}{1+\delta _n}`$ goes to 1 for $`n>5`$. Hence, the mass is determined from the relation in eq.(67) at $`Q^2=4m_c^2`$.<sup>1</sup><sup>1</sup>1The determination at different value of $`Q^2`$ was found to be not significant. $`m_{J/\psi }^2={\displaystyle \frac{M_{n1}(4m_c^2)}{M_n(4m_c^2)}}4m_c^2`$ (69) As before, we analyze this equation as a function of $`n`$. In Fig.2(a) , we plot the previous result, which includes the contribution only up to dimension 4. In Fig.2(b), we plot the present result which includes the total dimension 6 contribution. As can be seen from the comparison, the minimum occurs again at similar $`n`$ value and the change from the vacuum results are also similar. We avoid fine tuning of the bare charm quark mass $`m_c`$ to fit the vacuum $`J/\psi `$ mass to its vacuum value, because we are only interested in the shift of the $`J/\psi `$ mass, which is almost independent of this fine tuning. Comparing the two graphs, one notes that the minimum occurs at the same $`n`$ value and the graphs looks similar. Comparing the changes from the vacuum curve and the medium curve at the minimum point, we find $`\mathrm{\Delta }m_{J/\psi }4\mathrm{MeV}.`$ (70) This mass shift is smaller than our previous calculation including dimension 4 only. The main reason for a smaller mass shift compared to including dimension 4 only is as follows. In the vacuum, the dimension 6 operators tend to cancel the dimension 4 operators. This tendency is not only true but more effective in the medium. Therefore, including dimension 6 effects in medium would effectively correspond to a smaller change in the dimension 4 operators in our previous analysis in ref. This implies a smaller mass shift. Below, we will try to elaborate on this point and to explain the interplay of each contribution. * In Fig.3, we plot the contributions from dimension 4 and dimension 6 operators to the moments both in the vacuum and in nuclear medium. As expected, the dimension 6 contributions are smaller compared to dimension 4. Nevertheless, one notes that both in the vacuum and in the nuclear medium, their contributions are opposite to each other so that the contribution from dimension 6 operators tend to cancel the contributions from dimension 4 operators. * The mass shift is coming from the changes in the contributions from each dimension. Hence, for comparison, in Fig.4, we plot the contributions coming from each dimension in the vacuum and in the nuclear medium. One notes that effectively, each contributions becomes smaller in nuclear medium. Numerically, the ratio between the changes in dimension 4 operators and dimension 6 operators are roughly given by (dimension 4 : dimension 6 $`3:1`$). * The changes in dimension 4 operators are dominated by the changes in the scalar gluon condensate. This is evident by comparing $`\varphi _b^4`$ and $`\varphi _c^4`$ in table 1. The comparison between $`\varphi _b^4`$ and $`\varphi _c^4`$ is enough to compare its contribution to the moments, because the Wilson coefficients multiplying them are the same in the large $`n`$ limit. Moreover, $`b_n=c_nn^3`$ and $`b_nc_nn`$ such that $`b_nc_n`$ at moderate $`n`$. That is not quite so for the dimension 6 operators. Here, the Wilson coefficients multiplying $`\varphi _x^6,..`$’s given in eq.(34) are also defined such that they are equal in the large $`n`$ limit and are proportional to $`n^5`$. However, the difference in the Wilson coefficients differ by only one power of $`n`$, that is, the difference between Wilson coefficients are of order $`n^4`$. Therefore, the Wilson coefficients are substantially different at the value of our interest, which is smaller than $`n=10`$, and become similar only at very large value of $`n>100`$. Hence, we will analyze each contribution from dimension 6 more in detail to identify the important contributions. In Fig.5 (a), we have plotted the density dependent part of the scalar, twist-2 and twist-4 contributions. As can be seen, the important contributions are the scalar operators. And among the scalar operators, the contributions coming from $`\varphi _t^6`$ is more important than that from $`\varphi _s^6`$. It should be noted that this is so because of the relatively large Wilson coefficients at small $`n`$ values. The relative importance among the scalar operators are the same also in the vacuum. The next important set of operators are the twist-4 operators. In Fig.5 (b), we plot the contributions from twist-4 operators, which has no vacuum expectation values. The most important contribution is coming from $`G_{\mu \nu }^aG_{\mu \nu ;\alpha \beta }^a`$, which contributes to $`\varphi _x^6`$. Here, the Wilson coefficients are similar at least among the twist-4 operators and the relative importance can be read off from Table 1. * We finally compare the dominant contributions from dimension 4 and dimension 6 operators to the moments. In Fig. 6 (a), we show the total contributions of the scalar dimension 4 and scalar dimension 6 operators to the moments. As can be seen from the figure, the dominant contributions from dimension 6 operators have opposite sign from dimension 4 operator and reduces its contribution. This reduction is enhanced for the density dependent part. This is shown in Fig. 6 (b). The reason why the change in the dimension 6 operator are greater is clear. The gluon condensate eq.(46), which is the dominant dim 4 operator, changes in nuclear matter only by 6 %, whereas, the scalar four quark condensate eq.(48), which is the dominant dim 6 operator, changes by almost 70%. Therefore, although the Wilson coefficients are smaller for dimension 6 operator, the density dependent part of dimension 6 operator are large such that it cancels the the contribution from dimension 4 operator non-trivially. Hence, the reason why the mass shift gets smaller compared to just taking into account dimension 4 operators is because the density dependent changes coming from dimension 4 and dimension 6 operators tend to cancel each other. Therefore, taking into account dimension 6 operators would effectively be equal to a smaller change in the dimension 4 condensate in the previous result, which would have given a smaller mass shift. ## 5 Conclusion In this paper we have calculated the OPE of the correlation function between two vector currents made of heavy quarks up to dimension 6 operators with any tensor structure. The formidable task was to categorize and calculate the corresponding Wilson coefficients of dimension 6 twist-4 gluon operators. This is a first attempt to establish the three independent twist-4 gluon operators in terms of operator basis. Using this result, we have applied our OPE to analyze the mass shift of $`J/\psi `$ in nuclear medium using QCD moment sum rules for the heavy quark system. This is a generalization of our previous result, where we calculated the mass shift using the OPE only up to dimension 4 operators. Unfortunately, the nucleon expectation values of the dimension 6 operators are not as reliable as the dimension 4 operators. Nevertheless, using an order of magnitude estimate for the matrix elements, we find that the mass of $`J/\psi `$ would decrease by about $`4`$ MeV in the nuclear medium. This is $`3`$ MeV smaller than the previous result on including only dimension 4 operators, and shows that the dimension 6 effect is about 40% correction of the dimension 4 effects and goes in the opposite direction. This result seems consistent with the notion that the higher dimensional correction in the vacuum QCD sum rule for the heavy quark system goes like $`\left(G^2/m_c^4\right)^K`$ in the $`r_n=\frac{M_n(\xi =0)}{M_{n1}(\xi =0)}`$ with alternating signs with $`K`$. This also seems to be true in medium and the true mass shift is expected to lie between $`4`$ and $`7`$ MeV. The resulting value of mass shift, is also consistent with the more recent estimates using a totally different approach. This is a result for a $`J/\psi `$ at rest with respect to the nuclear medium. However, since we have calculated the OPE for a general external four momentum, our results can be easily and reliably generalized to study the moving $`J/\psi `$ and also to finite temperature, which would be also interesting in relation to the ongoing discussion of $`J/\psi `$ suppression in RHIC due to a comover model. ## 6 Acknowledgement This work was supported by KOSEF grant number 1999-2-111-005-5, by the BK 21 project of the Korean Ministry of Education and by the Yonsei university research grant. We would like to thank T. Hatsuda, A. Hayashigaki, and W. Weise for useful discussions. We particularly thank D. Kharzeev and P. Morath for pointing out the difference in reference and . ## Appendix A Identities for spin-2 dimension 6 gluon operators The spin-2 dimension 6 gluon operators in eq.(2) are not independent. The following relations holds among them. $`G_{\kappa \lambda ;\mu }^aG_{\nu \kappa ;\lambda }^a`$ $`=`$ $`G_{\kappa \lambda ;\mu }^a{\displaystyle \frac{1}{2}}(G_{\lambda \nu ;\kappa }^a+G_{\nu \kappa ;\lambda }^a)={\displaystyle \frac{1}{2}}G_{\kappa \lambda ;\mu }^aG_{\kappa \lambda ;\nu }^a`$ (71) $`G_{\mu \kappa ;\lambda }^aG_{\kappa \lambda ;\nu }^a`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(G_{\lambda \mu ;\kappa }^a+G_{\mu \kappa ;\lambda }^a\right)G_{\kappa \lambda ;\nu }^a={\displaystyle \frac{1}{2}}G_{\kappa \lambda ;\mu }^aG_{\kappa \lambda ;\nu }^a`$ (72) $`G_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a`$ $`=`$ $`G_{\mu \kappa ;\nu \lambda }^aG_{\kappa \lambda }^a=(G_{\mu \kappa ;\lambda \nu }^a+gf^{abc}G_{\mu \kappa }^bG_{\nu \lambda }^c)G_{\kappa \lambda }^a`$ $`=`$ $`()gfG_{\mu \nu }^3+G_{\mu \kappa ;\lambda }^aG_{\kappa \lambda ;\nu }^a`$ $`\mathrm{Or},`$ $`gfG_{\mu \nu }^3`$ $`=`$ $`G_{\mu \kappa ;\lambda }^aG_{\kappa \lambda ;\nu }^aG_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a`$ (73) $`=`$ $`\left({\displaystyle \frac{1}{2}}G_{\kappa \lambda ;\mu }^aG_{\kappa \lambda ;\nu }^a+G_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a\right)`$ $`G_{\kappa \lambda ;\lambda }^aG_{\mu \kappa ;\nu }^a`$ $`=`$ $`G_{\kappa \lambda ;\lambda \nu }^aG_{\mu \kappa }^a=\left(G_{\kappa \lambda ;\nu \lambda }^a+gf^{abc}G_{\kappa \lambda }^bG_{\lambda \nu }^c\right)G_{\mu \kappa }^a`$ (74) $`=`$ $`+G_{\kappa \lambda ;\nu }^aG_{\mu \kappa ;\lambda }^ag𝒢_{\mu \nu }^3=G_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a`$ $`G_{\mu \kappa ;\lambda }^aG_{\nu \lambda ;\kappa }^a`$ $`=`$ $`G_{\mu \kappa ;\lambda \kappa }^aG_{\nu \lambda }^a=\left(G_{\mu \kappa ;\kappa \lambda }^a+gf^{abc}G_{\mu \kappa }^bG_{\lambda \kappa }^c\right)G_{\nu \lambda }^a`$ (75) $`=`$ $`G_{\mu \kappa ;\kappa }^aG_{\nu \lambda ;\lambda }^ag𝒢_{\mu \nu }^3`$ $`=`$ $`G_{\mu \kappa ;\kappa }^aG_{\nu \lambda ;\lambda }^a+{\displaystyle \frac{1}{2}}G_{\kappa \lambda ;\mu }^aG_{\kappa \lambda ;\nu }^a+G_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a`$ $`G_{\mu \kappa ;\lambda }^aG_{\nu \kappa ;\lambda }^a`$ $`=`$ $`G_{\mu \kappa ;\lambda \lambda }^aG_{\nu \kappa }^a=g\left(2f^{abc}G_{\mu \lambda }^bG_{\kappa \lambda }^c+j_{\mu ;\kappa }^aj_{\kappa ;\mu }^a\right)G_{\nu \kappa }^a`$ (76) $`=`$ $`(2)g𝒢_{\mu \nu }^3+g^2j_\mu ^aj_\nu ^agj_\kappa ^aG_{\nu \kappa ;\mu }^a`$ $`=`$ $`G_{\kappa \lambda ;\mu }^aG_{\kappa \lambda ;\nu }^a+2G_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a+g^2j_\mu ^aj_\nu ^aG_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a`$ $`=`$ $`G_{\mu \kappa ;\nu }^aG_{\kappa \lambda ;\lambda }^a+g^2j_\mu ^aj_\nu ^a+G_{\kappa \lambda ;\mu }^aG_{\kappa \lambda ;\nu }^a`$ Using these, one can reduce the operators in eq.(2) to the three independent operators in eq.(5). ## Appendix B Current Conservation in the Polarization Operators Let us consider the polarization of vector currents defined in eq.(6). When the current is Hermitian, i.e. $`j_\mu =j_\mu ^{}`$, the contribution from spin-2 dimension 4 operator can in general be written in terms of the following four independent terms. $`\mathrm{\Pi }_{\mu \nu }\left(q\right)={\displaystyle \frac{1}{Q^2}}\left[\theta _{\mu \nu }^1+{\displaystyle \frac{q^\alpha q^\beta q_\mu q_\nu }{Q^4}}\theta _{\alpha \beta }^2+{\displaystyle \frac{q^\alpha q_\nu }{Q^2}}\theta _{\mu \alpha }^3+{\displaystyle \frac{q^\alpha q_\mu }{Q^2}}\theta _{\nu \alpha }^3+g_{\mu \nu }{\displaystyle \frac{q^\alpha q^\beta }{Q^2}}\theta _{\alpha \beta }^4\right]`$ (77) Imposing current conservation, on finds, $`\theta _{\mu \nu }^1`$ $`=`$ $`\theta _{\mu \nu }^3`$ $`\theta _{\mu \nu }^2`$ $`=`$ $`\theta _{\mu \nu }^1+\theta _{\mu \nu }^4`$ (78) Similar relations holds for higher dimensional operator. This is the operator form of showing the existence of two independent polarization directions in medium for eq.(6). Summarizing similar relations for different spins and dimensions, we have 1. dimension 4 and spin 2 $`\mathrm{\Pi }_{\mu \nu }^{4,2}\left(q\right)={\displaystyle \frac{1}{Q^2}}\left[_{\mu \nu }^2+{\displaystyle \frac{1}{Q^2}}\left(q_\rho q_\mu _{\rho \nu }^2+q_\rho q_\nu _{\rho \mu }^2\right)+g_{\mu \nu }{\displaystyle \frac{q_\rho q_\sigma }{Q^2}}𝒥_{\rho \sigma }^2+{\displaystyle \frac{q_\mu q_\nu q_\rho q_\sigma }{Q^4}}\left(_{\rho \sigma }^2+𝒥_{\rho \sigma }^2\right)\right]`$ 2. dimension 6 and spin 2 $`\mathrm{\Pi }_{\mu \nu }^{6,2}\left(q\right)={\displaystyle \frac{1}{\left(Q^2\right)^2}}\left[I_{\mu \nu }^2+{\displaystyle \frac{1}{Q^2}}\left(q_\rho q_\mu I_{\rho \nu }^2+q_\rho q_\nu I_{\rho \mu }^2\right)+g_{\mu \nu }{\displaystyle \frac{q_\rho q_\sigma }{Q^2}}J_{\rho \sigma }^2+{\displaystyle \frac{q_\mu q_\nu q_\rho q_\sigma }{Q^4}}\left(I_{\rho \sigma }^2+J_{\rho \sigma }^2\right)\right]`$ 3. dimension 6 and spin 4 $`\mathrm{\Pi }_{\mu \nu }^{6,4}\left(q\right)`$ $`=`$ $`{\displaystyle \frac{q_\kappa q_\lambda }{\left(Q^2\right)^3}}[I_{\kappa \lambda \mu \nu }^4+{\displaystyle \frac{1}{Q^2}}(q_\rho q_\mu I_{\kappa \lambda \rho \nu }^4+q_\rho q_\nu I_{\kappa \lambda \rho \mu }^4)+g_{\mu \nu }{\displaystyle \frac{q_\rho q_\sigma }{Q^2}}J_{\kappa \lambda \rho \sigma }^4`$ $`+{\displaystyle \frac{q_\mu q_\nu q_\rho q_\sigma }{Q^4}}(I_{\kappa \lambda \rho \sigma }^4+J_{\kappa \lambda \rho \sigma }^4)]`$ ## Appendix C Propagators Here we summarize the quark propagator in the presence of the external gauge field in the fixed point gauge. $`iS\left(p\right)`$ $``$ $`{\displaystyle d^4xe^{ipx}iS(x,0)}`$ (79) $`=`$ $`iS^{\left(0\right)}\left(p\right)+{\displaystyle d^4xe^{ipx}gd^4ziS^{\left(0\right)}\left(xz\right)i\overline{)}A\left(z\right)iS^{\left(0\right)}\left(z\right)}`$ $`+{\displaystyle d^4xe^{ipx}g^2d^4z^{}d^4ziS^{\left(0\right)}\left(xz^{}\right)i\overline{)}A\left(z^{}\right)iS^{\left(0\right)}\left(z^{}z\right)i\overline{)}A\left(z\right)iS^{\left(0\right)}\left(z\right)}`$ $`+\mathrm{}`$ and $`i\stackrel{~}{S}(p)`$ $``$ $`{\displaystyle d^4xe^{ipx}iS(0,x)}.`$ (80) In the fixed point gauge, we write the field in terms of covariant operators, $`A_\mu (x)`$ $`=`$ $`{\displaystyle \frac{1}{20!}}x_\rho G_{\rho \mu }(0)+{\displaystyle \frac{1}{31!}}x_\alpha x_\rho \left(D_\alpha G_{\rho \mu }(0)\right)`$ (81) $`+{\displaystyle \frac{1}{42!}}x_\alpha x_\beta x_\rho \left(D_\alpha D_\beta G_{\rho \mu }(0)\right)+\mathrm{}`$ Collecting terms, we can write the full propagator in terms of gauge covariant fields. A few symbols are used for convenience: $`\widehat{x}_\alpha i{\displaystyle \frac{\stackrel{}{}}{p}},\underset{\alpha }{\overset{}{x}}i{\displaystyle \frac{\stackrel{}{}}{p}},`$ (82) and $`\{\alpha ,\beta \}={\displaystyle \frac{1}{\overline{)}pm}}\gamma _\alpha {\displaystyle \frac{1}{\overline{)}pm}}\gamma _\beta {\displaystyle \frac{1}{\overline{)}pm}}`$ $`\{\alpha ,\beta ,\sigma \}={\displaystyle \frac{1}{\overline{)}pm}}\gamma _\alpha {\displaystyle \frac{1}{\overline{)}pm}}\gamma _\beta {\displaystyle \frac{1}{\overline{)}pm}}\gamma _\sigma {\displaystyle \frac{1}{\overline{)}pm}}`$ $`\mathrm{},`$ (83) and $`P(\alpha ,\beta ,\gamma \mathrm{})`$ means sum of all possible permutations in $`\alpha ,\beta ,\gamma \mathrm{}`$. Here we simply list the propagators. 1. $`iS^{(0)}(p)=i\stackrel{~}{S}^{(0)}(p)`$ $`={\displaystyle \frac{i}{\overline{)}pm}}`$ 2. $`iS_G(p)=i\stackrel{~}{S}_G(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\widehat{x}_\alpha G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)={\displaystyle \frac{1}{2}}G_{\alpha \beta }\{\beta ,\alpha \}`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\underset{\alpha }{\overset{}{x}}G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)={\displaystyle \frac{1}{2}}G_{\alpha \beta }\{\alpha ,\beta \}`$ 3. $`iS_{G^2}(p)=i\stackrel{~}{S}_{G^2}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\widehat{x}_\alpha G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{2}}\widehat{x}_\rho G_{\rho \sigma }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{4}}G_{\alpha \beta }G_{\rho \sigma }\left[\{\beta ,\alpha ,\sigma ,\rho \}+\{\beta ,\sigma ,P(\alpha ,\rho )\}\right]`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\underset{\alpha }{\overset{}{x}}G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{2}}\underset{\rho }{\overset{}{x}}G_{\rho \sigma }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{4}}G_{\alpha \beta }G_{\rho \sigma }\left[\{\alpha ,\beta ,\rho ,\sigma \}+\{P(\alpha ,\rho ),\beta ,\sigma \}\right]`$ 4. As for $`G^3`$ part, * $`iS_{G^3}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\widehat{x}_\alpha G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{2}}\widehat{x}_\rho G_{\rho \sigma }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\lambda {\displaystyle \frac{1}{2}}\widehat{x}_\kappa G_{\kappa \lambda }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{1}{8}}G_{\alpha \beta }G_{\rho \sigma }G_{\kappa \lambda }[\{\beta ,\alpha ,\sigma ,\rho ,\lambda ,\kappa \}+\{\beta ,\alpha ,\sigma ,\lambda ,P(\rho ,\kappa )\}`$ $`+\{\beta ,\sigma ,P(\alpha ,\rho ),\lambda ,\kappa \}+\{\beta ,\sigma ,\rho ,\lambda ,P(\alpha ,\kappa )\}`$ $`+\{\beta ,\sigma ,\alpha ,\lambda ,P(\rho ,\kappa )\}+\{\beta ,\sigma ,\lambda ,P(\alpha ,\rho ,\kappa )\}]`$ * $`i\stackrel{~}{S}_{G^3}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\stackrel{}{x}_\alpha G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{2}}\stackrel{}{x}_\rho G_{\rho \sigma }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\lambda {\displaystyle \frac{1}{2}}\stackrel{}{x}_\kappa G_{\kappa \lambda }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{1}{8}}G_{\alpha \beta }G_{\rho \sigma }G_{\kappa \lambda }[\{\alpha ,\beta ,\rho ,\sigma ,\kappa ,\lambda \}+\{P(\alpha ,\rho ),\beta ,\sigma ,\kappa ,\lambda \}`$ $`+\{\alpha ,\beta ,P(\kappa ,\rho ),\sigma ,\lambda \}+\{P(\alpha ,\kappa ,\rho ),\beta ,\sigma ,\lambda \}`$ $`+\{P(\alpha ,\rho ),\beta ,\kappa ,\sigma ,\lambda \}+\{P(\alpha ,\kappa ),\beta ,\rho ,\sigma ,\lambda \}]`$ 5. As for $`DG`$ part, * $`iS_{DG}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{3}}\widehat{x}_\alpha \widehat{x}_\beta D_\alpha G_{\beta \sigma }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{3}}D_\alpha G_{\beta \sigma }\{\sigma ,P(\alpha ,\beta )\}`$ * $`i\stackrel{~}{S}_{DG}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i(\gamma _\sigma {\displaystyle \frac{1}{3}}\stackrel{}{x}_\alpha \stackrel{}{x}_\beta D_\alpha G_{\beta \sigma })iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{3}}D_\alpha G_{\beta \sigma }\{P(\alpha ,\beta ),\sigma \}`$ 6. As for $`D^2G`$ part, * $`iS_{D^2G}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{8}}\widehat{x}_\alpha \widehat{x}_\beta \widehat{x}_\nu D_\alpha D_\beta G_{\nu \sigma }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{1}{8}}D_\alpha D_\beta G_{\nu \sigma }\{\sigma ,P(\alpha ,\beta ,\nu )\}`$ * $`i\stackrel{~}{S}_{D^2G}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i(\gamma _\sigma {\displaystyle \frac{1}{8}}\stackrel{}{x}_\alpha \stackrel{}{x}_\beta \stackrel{}{x}_\nu D_\alpha D_\beta G_{\nu \sigma })iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`+{\displaystyle \frac{1}{8}}D_\alpha D_\beta G_{\nu \sigma }\{P(\alpha ,\beta ,\nu ),\sigma \}`$ 7. As for $`(DG)^2`$ part, * $`iS_{(DG)^2}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\kappa {\displaystyle \frac{1}{3}}\widehat{x}_\alpha \widehat{x}_\beta D_\alpha G_{\beta \kappa }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\lambda {\displaystyle \frac{1}{3}}\widehat{x}_\rho \widehat{x}_\sigma D_\rho G_{\sigma \lambda }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{9}}D_\alpha G_{\beta \kappa }D_\rho G_{\sigma \lambda }[\{\kappa ,P(\alpha ,\beta ),\lambda ,P(\rho ,\sigma )\}`$ $`+\{\kappa ,\alpha ,\lambda ,P(\beta ,\rho ,\sigma )\}+\{\kappa ,\beta ,\lambda ,P(\alpha ,\rho ,\sigma )\}+\{\kappa ,\lambda ,P(\alpha ,\beta ,\rho ,\sigma )\}]`$ * $`i\stackrel{~}{S}_{(DG)^2}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i(\gamma _\kappa {\displaystyle \frac{1}{3}}\stackrel{}{x}_\alpha \stackrel{}{x}_\beta D_\alpha G_{\beta \kappa })iS^{\left(0\right)}\left(p\right)i(\gamma _\lambda {\displaystyle \frac{1}{3}}\stackrel{}{x}_\rho \stackrel{}{x}_\sigma D_\rho G_{\sigma \lambda })iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{9}}D_\alpha G_{\beta \kappa }D_\rho G_{\sigma \lambda }[\{P(\alpha ,\beta ),\kappa ,P(\rho ,\sigma ),\lambda \}`$ $`+\{P(\alpha ,\beta ,\rho ),\kappa ,\sigma ,\lambda \}+\{P(\alpha ,\beta ,\sigma ),\kappa ,\rho ,\lambda \}+\{P(\alpha ,\beta ,\rho ,\sigma ),\kappa ,\lambda \}]`$ 8. As for $`GD^2G`$ part, * $`iS_{GD^2G}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\widehat{x}_\alpha G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{8}}\widehat{x}_\rho \widehat{x}_\kappa \widehat{x}_\lambda D_\rho D_\kappa G_{\lambda \sigma }\right)iS^{\left(0\right)}`$ $`=`$ $`{\displaystyle \frac{i}{16}}G_{\alpha \beta }D_\rho D_\kappa G_{\lambda \sigma }\left[\{\beta ,\alpha ,\sigma ,P(\rho ,\kappa ,\lambda )\}+\{\beta ,\sigma ,P(\alpha ,\rho ,\kappa ,\lambda )\}\right]`$ * $`i\stackrel{~}{S}_{GD^2G}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i(\gamma _\beta {\displaystyle \frac{1}{2}}\stackrel{}{x}_\alpha G_{\alpha \beta })iS^{\left(0\right)}\left(p\right)i(\gamma _\sigma {\displaystyle \frac{1}{8}}\stackrel{}{x}_\rho \stackrel{}{x}_\kappa \stackrel{}{x}_\lambda D_\rho D_\kappa G_{\lambda \sigma })iS^{\left(0\right)}`$ $`=`$ $`{\displaystyle \frac{i}{16}}G_{\alpha \beta }D_\rho D_\kappa G_{\lambda \sigma }[\{\alpha ,\beta ,P(\rho ,\kappa ,\lambda ),\sigma \}+\{P(\alpha ,\rho ),\beta ,P(\kappa ,\lambda ),\sigma \}`$ $`+\{P(\alpha ,\kappa ),\beta ,P(\rho ,\lambda ),\sigma \}+\{P(\alpha ,\lambda ),\beta ,P(\kappa ,\rho ),\sigma \}`$ $`+\{P(\alpha ,\kappa ,\lambda ),\beta ,\rho ,\sigma \}+\{P(\alpha ,\rho ,\lambda ),\beta ,\kappa ,\sigma \}`$ $`+\{P(\alpha ,\kappa ,\rho ),\beta ,\lambda ,\sigma \}+\{P(\alpha ,\rho ,\kappa ,\lambda ),\beta ,\sigma \}]`$ 9. As for $`D^2GG`$ part, * $`iS_{D^2GG}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i\left(\gamma _\sigma {\displaystyle \frac{1}{8}}\widehat{x}_\rho \widehat{x}_\kappa \widehat{x}_\lambda D_\rho D_\kappa G_{\lambda \sigma }\right)iS^{\left(0\right)}\left(p\right)i\left(\gamma _\beta {\displaystyle \frac{1}{2}}\widehat{x}_\alpha G_{\alpha \beta }\right)iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{16}}D_\rho D_\kappa G_{\lambda \sigma }G_{\alpha \beta }[\{\sigma ,P(\rho ,\kappa ,\lambda ),\beta ,\alpha \}+\{\sigma ,P(\kappa ,\lambda ),\beta ,P(\alpha ,\rho )\}`$ $`+\{\sigma ,\rho ,\beta ,P(\kappa ,\lambda ,\alpha )\}+\{\sigma ,\beta ,P(\rho ,\kappa ,\lambda ,\alpha )\}`$ $`+\{\sigma ,P(\kappa ,\rho ),\beta ,P(\alpha ,\lambda )\}+\{\sigma ,\kappa ,\beta ,P(\alpha ,\lambda ,\rho )\}`$ $`+\{\sigma ,P(\lambda ,\rho ),\beta ,P(\alpha ,\kappa )\}+\{\sigma ,\lambda ,\beta ,P(\rho ,\alpha ,\kappa )\}]`$ * $`i\stackrel{~}{S}_{D^2GG}(p)`$ $`=`$ $`iS^{\left(0\right)}\left(p\right)i(\gamma _\sigma {\displaystyle \frac{1}{8}}\stackrel{}{x}_\rho \stackrel{}{x}_\kappa \stackrel{}{x}_\lambda D_\rho D_\kappa G_{\lambda \sigma })iS^{\left(0\right)}\left(p\right)i(\gamma _\beta {\displaystyle \frac{1}{2}}\stackrel{}{x}_\alpha G_{\alpha \beta })iS^{\left(0\right)}\left(p\right)`$ $`=`$ $`{\displaystyle \frac{i}{16}}D_\rho D_\kappa G_{\lambda \sigma }G_{\alpha \beta }\left[\{P(\rho ,\kappa ,\lambda ),\sigma ,\alpha ,\beta \}+\{P(\rho ,\kappa ,\lambda ,\alpha ),\sigma ,\beta \}\right]`$ ## Appendix D spin structure of operators of dimension 6 In computing the Wilson coefficients we used the following reduction of the Lorentz indices to the spin-2 operators. $`f^{abc}G_{\mu \nu }^aG_{\alpha \beta }^bG_{\rho \sigma }^c`$ $`=`$ $`A_{\mu \alpha }c_{\nu \rho \beta \sigma }A_{\mu \beta }c_{\nu \rho \alpha \sigma }A_{\nu \alpha }c_{\mu \rho \beta \sigma }+A_{\nu \beta }c_{\mu \rho \alpha \sigma }`$ (84) $`A_{\mu \rho }c_{\nu \alpha \sigma \beta }+A_{\mu \sigma }c_{\nu \alpha \rho \beta }+A_{\nu \rho }c_{\mu \alpha \sigma \beta }A_{\nu \sigma }c_{\mu \alpha \rho \beta }`$ $`+A_{\alpha \rho }c_{\beta \mu \sigma \nu }A_{\alpha \sigma }c_{\beta \mu \rho \nu }A_{\beta \rho }c_{\alpha \mu \sigma \nu }+A_{\beta \sigma }c_{\alpha \mu \rho \nu }`$ $`G_{\mu _1\nu _1}^aG_{\mu _2\nu _2;\alpha \beta }^a`$ $`=`$ $`K_{\alpha \beta }c_{\mu _1\mu _2\nu _1\nu _2}`$ (85) $`+P_{\beta \mu _1}c_{\alpha \mu _2\nu _1\nu _2}P_{\beta \nu _1}c_{\alpha \mu _2\mu _1\nu _2}`$ $`+J_{\alpha \mu _2}c_{\beta \mu _1\nu _2\nu _1}J_{\alpha \nu _2}c_{\beta \mu _1\mu _2\nu _1}`$ $`+Q_{\beta \mu _2}c_{\alpha \mu _1\nu _2\nu _1}Q_{\beta \nu _2}c_{\alpha \mu _1\mu _2\nu _1}`$ $`+W_{\alpha \mu _1}c_{\beta \mu _2\nu _1\nu _2}W_{\alpha \nu _1}c_{\beta \mu _2\mu _1\nu _2}`$ $`+L_{\mu _1\mu _2}d_{\beta \nu _1\alpha \nu _2}L_{\nu _1\mu _2}d_{\beta \mu _1\alpha \nu _2}`$ $`L_{\mu _1\nu _2}d_{\beta \nu _1\alpha \mu _2}+L_{\nu _1\nu _2}d_{\beta \mu _1\alpha \mu _2}`$ $`+M_{\mu _1\mu _2}d_{\beta \alpha \nu _1\nu _2}M_{\nu _1\mu _2}d_{\alpha \beta \mu _1\nu _2}`$ $`M_{\mu _1\nu _2}d_{\alpha \beta \nu _1\mu _2}+M_{\nu _1\nu _2}d_{\alpha \beta \mu _1\mu _2}`$ $`+T_{\mu _1\mu _2}d_{\beta \nu _2\alpha \nu _1}T_{\nu _1\mu _2}d_{\beta \nu _2\alpha \mu _1}`$ $`T_{\mu _1\nu _2}d_{\beta \mu _2\alpha \nu _1}+T_{\nu _1\nu _2}d_{\beta \mu _2\alpha \mu _1},`$ where $`c_{\alpha \beta \mu \nu }`$ $``$ $`g_{\alpha \beta }g_{\mu \nu }g_{\alpha \nu }g_{\mu \beta },`$ $`d_{\alpha \beta \mu \nu }`$ $``$ $`g_{\alpha \beta }g_{\mu \nu }`$ (86) and, if we take 3 operators, $`ddgg1G_{\kappa \lambda }^aG_{\kappa \lambda ;\mu \nu }^a,ddgg2g^2j_\mu ^aj_\nu ^a,`$ and $`ddgg3G_{\mu \kappa }^aG_{\kappa \lambda ;\lambda \nu }^a`$ as our basis, we get $`A=(ddgg1+2ddgg3)/4,`$ $`K=2P=2J=(13ddgg1+ddgg2+19ddgg3)/80,`$ $`Q=W=(27ddgg1+ddgg261ddgg3)/160,`$ $`L=(4ddgg1+3ddgg2+7ddgg3)/40,`$ $`M=(29ddgg1+13ddgg2+47ddgg3)/160,`$ $`T=(6ddgg1+3ddgg213ddgg3)/40.`$ (87) As for the spin 4 part, we have a simpler reduction: $`G_{\mu _1\nu _1}^aG_{\mu _2\nu _2;\alpha \beta }^a`$ $`=`$ $`W_{\alpha \beta \mu _1\mu _2}g_{\nu _1\nu _2}+W_{\alpha \beta \nu _1\nu _2}g_{\mu _1\mu _2}`$ (88) $`W_{\alpha \beta \mu _1\nu _2}g_{\nu _1\mu _2}W_{\alpha \beta \nu _1\mu _2}g_{\mu _1\nu _2},`$ where $`W_{\alpha \beta \mu _1\mu _2}={\displaystyle \frac{1}{2}}G_{\mu _1\kappa }^aG_{\mu _2\kappa ;\alpha \beta }^a.`$ (89) ## Appendix E Integrations with respect to Feynman Parameter In general, after Feynman integration, one can write the polarizations in terms of linear sums of $`J`$’s defined in eq.(13). $`\mathrm{\Pi }(q)=Q^a\left[a_0J_0+a_1J_1+a_2J_2+\mathrm{}+a_kJ_k\right].`$ (90) To reduce the polarization function into this final form, we use the following steps and identities. After the Feynman integral, the polarization function will be a sum of $`Imn`$’s: $`I_N^{mn}(Q^2,m^2){\displaystyle _0^1}𝑑x{\displaystyle \frac{x^n(1x)^m}{\left[m^2+Q^2x(1x)\right]^N}}=I_N^{nm}(Q^2,m^2)`$ (91) We then follow the following steps. 1. $`I^{mn}`$ can be expressed in terms of $`I_N^nI_N^{nn}`$ using the following identity. $`x^n+\left(1x\right)^n`$ $`=`$ $`1\{{}_{n}{}^{}C_{1}^{}x^{n1}(1x)+_nC_2x^{n2}(1x)^2+\mathrm{}+_nC_kx^{nk}(1x)^k`$ $`+\mathrm{}+_nC_{n1}x(1x)^{n1}\}`$ $`=`$ $`1_nC_1\left\{x^{n1}\left(1x\right)+x\left(1x\right)^{n1}\right\}`$ $`_nC_2\left\{x^{n2}\left(1x\right)^2+x^2\left(1x\right)^{n2}\right\}\mathrm{}`$ $`=`$ $`1_nC_1x\left(1x\right)\left\{x^{n2}+\left(1x\right)^{n2}\right\}`$ $`_nC_2x^2\left(1x\right)^2\left\{x^{n4}+\left(1x\right)^{n4}\right\}\mathrm{}`$ 2. Then we can reduce $`I_N^n`$ to $`I_N^0I_N`$ using $`I_N^n=\frac{1}{Q^2}(I_{N1}^{n1}m^2I_N^{n1})=\frac{1}{Q^2}I_{N1}^{n1}\frac{1}{y}I_{N1}^{n1}`$ 3. We then introduce the dimensionless function $`J_N(y)=\left(\frac{Q^2}{y}\right)^NI_N`$, where $`y=Q^2/m^2`$. 4. Finally, we use the recurrence relation. $`yJ_N(y)={\displaystyle \frac{2}{N1}}+{\displaystyle \frac{4N6}{N1}}J_{N1}(y)4J_N(y)`$ (92) 5. We can also integrate explicitly, $`J_N\left(y\right)={\displaystyle \frac{\left(2N3\right)!!}{\left(N1\right)!}}\left[\left({\displaystyle \frac{s1}{2s}}\right)^N\sqrt{s}\mathrm{log}{\displaystyle \frac{\sqrt{s}+1}{\sqrt{s}1}}+{\displaystyle \underset{k=1}{\overset{N1}{}}}{\displaystyle \frac{\left(k1\right)!}{\left(2k1\right)!!}}\left({\displaystyle \frac{s1}{2s}}\right)^{Nk}\right],`$ (93) where $`s=1+4/y`$. ## Appendix F An example: evaluation of a diagram Diagram 1a: We show the evaluation of the Feynman diagrams here. The first diagram in Fig. 1 is taken as an example. $`i{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\mathrm{Tr}\left[\gamma _\mu S^{\left(0\right)}\left(k+q\right)\gamma _\nu \stackrel{~}{S}_{G^3}\left(k\right)\right]}`$ (94) $`=`$ $`i{\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}\mathrm{Tr}[\gamma _\mu {\displaystyle \frac{1}{\overline{)}k+\overline{)}qm}}\gamma _\nu \left({\displaystyle \frac{i}{8}}\right)G_{\alpha \beta }^aG_{\rho \sigma }^bG_{\kappa \lambda }^ct^at^bt^c[\{\alpha ,\beta ,\rho ,\sigma ,\kappa ,\lambda \}+\{P(\alpha ,\rho ),\beta ,\sigma ,\kappa ,\lambda \}`$ $`+\{\alpha ,\beta ,P(\kappa ,\rho ),\sigma ,\lambda \}+\{P(\alpha ,\kappa ,\rho ),\beta ,\sigma ,\lambda \}+\{P(\alpha ,\rho ),\beta ,\kappa ,\sigma ,\lambda \}+\{P(\alpha ,\kappa ),\beta ,\rho ,\sigma ,\lambda \}]]`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{i}{4}}\left(f^{abc}G_{\alpha \beta }^aG_{\rho \sigma }^bG_{\kappa \lambda }^c\right){\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}{\displaystyle \frac{1}{\left(k+q\right)^2m^2}}\mathrm{Tr}[\gamma _\mu (\overline{)}k+\overline{)}q+m)\gamma _\nu `$ $`\times [\{\alpha ,\beta ,\rho ,\sigma ,\kappa ,\lambda \}+\{P(\alpha ,\rho ),\beta ,\sigma ,\kappa ,\lambda \}+\{\alpha ,\beta ,P(\kappa ,\rho ),\sigma ,\lambda \}`$ $`+\{P(\alpha ,\kappa ,\rho ),\beta ,\sigma ,\lambda \}+\{P(\alpha ,\rho ),\beta ,\kappa ,\sigma ,\lambda \}+\{P(\alpha ,\kappa ),\beta ,\rho ,\sigma ,\lambda \}]]`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{i}{4}}(A_{\mu \alpha }c_{\nu \rho \beta \sigma }A_{\mu \beta }c_{\nu \rho \alpha \sigma }A_{\nu \alpha }c_{\mu \rho \beta \sigma }+A_{\nu \beta }c_{\mu \rho \alpha \sigma }A_{\mu \rho }c_{\nu \alpha \sigma \beta }+A_{\mu \sigma }c_{\nu \alpha \rho \beta }`$ $`+A_{\nu \rho }c_{\mu \alpha \sigma \beta }A_{\nu \sigma }c_{\mu \alpha \rho \beta }+A_{\alpha \rho }c_{\beta \mu \sigma \nu }A_{\alpha \sigma }c_{\beta \mu \rho \nu }A_{\beta \rho }c_{\alpha \mu \sigma \nu }+A_{\beta \sigma }c_{\alpha \mu \rho \nu })`$ $`\times {\displaystyle }{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}}{\displaystyle \frac{1}{\left(k+q\right)^2m^2}}{\displaystyle \frac{1}{\left(k^2m^2\right)^7}}\mathrm{Tr}[\gamma _\mu (\overline{)}k+\overline{)}q+m)\gamma _\nu `$ $`\times \{(\overline{)}k+m)\gamma _\alpha (\overline{)}k+m)\gamma _\beta (\overline{)}k+m)\gamma _\rho (\overline{)}k+m)\gamma _\sigma (\overline{)}k+m)\gamma _\kappa (\overline{)}k+m)\gamma _\lambda (\overline{)}k+m)+\mathrm{}\}]`$ $``$ $`i{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\frac{1}{\left(k+q\right)^2m^2}\frac{1}{\left(k^2m^2\right)^7}f_{1a}\left(k\right)}`$ $`=`$ $`i{\displaystyle \frac{d^4k}{\left(2\pi \right)^4}\left[_0^1𝑑x\frac{x^6}{\left[\left\{k+\left(1x\right)q\right\}^2s^2\right]^8}\frac{\mathrm{\Gamma }\left(8\right)}{\mathrm{\Gamma }\left(7\right)}\right]f_{1a}\left(k\right)}`$ $`=`$ $`{\displaystyle \frac{i}{8}}{\displaystyle _0^1}dx{\displaystyle }{\displaystyle \frac{d^4l}{\left(2\pi \right)^4}}{\displaystyle \frac{1}{\left[l^2s^2\right]^8}}{\displaystyle \frac{\mathrm{\Gamma }\left(8\right)}{\mathrm{\Gamma }\left(7\right)}}[\{{\displaystyle \frac{28}{3}}l^8+({\displaystyle \frac{116}{3}}m^2{\displaystyle \frac{220}{3}}q^2t^2)l^6`$ $`+({\displaystyle \frac{148}{3}}m^4+{\displaystyle \frac{416}{3}}m^2q^2t^2{\displaystyle \frac{268}{3}}q^4t^4)l^4+(20m^660m^4q^2t^2+60m^2q^4t^420q^6t^6)l^2\}A_{\mu \nu }`$ $`+\{(640t^3+896t^4)l^4+(640m^2t^3768m^2t^4)l^2`$ $`128m^4(t^3t^4)+256m^2q^2(t^5t^6)128q^4(t^7t^8)\}q_\mu q_\nu q_\alpha q_\beta A_{\alpha \beta }`$ $`+\{{\displaystyle \frac{140}{3}}tl^6+({\displaystyle \frac{464}{3}}m^2t+{\displaystyle \frac{400}{3}}m^2t^2+{\displaystyle \frac{80}{3}}q^2t^3)l^4`$ $`+\left(148m^4t112m^4t^2+96m^2q^2t^3+296m^2q^2t^4+116q^4t^5184q^4t^6\right)l^2`$ $`+(40t+24t^2)m^6(56q^2t^372q^2t^4)m^4+(8q^4t^5+72q^4t^6)m^2+24q^6(t^7t^8)\}g_{\mu \nu }q_\alpha q_\beta A_{\alpha \beta }`$ $`+\{({\displaystyle \frac{140}{3}}t40t^2)l^6+({\displaystyle \frac{464}{3}}m^2t+{\displaystyle \frac{560}{3}}m^2t^2)l^4`$ $`+\left(148m^4t176m^4t^2416m^2q^2t^3+472m^2q^2t^4+268q^4t^5296q^4t^6\right)l^2`$ $`40m^6(tt^2)+120m^4q^2(t^3t^4)120m^2q^4(t^5t^6)+40q^6(t^7t^8)\}(q_\mu q_\alpha A_{\alpha \nu }+q_\nu q_\alpha A_{\alpha \mu })]`$ $`=`$ $`A_{\mu \nu }{\displaystyle \frac{\pi ^2}{Q^4}}\left({\displaystyle \frac{7}{3}}y+\left({\displaystyle \frac{1}{3}}+{\displaystyle \frac{1}{6}}y\right)J_1{\displaystyle \frac{8}{3}}J_2\right)`$ $`+g_{\mu \nu }q_\alpha q_\beta A_{\alpha \beta }{\displaystyle \frac{\pi ^2}{Q^6}}\left({\displaystyle \frac{5}{3}}+y\left(2+y\right)J_19J_2+{\displaystyle \frac{52}{3}}J_38J_4\right)`$ $`+q_\mu q_\nu q_\alpha q_\beta A_{\alpha \beta }{\displaystyle \frac{\pi ^2}{Q^8}}\left({\displaystyle \frac{68}{3}}\left(8+4y\right)J_136J_2+{\displaystyle \frac{88}{3}}J_38J_4\right)`$ $`+\left(q_\mu q_\alpha A_{\alpha \nu }+q_\nu q_\alpha \right)A_{\alpha \mu }{\displaystyle \frac{\pi ^2}{Q^6}}\left({\displaystyle \frac{23}{3}}y\left({\displaystyle \frac{2}{3}}+{\displaystyle \frac{1}{3}}y\right)J_1{\displaystyle \frac{41}{3}}J_2+{\displaystyle \frac{20}{3}}J_3\right)`$ ## Appendix G Moments in terms of Hypergeometric Functions $`{}_{2}{}^{}F_{1}^{}`$ In order to obtain moments written in terms of hypergeometric functions $`{}_{2}{}^{}F_{1}^{}`$ from polarizations eq.(90), we expand the polarizations in terms of $`Q^2/m^2`$. Here we divide $`a_0`$ and $`a_1`$ into $`a_0=d_1+d_2y`$ and $`a_1=c_1+c_2y`$. The other higher coefficients, $`a_2,a_3,\mathrm{},a_k`$ don’t contain $`y`$, where $`y=Q^2/m^2`$. $`\mathrm{\Pi }\left(q\right)`$ $`=`$ $`Q^a[a_0+(c_1+c_2y){\displaystyle _0^1}{\displaystyle \frac{dx}{\left[1+x\left(1x\right)Q^2/m^2\right]}}+a_2{\displaystyle _0^1}{\displaystyle \frac{dx}{\left[1+x\left(1x\right)Q^2/m^2\right]^2}}+`$ (95) $`\mathrm{}`$ $`+a_k{\displaystyle _0^1}{\displaystyle \frac{dx}{\left[1+x\left(1x\right)Q^2/m^2\right]^k}}]`$ $`=`$ $`Q^a[a_0+(c_1+c_2y){\displaystyle _0^1}dx{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\{x(1x)\}^j{\displaystyle \frac{Q^{2j}}{m^{2j}}}`$ $`+a_2{\displaystyle _0^1}𝑑x{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\left\{x\left(1x\right)\right\}^j{\displaystyle \frac{\left(1+j\right)!}{j!}}{\displaystyle \frac{Q^{2j}}{m^{2j}}}+`$ $`\mathrm{}`$ $`+a_k{\displaystyle _0^1}dx{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}\{x(1x)\}^j{\displaystyle \frac{\left(k1+j\right)!}{\left(k1\right)!j!}}{\displaystyle \frac{Q^{2j}}{m^{2j}}}]`$ $`=`$ $`Q^a{\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}[c_1{\displaystyle _0^1}dx(1)^jx^j(1x)^j{\displaystyle \frac{Q^{2j}}{m^{2j}}}+a_2{\displaystyle _0^1}dx(1)^jx^j(1x)^j(1+j){\displaystyle \frac{Q^{2j}}{m^{2j}}}+`$ $`\mathrm{}`$ $`+a_k{\displaystyle _0^1}dx(1)^jx^j(1x)^j{\displaystyle \frac{\left(k1+j\right)!}{\left(k1\right)!j!}}{\displaystyle \frac{Q^{2j}}{m^{2j}}}c_2{\displaystyle _0^1}dx(1)^jx^{j1}(1x)^{j1}{\displaystyle \frac{Q^{2j}}{m^{2j}}}]`$ $`=`$ $`{\displaystyle \underset{j=a/2}{\overset{\mathrm{}}{}}}()^{a/21}{\displaystyle _0^1}dx\left\{x(1x)\right\}^{j1}{\displaystyle \frac{\left(Q^2\right)^{ja/2}}{m^{2j}}}[x(1x)\{c_1+a_2(1+j)`$ $`+{\displaystyle \frac{a_3}{2!}}(1+j)(2+j)+\mathrm{}+{\displaystyle \frac{a_k}{\left(k1\right)!}}{\displaystyle \underset{l=1}{\overset{k1}{}}}(l+j)\}+c_2]`$ Now the $`n`$’s moment $`M_n`$ can be obtained from above by differentiation. $`M_n\left(Q_0\right)`$ $`=`$ $`{\displaystyle \frac{1}{n!}}\left({\displaystyle \frac{d}{dQ^2}}\right)^n\mathrm{\Pi }\left(Q^2\right)`$ (104) $`=`$ $`()^{a/2}{\displaystyle \frac{1}{\left(4m^2\right)^{a/2+n}}}{\displaystyle \frac{1}{\left(1+\xi \right)^{n+1}}}[c_1d_1+c_2d_2`$ $`+{\displaystyle \frac{1}{2}}\{a_2{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+\frac{a}{2}+2\right)}{\mathrm{\Gamma }\left(n+\frac{a+3}{2}\right)}}_2F_1(n+1,{\displaystyle \frac{1}{2}},n+{\displaystyle \frac{a+3}{2}};{\displaystyle \frac{\xi }{1+\xi }})`$ $`+{\displaystyle \frac{a_3}{2!}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+\frac{a}{2}+3\right)}{\mathrm{\Gamma }\left(n+\frac{a+3}{2}\right)}}_2F_1(n+1,{\displaystyle \frac{3}{2}},n+{\displaystyle \frac{a+3}{2}};{\displaystyle \frac{\xi }{1+\xi }})`$ $`+{\displaystyle \frac{a_4}{3!}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+\frac{a}{2}+4\right)}{\mathrm{\Gamma }\left(n+\frac{a+3}{2}\right)}}_2F_1(n+1,{\displaystyle \frac{5}{2}},n+{\displaystyle \frac{a+3}{2}};{\displaystyle \frac{\xi }{1+\xi }})`$ $`+\mathrm{}`$ $`+{\displaystyle \frac{a_k}{\left(k1\right)!}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+\frac{a}{2}+k\right)}{\mathrm{\Gamma }\left(n+\frac{a+3}{2}\right)}}{}_{2}{}^{}F_{1}^{}(n+1,k+{\displaystyle \frac{3}{2}},n+{\displaystyle \frac{a+3}{2}};{\displaystyle \frac{\xi }{1+\xi }})\}`$ $`=`$ $`{\displaystyle \frac{1}{\left(4m^2\right)^{a/2+n}}}{\displaystyle \frac{1}{\left(1+\xi \right)^{n+1}}}{\displaystyle \frac{()^{a/2}\mathrm{\Gamma }\left(1/2\right)}{2\mathrm{\Gamma }\left(n+\frac{a+3}{2}\right)}}\mathrm{\Gamma }\left(n+{\displaystyle \frac{a}{2}}\right)\left(\begin{array}{c}c_1,c_2,a_2,a_3,\mathrm{},a_k\end{array}\right)`$ $`\times \left(\begin{array}{cc}\left(n+\frac{a}{2}\right)/0!{}_{2}{}^{}F_{1}^{}(n+1,\frac{1}{2},n+\frac{a+3}{2};\rho )& \\ 2\left(2n+a+1\right){}_{2}{}^{}F_{1}^{}(n+1,\frac{1}{2},n+\frac{a+1}{2};\rho )& \\ \left(n+\frac{a}{2}\right)\left(n+\frac{a+2}{2}\right)/1!{}_{2}{}^{}F_{1}^{}(n+1,\frac{1}{2},n+\frac{a+3}{2};\rho )& \\ \left(n+\frac{a}{2}\right)\left(n+\frac{a+2}{2}\right)\left(n+\frac{a+4}{2}\right)/2!{}_{2}{}^{}F_{1}^{}(n+1,\frac{3}{2},n+\frac{a+3}{2};\rho )& \\ \mathrm{}& \\ \left(n+\frac{a}{2}\right)\left(n+\frac{a+2}{2}\right)\mathrm{}\left(n+\frac{a}{2}+k1\right)/\left(k1\right)!{}_{2}{}^{}F_{1}^{}(n+1,k+\frac{3}{2},n+\frac{a+3}{2};\rho )& \end{array}\right)`$ where $`\left(\begin{array}{c}d_1\\ d_2\end{array}\right)=\left(\begin{array}{cc}\frac{1}{2}\frac{\mathrm{\Gamma }[1/2]\mathrm{\Gamma }[n+(a+2)/2]}{\mathrm{\Gamma }[n+(a+3)/2]}_2F_1(n+1,\frac{1}{2},n+\frac{a+3}{2};\frac{\xi }{1+\xi })& \\ 2\frac{\mathrm{\Gamma }[1/2]\mathrm{\Gamma }[n+a/2]}{\mathrm{\Gamma }[n+(a+1)/2]}_2F_1(n+1,\frac{1}{2},n+\frac{a+1}{2};\frac{\xi }{1+\xi })& \end{array}\right).`$ (110) As an example, when $`\xi =0`$ and $`a=6`$, we have, $`M_n(Q_0^2,\xi `$ $`=`$ $`0)={\displaystyle \frac{\left(1\right)}{\left(4m_c^2\right)^{n+3}}}[c_1{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+4\right)}{\mathrm{\Gamma }\left(n+\frac{9}{2}\right)}}+c_2(2){\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+3\right)\left(n+\frac{7}{2}\right)}{\mathrm{\Gamma }\left(n+\frac{9}{2}\right)}}+`$ (118) $`+{\displaystyle \frac{1}{2}}\{a_2{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+5\right)}{\mathrm{\Gamma }\left(n+\frac{9}{2}\right)}}+{\displaystyle \frac{a_3}{2!}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+6\right)}{\mathrm{\Gamma }\left(n+\frac{9}{2}\right)}}+\mathrm{}+{\displaystyle \frac{a_k}{\left(k1\right)!}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(n+k+3\right)}{\mathrm{\Gamma }\left(n+\frac{9}{2}\right)}}\}]`$ $`=`$ $`{\displaystyle \frac{\left(1\right)}{\left(4m_c^2\right)^{n+3}}}{\displaystyle \frac{2^{n+3}\left(n+2\right)!}{\left(2n+7\right)!!}}\left(\begin{array}{c}c_1,c_2,a_2,a_3,\mathrm{},a_k\end{array}\right)\left(\begin{array}{cc}\left(n+3\right)/0!& \\ 2\left(2n+7\right)& \\ \left(n+3\right)\left(n+4\right)/1!& \\ \left(n+3\right)\left(n+4\right)\left(n+5\right)/2!& \\ \mathrm{}& \\ \frac{\left(n+2+k\right)!}{\left(n+2\right)!}/\left(k1\right)!& \end{array}\right)`$ $`=`$ $`{\displaystyle \frac{M_n^0}{\left(4m_c^2\right)^3}}\left[{\displaystyle \frac{2^5\pi ^2}{3}}{\displaystyle \frac{n\left(n+2\right)}{\left(2n+5\right)\left(2n+7\right)}}\right]\left(\begin{array}{c}c_1,c_2,a_2,\mathrm{},a_k\end{array}\right)\left(\begin{array}{cc}\left(n+3\right)/0!& \\ 2\left(2n+7\right)& \\ \left(n+3\right)\left(n+4\right)/1!& \\ \left(n+3\right)\left(n+4\right)\left(n+5\right)/2!& \\ \mathrm{}& \\ \frac{\left(n+2+k\right)!}{\left(n+2\right)!}/\left(k1\right)!& \end{array}\right),`$ (126) where $`M_n^0=3\times 2^n\left(n+1\right)\left(n1\right)!/\left(4\pi ^2\left(2n+3\right)!!\left(4m_c^2\right)^n\right)=A_n\left(\xi =0\right).`$ ## Appendix H Consistency check As discussed in the text, there are few checks we can perform to confirm our calculation. The first is the current conservation, which we checked explicitly. The second is the regularity at $`Q^2=0`$. This can be checked from the following equation. $`\mathrm{\Pi }\left(q\right)`$ $`=`$ $`{\displaystyle \underset{j=a/2}{\overset{\mathrm{}}{}}}()^{a/21}{\displaystyle _0^1}dx\left\{x(1x)\right\}^{j1}{\displaystyle \frac{\left(Q^2\right)^{ja/2}}{m^{2j}}}[x(1x)\{c_1+a_2(1+j)`$ (128) $`+{\displaystyle \frac{a_3}{2!}}(1+j)(2+j)+\mathrm{}+{\displaystyle \frac{a_k}{\left(k1\right)!}}{\displaystyle \underset{l=1}{\overset{k1}{}}}(l+j)\}+c_2]`$ holds only after making the following checks: * $`j=0`$ $$0=d_1+c_1+a_2+a_3+\mathrm{}+a_k$$ * $`j=1`$ $`0=c_2+d_2+{\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}\left\{c_1+2a_2+3a_3+\mathrm{}+ka_k\right\}`$ * $`j=2`$ $`0={\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}c_2+{\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}^2\left\{c_1+\mathrm{}+{\displaystyle \frac{k\left(k+1\right)}{2}}a_k\right\}`$ * $`j=3`$ $`0={\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}^2c_2+{\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}^3\left\{c_1+\mathrm{}+{\displaystyle \frac{k\left(k+1\right)\left(k+2\right)}{3}}a_k\right\}`$ * $`j=\frac{a}{2}1`$ $`0={\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}^{\frac{a}{2}2}c_2`$ $`+{\displaystyle _0^1}𝑑x\left\{x\left(1x\right)\right\}^{\frac{a}{2}1}\left\{c_1+\mathrm{}+{\displaystyle \frac{k\mathrm{}\left(k+\frac{a}{2}3\right)\left(k+\frac{a}{2}2\right)}{\frac{a}{2}1}}a_k\right\}`$ So we know the summation begins at $`j=a/2`$. Substituting the values for the coefficients multiplying $`J`$’s, we have explicitly checked that the above constraints are satisfied in our calculation.
warning/0002/hep-lat0002011.html	ar5iv	text	# Lattice monopole action in pure SU(3) QCD ## 1 Introduction The quark confinement is a key problem to understand non-perturbative phenomena of QCD. The dual Meissner effect is a promising candidate . ’t Hooft suggested the idea of abelian projection . $`SU(N)`$ QCD is reduced to a $`U(1)^{N1}`$ abelian gauge theory with magnetic monopoles by a partial gauge fixing (abelian projection). The dual Meissner effect is caused by the monopole condensation. Monopoles are responsible for the confinement as in compact QED \[3-8\]. Monte-Carlo simulations of lattice QCD are a most powerful method to study the non-perturbative phenomena. Numerical studies of the abelian projection in the maximally abelian (MA) gauge have confirmed the ’t Hooft conjecture. In MA gauge , the string tension derived from abelian Wilson loops gives almost the same value as that of non-abelian Wilson loop (abelian dominance) \[11-13\]. Moreover, the monopole contribution to the abelian Wilson loops alone reproduces the string tension in SU(2) QCD (monopole dominance) \[14-16\]. The abelian and the monopole dominances are seen also in the behavior of the Polyakov loop in $`T0`$ SU(2) and SU(3) QCD \[17-19\]. These results support the ’t Hooft conjecture. Wilson’s idea of a block-spin transformation and a renormalized trajectory (RT) are useful when we study the continuum limit on available lattices . The lattice action on RT has no lattice artifact and hence it is called quantum perfect action. It reproduces the same physical results as in the continuum limit. It is challenging to get the perfect lattice action for the infrared region of QCD. For that purpose, we have to extract a dynamical variable which plays a dominant role in the infrared region. The numerical evidence for the monopole dominance suggests that low-energy QCD can be described by an effective action on the dual lattice in terms of a dual quantity like monopole currents. It is very interesting to obtain such an effective monopole action. Also it is challenging to set the perfect monopole action with the help of block-spin transformations for monopole currents. The monopole action can be obtained numerically by an inverse Monte-Carlo method . A block-spin transformation on the dual lattice can be performed by considering an $`n`$-blocked monopole current . The detailed study of the effective monopole action in pure SU(2) QCD has been done \[23-25\]. The monopole action determined in SU(2) has the following important features: 1. Two-point interactions are dominant in the infrared region and coupling constants decrease rapidly as the distance between two monopole currents increases. 2. The action fixed seems to satisfy the scaling behavior, that is, it depends only on a physical scale $`b=na(\beta )`$. This suggests that the action is near to RT. 3. Monopole condensation seems to occur for large $`b`$ region from the energy-entropy balance. In order to test the validity of the statement that the action is near to RT, we need to check the restoration of the continuum rotational invariance. Then, we have to determine first the correct forms of physical operators (quantum perfect operators) on the blocked lattice. In the above pure SU(2) study, Fujimoto et al. have taken the following steps: 1. Study the renormalization flow on the projected space of two-point monopole interactions alone, considering that they are dominant numerically in the infrared region. 2. The static potential between quark and antiquark can be evaluated by the expectation value of the Wilson loop in the continuum limit. Hence the Wilson loop can be regarded as the correct operator evaluating the static potential on the fine $`a`$-lattice also. 3. Perform the block spin transformation analytically, starting from the two-point monopole interactions with the Wilson loop on the fine $`a`$-lattice. The quantum perfect action and the quantum perfect operator on the coarse $`b`$-lattice can be obtained analytically when we take the $`a0`$ limit for fixed $`b=na`$. 4. Compare the quantum perfect action composed of two-point monopole interactions with the effective action numerically determined from the inverse Monte-Carlo method. The parameters in the perfect monopole action and the quantum perfect operator are fixed then. 5. The effective action with the quantum perfect operator can be transformed into that of the string model. Since the strong-coupling expansion is found to work well in the string model, we see that the static potential is estimated analytically by the classical part alone and the continuum rotational invariance is restored. It is not so straightforward to extend these studies to pure SU(3) QCD, since there are three monopole currents $`k_\mu ^{(a)}(s)`$ satisfying one constraint $`_{a=1}^3k_\mu ^{(a)}(s)=0`$. So far, SU(3) monopole action composed of only one kind of monopole current after integrating out the other two has been derived in the case of two-point interactions . In this case, the same method can be applied as done in pure SU(2) QCD and monopole condensation seems to occur from the energy-entropy balance in rather strong coupling region. But the scaling seen in the pure SU(2) case was not seen. Also it is important to study effective action for two independent monopole currents in order to see the characteristic features of pure SU(3) QCD. The purpose of this paper is to report the results of the extensive studies of SU(3) monopole actions, especially in terms of two independent monopole currents. In Section 2 we briefly review the SU(3) monopole current on the lattice and the inverse Monte-Carlo method. The numerical results of the SU(3) monopole action are shown in Section 3. In Section 4 we perform a block-spin transformation of the monopole current analytically, considering only the infrared dominant two-point monopole interactions. In Section 5, transforming the monopole action into that of the string model, we calculate the static potential analytically. In Section 6 the string tension is estimated from the results of previous Sections. The conclusions are given in Section 7. The correspondence between the monopole action and the dual abelian Higgs theory (dual Ginzburg-Landau theory) is given in Appendix. ## 2 Monopole current and the inverse Monte-Carlo method We extract a $`U(1)^2`$ link field $`u_\mu (s)`$ from a SU(3) link field $`U_\mu (s)`$ after abelian projection called MA gauge. In the SU(2) case, a $`U(1)`$ link field is defined as $`u_\mu (s)=\mathrm{diag}(\mathrm{e}^{\mathrm{i}\theta _\mu ^{(1)}},\mathrm{e}^{\mathrm{i}\theta _\mu ^{(2)}})`$, where $`\theta _\mu ^{(a)}(s)\mathrm{arg}[U_\mu (s)]_{aa}(a=1,2)`$. It satisfies $`det(u_\mu (s))=1`$ due to $`_{a=1}^2\theta _\mu ^{(a)}=0`$. Since $`_{a=1}^3\theta _\mu ^{(a)}0`$ in the SU(3) case, the definition of $`U(1)^2`$ link field necessarily becomes more complicated. In this study, we use the definition of Ref. .<sup>1</sup><sup>1</sup>1 Another definition is seen in Ref. . Both definitions are equivalent in the continuum limit. The fields transforming as photon fields under $`U(1)^2`$ are defined as follows: $$\theta _\mu ^{(a)}\mathrm{arg}[U_\mu ]_{aa}\frac{1}{3}\varphi _\mu ,\varphi _\mu \underset{a=1}{\overset{3}{}}\mathrm{arg}[U_\mu ]_{aa}\|_{\mathrm{mod2}\pi }[\pi ,\pi ).$$ (1) The $`U(1)^2`$ link field defined by $`u_\mu (s)=\mathrm{diag}(\mathrm{e}^{\mathrm{i}\theta _\mu ^{(1)}},\mathrm{e}^{\mathrm{i}\theta _\mu ^{(2)}},\mathrm{e}^{\mathrm{i}\theta _\mu ^{(3)}})`$ satisfies $`det(u_\mu (s))=1`$. If the $`U(1)^2`$ field strength is defined as $`\overline{\mathrm{\Theta }}_{\mu \nu }^{(a)}_\mu \theta _\nu ^{(a)}_\nu \theta _\mu ^{(a)}(\mathrm{mod2}\pi )`$, then $`_a\overline{\mathrm{\Theta }}_{\mu \nu }^{(a)}=2\pi l(l=0,\pm 1)`$ and is not always zero. When $`l=+1(1)`$$`\overline{\mathrm{\Theta }}_{\mu \nu }^{(a)}`$ is redefined. If $`\overline{\mathrm{\Theta }}_{\mu \nu }^{(a)}`$ is the maximum(minimum) of $`(\overline{\mathrm{\Theta }}_{\mu \nu }^{(1)},\overline{\mathrm{\Theta }}_{\mu \nu }^{(2)},\overline{\mathrm{\Theta }}_{\mu \nu }^{(3)})`$, it is redefined as $`\overline{\mathrm{\Theta }}_{\mu \nu }^{(a)}2\pi (+2\pi )`$. Others do not change. Then new $`U(1)^2`$ field strengths satisfy $`_a\overline{\mathrm{\Theta }}_{\mu \nu }^{(a)}=0`$. Monopole currents are defined as $`k_\mu ^{(a)}1/4\pi ϵ_{\mu \nu \rho \sigma }_\nu \overline{\mathrm{\Theta }}_{\rho \sigma }^{(a)}`$ by DeGrand-Toussaint(D-T) which satisfies the constraint $$\underset{a}{}k_\mu ^{(a)}=0.$$ (2) We want to get an effective monopole action on the dual lattice integrating out the degrees of freedom except for the monopole currents: $`Z`$ $`=`$ $`{\displaystyle 𝒟U\delta (X^{off})\mathrm{\Delta }_F(U)e^{S(U)}}`$ (3) $`=`$ $`{\displaystyle 𝒟u[𝒟c\delta (X^{off})\mathrm{\Delta }_F(U)e^{S(U)}]}`$ (4) $`=`$ $`{\displaystyle 𝒟ue^{S_{eff}(u)}}`$ (5) $`=`$ $`{\displaystyle \underset{k^{(a)}Z}{}}\delta _{_\mu ^{}k_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}{\displaystyle 𝒟u\delta (k^{(a)},u)e^{S_{eff}(u)}}`$ (6) $`=`$ $`{\displaystyle \underset{k^{(a)}Z}{}}\delta _{_\mu ^{}k_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}e^{S[k^{(a)}]},`$ (7) where $`U_\mu (s)=c_\mu (s)u_\mu (s)`$ and $`X^{off}`$ are the off-diagonal part of the following quantity: $$X(s)=\underset{\mu ,a}{}[U_\mu (s)\mathrm{\Lambda }_aU_\mu ^{}(s)+U_\mu ^{}(s\widehat{\mu })\mathrm{\Lambda }_aU_\mu (s\widehat{\mu }),\mathrm{\Lambda }_a],$$ (8) $`\mathrm{\Lambda }_1`$ $`=`$ $`\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right),\mathrm{\Lambda }_2=\left(\begin{array}{ccc}1& 0& 0\\ 0& 0& 0\\ 0& 0& 1\end{array}\right),\mathrm{\Lambda }_3=\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).`$ (18) $`\mathrm{\Delta }_F(U)`$ is the Faddeev-Popov determinant and $`\delta (k^{(a)},u)`$ is the delta function corresponding to the D-T definition of the monopole current. We can perform the above integration numerically. We create vacuum ensembles of monopole currents $`\{k_\mu ^{(a)}(s)\}`$ using the Monte-Carlo method and the above definition of the monopole current. Then we derive the effective monopole action using the Swendsen method which is one of the inverse Monte-Carlo methods. In order to take the continuum limit, we perform the block-spin transformation on the dual lattice by defining an $`n`$-blocked monopole current : $`K_\mu ^{(a)}(s^{(n)})`$ $``$ $`{\displaystyle \underset{i,j,m=0}{\overset{n1}{}}}k_\mu ^{(a)}(ns+(n1)\widehat{\mu }+i\widehat{\nu }+j\widehat{\rho }+m\widehat{\sigma }).`$ (19) If the action numerically obtained satisfies a scaling behavior $`S[K^{(a)},n,a(\beta )]=S[K^{(a)},b=na(\beta )]`$, that is, the action depends only on $`b`$, the continuum limit can be taken as $`n\mathrm{},a(\beta )0`$ for a fixed physical length $`b=na(\beta )`$. The original Swendsen method must be extended due to the conservation law of the monopole current $`_\mu ^{}k_\mu ^{(a)}(s)=0`$ as in the SU(2) case . Let us assume the form of the monopole action as $`S[k^{(a)}]=_iG_iS_i[k^{(a)}]`$ ($`G_i`$ is the coupling constant for the operator $`S_i[k^{(a)}]`$) and define $`\stackrel{~}{S}_i[k^{(a)}]`$ as a part of $`S_i[k^{(a)}]`$ which contains the currents around a specific plaquette $`(s^{},\widehat{\mu }^{},\widehat{\nu }^{})`$. When we consider the expectation value of some operator $`O_i[k^{(a)}]`$, the following identity holds as in the SU(2) case: $$O_i[k^{(a)}]=\overline{O}_i[k^{(a)}],$$ (20) $$O_i[k^{(a)}]\frac{_{k^{(a)}Z}\delta _{_\mu ^{}k_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}O_i[k^{(a)}]\mathrm{exp}(_jG_jS_j[k^{(a)}])}{_{k^{(a)}Z}\delta _{_\mu ^{}k_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}\mathrm{exp}(_jG_jS_j[k^{(a)}])},$$ (21) $$\overline{O}_i[k^{(a)}]\frac{_{M^{(a)}Z}\delta _{\mathrm{\Sigma }_aM^{(a)},0}O_i[k^{(a)}]\mathrm{exp}(_jG_j\stackrel{~}{S}_j[k^{(a)}])}{_{M^{(a)}Z}\delta _{\mathrm{\Sigma }_aM^{(a)},0}\mathrm{exp}(_jG_j\stackrel{~}{S}_j[k^{(a)}])},$$ (22) $$k_{}^{}{}_{\mu }{}^{(a)}(s)k_\mu ^{(a)}(s)+M^{(a)}(\delta _{s,s^{}}\delta _{\mu ,\mu ^{}}+\delta _{s,s^{}+\widehat{\mu }^{}}\delta _{\mu ,\nu ^{}}\delta _{s,s^{}+\widehat{\nu }^{}}\delta _{\mu ,\mu ^{}}\delta _{s,s^{}}\delta _{\mu ,\nu ^{}}).$$ (23) We use this identity to determine the values of the couplings $`G_i`$ iteratively. Choose trial couplings $`\stackrel{~}{G}_i`$ suitably. If $`\stackrel{~}{G}_i`$ are not equal to $`G_i`$ for all $`i`$, then $`\mathrm{\Delta }\stackrel{~}{G}_iG_i\stackrel{~}{G}_i`$ are estimated from the expansion of Eq. (2.11) up to the first order of $`\mathrm{\Delta }\stackrel{~}{G}_i`$: $$O_i\overline{O}_i\overline{O}_i\overline{S}_j\overline{O_iS_j}\mathrm{\Delta }\stackrel{~}{G}_j.$$ (24) $`\stackrel{~}{G^{}}_i=\stackrel{~}{G}_i+\mathrm{\Delta }\stackrel{~}{G}_i`$ are used as the next trial couplings. These procedures continue iteratively until the couplings converge. ## 3 Numerical results Practically, we must truncate the form of the monopole action to derive it numerically. We know that short-distant and two-point interactions are dominant in the SU(2) case. Here, we assume the following form of the SU(3) monopole action: $`S[k^{(a)}]`$ $`=`$ $`{\displaystyle \underset{i}{}}G_i(S_i[k^{(1)}]+S_i[k^{(2)}]+S_i[k^{(3)}])`$ (25) $`=`$ $`{\displaystyle \underset{i}{}}G_i(S_i[k^{(1)}]+S_i[k^{(2)}]+S_i[k^{(1)}k^{(2)}]).`$ (26) The Weyl symmetry (the species permutation symmetry of the monopole currents) remains after abelian projection. We adopt 27 two-point interactions whose distances are up to $`3na(\beta )`$ and 4, 6-point interactions of the following simple form: $$\underset{a,s}{}\{\underset{\mu =4}{\overset{4}{}}(k_\mu ^{(a)}(s))^2\}^2,\underset{a,s}{}\{\underset{\mu =4}{\overset{4}{}}(k_\mu ^{(a)}(s))^2\}^3.$$ (27) The lattice size is $`48^4`$ from $`\beta =5.6`$ to $`\beta =6.4`$. After thermalization, 30 configurations are used for the average. We determine the lattice spacing $`a(\beta )`$ without using the theoretical asymptotic beta function. It is given by the relation $`a(\beta )=\sqrt{\sigma (\beta )/\sigma _{ph}}`$ , where $`\sigma _{ph}`$ is the physical string tension <sup>2</sup><sup>2</sup>2If we use $`\sigma _{ph}^{1/2}=440`$ \[MeV\], $`b=2\sigma _{ph}^{1/2}`$ is approximately equal to 0.9 \[fm\].. The results are summarized as follows: 1. The monopole action for two independent types of monopole currents is obtained clearly (Figure 1). The qualitative behaviors are the same as in the SU(2) monopole action: 1. The monopole action has a compact form. The self-interaction $`G_1_{a,s,\mu }(k_\mu ^{(a)}(s))^2`$ is dominant and the coupling constants $`G_i`$ decrease rapidly as the distance between the two monopole currents increases. 2. The coupling constants have a direction dependence which is expected after blocking. Two nearest-neighbor interactions $`G_2_{a,s,\mu }k_\mu ^{(a)}(s)k_\mu ^{(a)}(s+\widehat{\mu })`$ and $`G_3_{a,s,\mu \nu }k_\mu ^{(a)}(s)k_\mu ^{(a)}(s+\widehat{\nu })`$ are quite different for small $`b`$ region. 3. The coupling constants $`G_i`$ become very small for large $`b`$ region. 4. The simple 4, 6-point interactions become negligibly small for large $`b>2\sigma _{ph}^{1/2}`$ . Two-point interactions are relatively dominant for large $`b`$ region. 5. The scaling behavior holds well for $`n=4,6,8`$ data, if the physical scale $`b=na(\beta )`$ is taken in unit of the string tension $`\sqrt{\sigma _{ph}}`$. The action seems to be very near to RT on which one can take the continuum limit. 2. In order to study if monopole condensation occurs by energy-entropy balance, we derive the monopole action, considering only one type of the monopole current. For simplicity, only two-point interactions are taken into account. The scaling was not seen in the previous study where the two loop perturbation value was used for $`a(\beta )`$ . When $`a(\beta )`$ is fixed by the string tension, the scaling is seen beautifully in Figure 2. Since we restrict ourselves to one type of the monopole current, the entropy of the monopole current loop is given approximately by $`\mathrm{ln}7\times L`$ ($`L`$ is a length of one long loop). From the previous study, we know that only one long loop and some short loops of monopoles exist in the vacuum and the value of the action is well approximated as $`G_1\times L`$. Figure 2 shows that the entropy dominates over the energy in the large $`b`$ region, i.e., $`G_1(b)<\mathrm{ln}7`$. ## 4 Perfect action and perfect operator for monopole current The monopole action seems to satisfy the scaling behavior. The scaling is the characteristic of the perfect action. The perfect action reproduces the continuum rotational invariance. To test the rotational invariance, we have to determine the form of physical operators on the blocked lattice. For that purpose, as done in the SU(2) case, we perform a block-spin transformation from the small $`a`$-lattice ($`a0`$) to the finite $`b=na`$ lattice analytically, restricting ourselves to a simple case of two-point monopole interactions with a monopole Wilson loop. Note that two-point interactions are dominant for large $`b`$ region also in pure SU(3). We start from the following action on the $`a`$-lattice: $`Z[J]`$ $`=`$ $`{\displaystyle \underset{k^{(a)}Z}{}}\delta _{_\mu ^{}k_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}\mathrm{exp}\{{\displaystyle \underset{a,s,s^{}}{}}k_\mu ^{(a)}(s)D_0(ss^{})k_\mu ^{(a)}(s^{})`$ (28) $`+2\pi i{\displaystyle \underset{s}{}}N_\mu (s)k_\mu ^{(1)}(s)\},`$ $`N_\mu (s)`$ $`=`$ $`{\displaystyle \underset{s^{}}{}}\mathrm{\Delta }^1(ss^{}){\displaystyle \frac{1}{2}}ϵ_{\mu \alpha \beta \gamma }_\alpha S_{\beta \gamma }(s^{}+\widehat{\mu }),_{}^{}{}_{\alpha }{}^{}S_{\alpha \beta }(s)=J_\beta (s).`$ (29) Here $`J_\beta (s)`$ is an electric current around the Wilson loop. Note that in Eq. (4.1), the static potential between quark and antiquark with $`a=1`$ is considered. We have used the monopole Wilson loop operator as in Ref. . As the surface $`S_{\alpha \beta }`$, we can take any open surface with the fixed boundary $`C`$, since the monopole Wilson loop operator with the closed surface is unity due to 4-dimensional linking number in the continuum limit. Here we take the flat surface as $`S_{\alpha \beta }`$. We expand $`D_0(s)`$ and adopt the following first three terms for simplicity: $`D_0(ss^{})`$ $`=`$ $`\alpha \mathrm{\Delta }^1(ss^{})+\beta \delta _{s,s^{}}+\gamma \mathrm{\Delta }(ss^{}).`$ (30) The above form of the monopole action is derived from the dual Ginzburg-Landau (DGL) theory also (for details, see Appendix). We perform the block-spin transformation, using the definition of $`n`$-blocked monopole currents $`K^{(a)}`$ in Eq. (2.10): $`Z[K^{(a)},J]`$ $`=`$ $`{\displaystyle \underset{k^{(a)}Z}{}}\delta _{_\mu ^{}k_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}\mathrm{exp}\{{\displaystyle \underset{a,s,s^{}}{}}k_\mu ^{(a)}(s)D_0(ss^{})k_\mu ^{(a)}(s^{})`$ (32) $`+2\pi i{\displaystyle \underset{s}{}}N_\mu (s)k_\mu ^{(1)}(s)\}`$ $`\times {\displaystyle \underset{a}{}}\delta (K_\mu ^{(a)}(s^{(n)}){\displaystyle \underset{i,j,k=0}{\overset{n1}{}}}k_\mu ^{(a)}(ns^{(n)}+(n1)\widehat{\mu }+i\widehat{\nu }+j\widehat{\rho }+k\widehat{\sigma }))`$ $`=`$ $`{\displaystyle _\pi ^\pi }𝒟𝜸{\displaystyle _\pi ^\pi }𝒟𝑩{\displaystyle _\pi ^\pi }𝒟\varphi {\displaystyle \underset{k^{(a)}Z}{}}\mathrm{exp}\{{\displaystyle \underset{a,s,s^{}}{}}k_\mu ^{(a)}(s)D_0(ss^{})k_\mu ^{(a)}(s^{})`$ $`+2\pi i{\displaystyle \underset{s}{}}N_\mu (s)k_\mu ^{(1)}(s)+i{\displaystyle \underset{a,s}{}}B^{(a)}(s)_{\mu }^{}{}_{}{}^{}k_\mu ^{(a)}(s)+i{\displaystyle \underset{s}{}}\varphi _\mu (s){\displaystyle \underset{a}{}}k_\mu ^{(a)}(s)`$ $`+i{\displaystyle \underset{a,s^{(n)}}{}}\gamma _\mu ^{(a)}(s^{(n)})(K_\mu ^{(a)}(s^{(n)}){\displaystyle \underset{i,j,k=0}{\overset{n1}{}}}k_\mu ^{(a)})\}`$ where we have introduced auxiliary fields $`\varphi `$$`B^{(a)}`$ and $`\gamma ^{(a)}`$ for the constraints $`_ak^{(a)}=0`$$`_{}^{}{}_{\mu }{}^{}k_\mu ^{(a)}(s)=0`$ and the definition of $`K^{(a)}`$, respectively. Here we have used such notations as $`𝒟𝜸_a𝒟\gamma ^{(a)}`$. There are following identities for integers $`m^{(a)}(s),n_\mu ^{(a)}(s^{(n)})`$ and $`P_\mu (s)`$: $`\mathrm{exp}\{2\pi im^{(a)}(s)_{\mu }^{}{}_{}{}^{}k_\mu ^{(a)}(s)\}=1,`$ (33) $`\mathrm{exp}\{2\pi in_\mu ^{(a)}(s^{(n)})(K_\mu ^{(a)}(s^{(n)}){\displaystyle \underset{i,j,k=0}{\overset{n1}{}}}k_\mu ^{(a)})\}=1,`$ (34) $`\mathrm{exp}\{2\pi iP_\mu (s){\displaystyle \underset{a}{}}k_\mu ^{(a)}(s)\}=1.`$ (35) Then we can change the integral region of $`\gamma ^{(a)},B^{(a)}`$ and $`\varphi `$ from the first Brillouin zone to the infinite region. Using the Poisson summation formula $`{\displaystyle \underset{k_\mu ^{(a)}Z}{}}f[k_\mu ^{(a)}]`$ $`=`$ $`\mathrm{const}.{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑F_\mu ^{(a)}{\displaystyle \underset{l_\mu ^{(a)}Z}{}}\mathrm{exp}\{2\pi iF_\mu ^{(a)}l_\mu ^{(a)}\}f[F_\mu ^{(a)}],`$ (36) $`Z[K^{(a)},J]`$ becomes $`Z[K^{(a)},J]`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝜸{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\varphi {\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝑩{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝑭{\displaystyle \underset{l^{(a)}Z}{}}`$ (37) $`\times \mathrm{exp}\{{\displaystyle \underset{a,s,s^{}}{}}F_\mu ^{(a)}(s)D_0(ss^{})F_\mu ^{(a)}(s^{})+2\pi i{\displaystyle \underset{s}{}}N_\mu (s)F_\mu ^{(1)}(s)`$ $`+i{\displaystyle \underset{a,s}{}}(_\mu B^{(a)}(s)+\varphi _\mu (s)+2\pi l_\mu ^{(a)}(s))F_\mu ^{(a)}(s)`$ $`i{\displaystyle \underset{a,s^{(n)}}{}}\gamma _\mu ^{(a)}(s^{(n)}){\displaystyle \underset{i,j,k=0}{\overset{n1}{}}}F_\mu ^{(a)}\}.`$ Writing the lattice spacing explicitly, we get $`ia^4n{\displaystyle \underset{a,s^{(n)}}{}}\gamma _\mu ^{(a)}(nas^{(n)}){\displaystyle \underset{i,j,k=0}{\overset{n1}{}}}F_\mu ^{(a)}(nas^{(n)}+(n1)a\widehat{\mu }+ia\widehat{\nu }+ja\widehat{\rho }+ka\widehat{\sigma })`$ $`=ia^4{\displaystyle \underset{a,s}{}}X_\mu ^{(a)}(as)F_\mu ^{(a)}(as),`$ (38) $`X_\mu ^{(a)}(as)`$ $``$ $`na^4{\displaystyle \underset{s^{(n)}}{}}\gamma _\mu ^{(a)}(nas^{(n)})`$ (39) $`\times \delta (nas_\mu ^{(n)}+(n1)aas_\mu ){\displaystyle \underset{i(\mu )}{}}{\displaystyle \underset{I=0}{\overset{n1}{}}}\delta (nas_i^{(n)}+Iaas_i).`$ Integrating out $`F^{(a)}`$ and $`B^{(a)}`$, we find $`Z[K^{(a)},J]`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝜸\mathrm{exp}\{ib^4{\displaystyle \underset{a,s^{(n)}}{}}\gamma _\mu ^{(a)}(bs^{(n)})K_\mu ^{(a)}(bs^{(n)})\}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\varphi `$ (40) $`\times \mathrm{exp}\{{\displaystyle \frac{1}{4}}a^8{\displaystyle \underset{s,s^{}}{}}V_\mu ^{(a)}(as)A_{\mu \nu }(asas^{})V_\nu ^{(a)}(as^{})\}`$ $`\times {\displaystyle \underset{al^{(a)}Z}{}}\mathrm{exp}\{\pi ^2a^8{\displaystyle \underset{a,s,s^{}}{}}l_\mu ^{(a)}(as)A_{\mu \nu }(asas^{})l_\nu ^{(a)}(as^{})`$ $`+\pi ^2a^8{\displaystyle \underset{a,s,s^{}}{}}V_\mu ^{(a)}(as)A_{\mu \nu }(asas^{})l_\nu ^{(a)}(as)\},`$ where $`V_\mu ^{(a)}(as)X_\mu ^{(a)}(as)\varphi _\mu (as)2\pi N_\mu (as)\delta _{a,1},`$ (41) $`A_{\mu \nu }(asas^{})\left\{\delta _{\mu \nu }{\displaystyle \frac{_\mu _\nu ^{}}{_\rho _\rho _\rho ^{}}}\right\}D_0^1(asas^{}).`$ (42) Here we consider $`l^{(a)}=0`$ alone, since we finally take the $`a0`$ limit in the sum with respect to $`al^{(a)}Z`$. We change the variables from $`\gamma ^{(a)}`$ to $`\gamma _{}^{}{}_{}{}^{(a)}`$: $`\gamma _{}^{}{}_{}{}^{(1)}`$ $``$ $`\gamma ^{(2)}\gamma ^{(3)},\gamma _{}^{}{}_{}{}^{(2)}\gamma ^{(3)}\gamma ^{(1)},\gamma _{}^{}{}_{}{}^{(3)}\gamma ^{(1)}\gamma ^{(2)}.`$ (43) Integrating out $`\gamma _{}^{}{}_{}{}^{(a)}`$ and $`\varphi _\mu `$, we get the monopole action in terms of $`K^{(a)}`$ as follows: $`Z[J]`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{2}{3}}\pi ^2a^8{\displaystyle \underset{s,s^{}}{}}N_\mu (as)A_{\mu \nu }(asas^{})N_\nu (as^{})`$ (44) $`+{\displaystyle \frac{2}{3}}\pi ^2b^8{\displaystyle \underset{s^{(n)},s^{(n)}}{}}B_\mu (bs^{(n)})A_{}^{}{}_{\mu \nu }{}^{1}(bs^{(n)}bs^{(n)})B_\nu (bs^{(n)})\}Z_{mon}[J],`$ where $`B_\mu (bs^{(n)})\underset{\genfrac{}{}{0pt}{}{a0}{n\mathrm{}}}{lim}{\displaystyle \frac{1}{n^3}}a^8{\displaystyle \underset{s,s^{},\nu }{}}\delta (nas_\mu ^{(n)}+(n1)aas_\mu ){\displaystyle \underset{i(\mu )}{}}{\displaystyle \underset{I=0}{\overset{n1}{}}}\delta (nas_i^{(n)}+Iaas_i)`$ $`\times \left\{\delta _{\mu \nu }{\displaystyle \frac{_\mu _\nu ^{}}{_\rho _\rho _\rho ^{}}}\right\}D_0^1(asas^{})N_\nu (as^{}),`$ (45) $`A_{}^{}{}_{\mu \nu }{}^{}(bs^{(n)}bs^{(n)})`$ $`\underset{\genfrac{}{}{0pt}{}{a0}{n\mathrm{}}}{lim}{\displaystyle \frac{1}{n^6}}a^8{\displaystyle \underset{s,s^{}}{}}\delta (nas_\mu ^{(n)}+(n1)aas_\mu ){\displaystyle \underset{i(\mu )}{}}{\displaystyle \underset{I=0}{\overset{n1}{}}}\delta (nas_i^{(n)}+Iaas_i)`$ $`\times \left\{\delta _{\mu \nu }{\displaystyle \frac{_\mu _\nu ^{}}{_\rho _\rho _\rho ^{}}}\right\}D_0^1(asas^{})`$ $`\times \delta (nas_{}^{}{}_{\nu }{}^{(n)}+(n1)aas_\nu ^{}){\displaystyle \underset{j(\nu )}{}}{\displaystyle \underset{J=0}{\overset{n1}{}}}\delta (nas_{}^{}{}_{j}{}^{(n)}+Jaas_j^{}).`$ (46) $`Z_{mon}[J]`$ is the dynamical monopole part: $`Z_{mon}[J]`$ $`=`$ $`{\displaystyle \underset{K^{(a)}Z}{}}\delta _{_\mu ^{}K_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_aK^{(a)},0}`$ (47) $`\times \mathrm{exp}\{b^8{\displaystyle \underset{a,s^{(n)},s^{(n)}}{}}K_\mu ^{(a)}(bs^{(n)})A_{\mu \nu }^1(bs^{(n)}bs^{(n)})K_\nu ^{(a)}(bs^{(n)})`$ $`+2\pi ib^8{\displaystyle \underset{s^{(n)},s^{(n)}}{}}B_\mu (bs^{(n)})A_{\mu \nu }^1(bs^{(n)}bs^{(n)})K_\nu ^{(1)}(bs^{(n)})\}.`$ The interactions of the perfect action on the $`b`$-lattice in Eq. (4.20) depend on directions. This is consistent with the numerical data. The spectrum of $`K_\mu ^{(a)}(bs^{(n)})`$ is found to be equivalent to that in the continuum theory as discussed in the SU(2) case . ## 5 String representation and rotational invariance When we transform $`Z_{mon}[J]`$ in Eq. (4.20) into the string representation as in SU(2) , we can estimate the static potential analytically. First we change the variables as follows: $`K^{(1)}`$ $`=`$ $`j^{(2)}j^{(3)},K^{(2)}=j^{(3)}j^{(1)},K^{(3)}=j^{(1)}j^{(2)}.`$ (48) The conservation laws of monopole currents $`_\mu ^{}K_\mu ^{(a)}=0`$ leads us to $`_\mu ^{}j_\mu ^{(1)}=_\mu ^{}j_\mu ^{(2)}=_\mu ^{}j_\mu ^{(3)}`$. The conditions are expressed as $`{\displaystyle _\pi ^\pi }𝒟𝝋\delta ({\displaystyle \underset{a}{}}\phi ^{(a)})\mathrm{exp}\{i{\displaystyle \underset{a,s}{}}j_\mu ^{(a)}(s)_\mu \phi ^{(a)}(s)\}.`$ (49) Then $`Z_{mon}[J]`$ is reduced to the following: $`Z_{mon}[J]`$ $`=`$ $`{\displaystyle _\pi ^\pi }𝒟𝝋\delta ({\displaystyle \underset{a}{}}\phi ^{(a)}){\displaystyle \underset{j^{(a)}Z}{}}\mathrm{exp}\{2{\displaystyle \underset{a,s,s^{}}{}}j_\mu ^{(a)}(s)A_{}^{}{}_{\mu \nu }{}^{1}(ss^{})j_\nu ^{(a)}(s^{})`$ (51) $`+2{\displaystyle \underset{a<b,s,s^{}}{}}j_\mu ^{(a)}(s)A_{}^{}{}_{\mu \nu }{}^{1}(ss^{})j_\nu ^{(b)}(s^{})+i{\displaystyle \underset{a,s}{}}j_\mu ^{(a)}(s)_\mu \phi ^{(a)}(s)`$ $`+2\pi i{\displaystyle \underset{s,s^{}}{}}B_\mu (s)A_{}^{}{}_{\mu \nu }{}^{1}(ss^{})(j_\nu ^{(2)}(s^{})j_\mu ^{(3)}(s^{}))\}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\stackrel{}{C}{\displaystyle _\pi ^\pi }𝒟𝝋\delta ({\displaystyle \underset{a}{}}\phi ^{(a)}){\displaystyle \underset{j^{(a)}Z}{}}\mathrm{exp}\{{\displaystyle \frac{1}{8}}{\displaystyle \underset{s,s^{}}{}}\stackrel{}{C}_\mu (s)A_{}^{}{}_{\mu \nu }{}^{}(ss^{})\stackrel{}{C}_\nu (s^{})`$ $`+i{\displaystyle \underset{a,s}{}}j_\mu ^{(a)}(s)(_\mu \phi ^{(a)}(s)+\stackrel{}{ϵ}_a\stackrel{}{C}_\mu (s))`$ $`+2\pi i{\displaystyle \underset{s,s^{}}{}}B_\mu (s)A_{}^{}{}_{\mu \nu }{}^{1}(ss^{})(j_\nu ^{(2)}(s^{})j_\nu ^{(3)}(s^{}))\},`$ where auxiliary fields $`\stackrel{}{C}_\mu (C_\mu ^3,C_\mu ^8)`$ are introduced and $`\stackrel{}{ϵ}_a`$ are the SU(3) root vectors: $`\stackrel{}{ϵ}_1=(1,0)`$$`\stackrel{}{ϵ}_2=(1/2,\sqrt{3}/2)`$$`\stackrel{}{ϵ}_3=(1/2,\sqrt{3}/2)`$. Since $`_a\phi ^{(a)}=0`$ and $`_a\stackrel{}{ϵ}_a=0`$, we get an identity $`\mathrm{exp}\{{\displaystyle \frac{i}{2\pi }}{\displaystyle \underset{s}{}}\varphi _\mu (s){\displaystyle \underset{a}{}}(_\mu \phi ^{(a)}(s)+\stackrel{}{ϵ}_a\stackrel{}{C}_\mu (s))\}=1.`$ (52) The Poisson summation formula $`{\displaystyle \underset{l_\mu ^{(a)}Z}{}}\mathrm{exp}\{2\pi i{\displaystyle \underset{a}{}}F_\mu ^{(a)}l_\mu ^{(a)}+i\varphi _\mu {\displaystyle \underset{a}{}}l_\mu ^{(a)}\}`$ $`=`$ $`{\displaystyle \underset{j_\mu ^{(a)}Z}{}}{\displaystyle \underset{a}{}}\delta (F_\mu ^{(a)}+{\displaystyle \frac{\varphi _\mu }{2\pi }}j_\mu ^{(a)})`$ (53) gives us $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝑭{\displaystyle \underset{l^{(a)}Z}{}}\delta _{\mathrm{\Sigma }_al^{(a)},0}\mathrm{exp}\{2\pi i{\displaystyle \underset{a,s}{}}F_\mu ^{(a)}(s)l_\mu ^{(a)}(s)\}f(F_\mu ^{(a)})=`$ $`\mathrm{const}.{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\varphi {\displaystyle \underset{j^{(a)}Z}{}}f(j_\mu ^{(a)}{\displaystyle \frac{\varphi _\mu }{2\pi }}).`$ (54) Then $`Z_{mon}[J]`$ becomes $`Z_{mon}[J]`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝑭{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\stackrel{}{C}{\displaystyle _\pi ^\pi }𝒟𝝋\delta ({\displaystyle \underset{a}{}}\phi ^{(a)}){\displaystyle \underset{l^{(a)}Z}{}}\delta _{\mathrm{\Sigma }_al^{(a)},0}`$ (55) $`\times \mathrm{exp}\{{\displaystyle \frac{1}{8}}{\displaystyle \underset{s,s^{}}{}}\stackrel{}{C}_\mu (s)A_{}^{}{}_{\mu \nu }{}^{}(ss^{})\stackrel{}{C}_\nu (s^{})`$ $`+i{\displaystyle \underset{a,s}{}}F_\mu ^{(a)}(s)(_\mu \phi ^{(a)}(s)+\stackrel{}{ϵ}_a\stackrel{}{C}_\mu (s)+2\pi l_\mu ^{(a)}(s))`$ $`+2\pi {\displaystyle \underset{s^{}}{}}(F_\mu ^{(2)}(s)F_\mu ^{(3)}(s))A_{}^{}{}_{\mu \nu }{}^{1}(ss^{})B_\nu (s^{})\}.`$ Here we perform the Berezinski-Kosterlitz-Thouless (BKT) transformation : $`l_\nu ^{(a)}(s)=s_\mu ^{(a)}(s)+_\mu r^{(a)}(s),_{[\mu }s_{\nu ]}^{(a)}(s)=\sigma _{\mu \nu }^{(a)}(s),`$ $`_\mu \phi ^{(a)}(s)+2\pi l_\mu ^{(a)}(s)=_\mu \phi _{nc}^{(a)}(s)2\pi {\displaystyle \underset{s^{}}{}}_\nu ^{}\mathrm{\Delta }^1(ss^{})\sigma _{\nu \mu }^{(a)}(s^{}),`$ (56) where $`\phi _{nc}^{(a)}(s)\phi ^{(a)}(s)2\pi _s^{}\mathrm{\Delta }^1(ss^{})_\nu ^{}s_\nu ^{(a)}(s^{})+2\pi r^{(a)}(s)`$ is non-compact. The plaquette variable $`\sigma _{\mu \nu }^{(a)}`$ satisfies a conservation law $`{}_{[\alpha }{}^{}\sigma _{\mu \nu ]}^{(a)}(s)=0`$ and a constraint $`_a\sigma _{\mu \nu }^{(a)}=0`$ due to $`_al_\mu ^{(a)}=0`$. Using the condition $`_\mu ^{}B_\mu =0`$, we integrate out $`F_\mu ^{(a)}`$$`\stackrel{}{C}_\mu `$ and $`\phi _{nc}^{(a)}`$. We get the string representation: $`Z_{str}[J]`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{2}{3}}\pi ^2{\displaystyle \underset{s,s^{}}{}}B_\mu (s)A_{}^{}{}_{\mu \nu }{}^{1}(ss^{})B_\nu (s^{})\}{\displaystyle \underset{\sigma ^{(a)}Z}{}}\delta _{_{[\alpha }\sigma _{\mu \nu ]}^{(a)},0}\delta _{\mathrm{\Sigma }_a\sigma ^{(a)},0}`$ (57) $`\times \mathrm{exp}\{{\displaystyle \frac{\pi ^2}{3}}{\displaystyle \underset{a,s,s^{}}{}}_\alpha ^{}\sigma _{\mu \alpha }^{(a)}(s)H_{\mu \nu }(ss^{})_\beta ^{}\sigma _{\nu \beta }^{(a)}(s^{})`$ $`+{\displaystyle \frac{2\pi ^2}{3}}{\displaystyle \underset{s,s^{}}{}}(\sigma _{\mu \nu }^{(2)}(s)\sigma _{\mu \nu }^{(3)}(s))_\nu \mathrm{\Delta }^1(ss^{})B_\mu (s^{})\},`$ $$H_{\mu \nu }(ss^{})=\underset{s_1}{}A_{}^{}{}_{\mu \nu }{}^{}(ss_1)\mathrm{\Delta }^2(s_1s^{}).$$ (58) Since the numerical results show that the couplings of the monopole action are weak for large $`b`$ region, we can use the strong coupling expansion in the dual string model. Then we see that the quantum fluctuations can be neglected as in SU(2) case . The classical part in Eq. (5.10) cancels the second classical term in Eq. (4.17). As a result, the classical part of the expectation value of the Wilson loop is reduced to $`W(C)_{cl}`$ $`=`$ $`\mathrm{exp}\{{\displaystyle \frac{2}{3}}\pi ^2{\displaystyle _{\mathrm{}}^{\mathrm{}}}d^4xd^4yN_\mu (x)A_{\mu \nu }(xy)N_\nu (y)\}.`$ (59) Since the classical part is written in the continuum form, the continuum rotational invariance is trivial. The static potential $`V(Ib,0,0)`$ and $`V(Ib,Ib,0)`$ can be evaluated as in SU(2). We take the following plaquette sources as $`S_{\alpha \beta }(x)`$ for $`V(Ib,0,0)`$: $`S_{\alpha <\beta }(x)`$ $`=`$ $`\delta _{\alpha 1}\delta _{\beta 4}\theta (x_1)\theta (Ibx_1)\theta (x_4)\theta (Tbx_4)\delta (x_2)\delta (x_3)`$ (60) and for $`V(Ib,Ib,0)`$: $`S_{\alpha <\beta }(x)`$ $`=`$ $`(\delta _{\alpha 1}\delta _{\beta 4}+\delta _{\alpha 2}\delta _{\beta 4})\delta (x_3)\theta (x_4)\theta (Tbx_4)`$ (61) $`\times \theta (x_1)\theta (Ibx_1)\theta (x_2)\theta (Ibx_2)\delta (x_1x_2),`$ respectively. Using the following formula $`\underset{T\mathrm{}}{lim}\left({\displaystyle \frac{\mathrm{sin}\alpha T}{\alpha }}\right)^2`$ $`=`$ $`\pi T\delta (\alpha ),`$ (62) we get the static potentials $`V(Ib,0,0)`$ $`=`$ $`\underset{T\mathrm{}}{lim}{\displaystyle \frac{1}{Tb}}\mathrm{ln}W(Ib,0,0,Tb)_{cl}`$ (63) $`\underset{I\mathrm{}}{}`$ $`{\displaystyle \frac{2\pi ^2}{3}}(Ib)^2{\displaystyle \frac{d^2p}{2\pi ^2}\left[\frac{1}{\mathrm{\Delta }D_0}\right](0,p_2,p_3,0)}`$ $`=`$ $`{\displaystyle \frac{\pi \kappa Ib}{3}}\mathrm{ln}{\displaystyle \frac{m_1}{m_2}},`$ $`V(Ib,Ib,0)`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}\pi \kappa Ib}{3}}\mathrm{ln}{\displaystyle \frac{m_1}{m_2}},`$ (64) where $`\kappa (m_1^2m_2^2)=\gamma ^1,m_1^2+m_2^2=\beta /\gamma ,m_1^2m_2^2=\alpha /\gamma `$. The static potential has the linear term alone and the rotational invariance is recovered completely. The string tension is evaluated as $$\sigma _{cl}=\frac{\pi \kappa }{3}\mathrm{ln}\frac{m_1}{m_2}.$$ (65) This is consistent with the results . The $`m_1^1`$ and the $`m_2^1`$ could be regarded as the coherence and the penetration lengths in Type-2 superconductor. ## 6 Estimate of the string tension Ths string tension in Eq. (5.18) is written by the parameters of the monopole action taken on the $`a`$-lattice. The parameters can be determined by comparing the theoretical perfect action in Eq. (4.20) with the numerical results of the monopole action on the $`b`$-lattice. $`A_{}^{}{}_{\mu \nu }{}^{1}`$ in Eq. (4.19) fixed the gauge is calculated in ref. . The results of the comparison are seen in Figure 3. The agreement is rather good. The value of the string tension is $`\sqrt{\sigma _{cl}/\sigma _{ph}}1.22`$ for large $`b`$ as in Figure 4. The result is not so bad, although the discrepancy is not negligible. In order to study whether the difference is due to the ambiguity of the fit of $`D_0`$, we try to estimate the string tension without determining the parameters on the $`a`$-lattice. The following monopole Wilson loop on the $`b`$-lattice is considered: $`W(C)`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle \underset{K^{(a)}Z}{}}\delta _{_\mu ^{}K_\mu ^{(a)},0}\delta _{\mathrm{\Sigma }_aK^{(a)},0}\mathrm{exp}\{{\displaystyle \underset{a,s,s^{}}{}}K_\mu ^{(a)}(s)D_{\mu \nu }^1(ss^{})K_\nu ^{(a)}(s^{})`$ (66) $`+2\pi i{\displaystyle \underset{s}{}}N_\mu (s)K_\mu ^1(s)\},`$ where the naive monopole Wilson loop operator (A.15) on the coarse $`b`$-lattice is used and $`D_{\mu \nu }^1`$ corresponds to $`A_{}^{}{}_{\mu \nu }{}^{1}`$. Transforming Eq. (6.1) to the string representation and neglecting the quantum fluctuations as in the previous section, we get the classical part of the Wilson loop as follows: $$W(C)_{cl}=\mathrm{exp}\{\frac{2}{3}\pi ^2\underset{s,s^{}}{}N_\mu (s)D_{\mu \nu }(ss^{})N_\nu (s^{})\}.$$ (67) It is found theoretically that the static potential $`V(Ib,0,0)`$ evaluated by Eq. (6.2) agrees with Eq. (5.16) in the $`I\mathrm{}`$ case. Using Eq. (6.2), the string tension becomes $`\sigma _{cl}`$ $`=`$ $`{\displaystyle \frac{\pi ^2}{6}}{\displaystyle _\pi ^\pi }{\displaystyle \frac{d^2k}{(2\pi )^2}}{\displaystyle \frac{1}{(\mathrm{sin}^2\frac{k_2}{2}+\mathrm{sin}^2\frac{k_3}{2})^2}}`$ (68) $`\times [\mathrm{sin}^2{\displaystyle \frac{k_3}{2}}D(0,k_2,k_3,0;\widehat{2})+\mathrm{sin}^2{\displaystyle \frac{k_2}{2}}D(0,k_2,k_3,0;\widehat{3})],`$ $$D_{\mu \nu }(ss^{})=\delta _{\mu \nu }_\pi ^\pi \frac{dk^4}{(2\pi )^4}D(k_1,k_2,k_3,k_4;\widehat{\mu })e^{ik(ss^{})}.$$ (69) We estimate the string tension from Eq. (6.3) with the numerical data of the monopole action on the $`b`$-lattice. The results are seen in Figure 4. The results are not so different from those of the first method for large $`b`$. The $`b`$-independence holds roughly for $`b>2`$ and $`\sqrt{\sigma _{cl}/\sigma _{ph}}1.27`$ for large $`b`$. The difference from the physical value is also seen in this method. Using the same method as in SU(2) , the effect of the quantum fluctuation can be evaluated. The order of the correction for the string tension becomes $$\frac{8}{b^2}\mathrm{exp}\{5\mathrm{\Pi }(0)b^2\},$$ (70) where $`\mathrm{\Pi }(0)`$ represents the coupling of the self-interaction in the string model. For example, the value of Eq. (6.5) is $`2.6\times 10^{10}`$ for $`b=3.27`$. The correction is very small in large $`b`$ region. In order to clarify the origin of the difference, we have to study the effects of systematic errors carefully. As a systematic error, the difference between D-T monopoles used in the simulation and the monopoles in the analytical calculations may be important. The magnetic charge of the latter ranges from $`\mathrm{}`$ to $`\mathrm{}`$, while D-T monopoles take only the restricted values of magnetic charges. This is under investigation. The glueball spectrum can be evaluated similarly as in SU(2) QCD . The lightest $`0^{++}`$ glueball mass is determined as $`M_{0^{++}}=2m_2`$ . For $`b=4.23`$$`2m_2/\sqrt{\sigma _{ph}}=3.90`$. The result is not so different from $`3.64\pm 0.09`$ . ## 7 Conclusions We have studied an effective monopole action in pure SU(3) lattice QCD. The following results are obtained: 1. We have found that the SU(3) monopole action can be derived clearly for various $`n`$-blocked monopoles in MA gauge numerically, using the extended Swendsen method in the two-current and the one-current cases. We perform $`n=28`$ step block-spin transformations on the $`48^4`$ lattice. In the two-current case, we obtain an almost perfect lattice monopole action for the infrared region of SU(3) QCD, since the action seems to depend on a physical scale $`b`$ alone. The simple 4, 6-point interactions become negligibly small as compared with two-point interactions for large $`b`$ region as in SU(2). Thus two-point monopole interactions are dominant in the infrared region. In the one-current case, if the physical scale $`b`$ is taken in unit of the string tension $`\sqrt{\sigma _{ph}}`$, the scaling which was not seen in the previous study holds beautifully. Monopole condensation occurs in the large $`b`$ region from the energy-entropy balance. 2. The perfect action satisfies the continuum rotational invariance. When we try to study the restoration of the rotational invariance of the static potential, we have to know the correct form of the perfect operator. For that purpose, a block-spin transformation from the small $`a`$-lattice ($`a0`$) to the finite $`b=na(\beta )`$ lattice is performed analytically in a simple restricted case of a quadratic monopole action as in SU(2) case. The perfect monopole action and the perfect operator evaluating the static potential between electric charges on the $`b`$-lattice are derived. The quadratic SU(3) monopole action is simple, but it corresponds to the London limit of the DGL theory which is non-trivially interacting field theory. 3. The SU(3) monopole action can be transformed into the string model. There are three strings $`\sigma ^{(a)}`$ satisfying one constraint $`_a\sigma ^{(a)}=0`$. This is consistent with the results in the continuum limit . Since the monopole interactions are weak for large $`b`$ region, the dual string interactions are strong. Using the strong coupling expansion, the static potential in the long distance is calculated by the classical part alone analytically. The static potential has a linear term alone and the restoration of the rotational invariance can be seen explicitly. The string tension is estimated from the numerical data of the monopole action. The result is rather good, but the difference from the physical string tension is not negligible. The same thing happens in SU(2) case . It seems that there still exist some systematic errors. ## Acknowledgments This work is supported by the Supercomputer Project (No.98-33 and No.99-47) of High Energy Accelerator Research Organization (KEK) and the Supercomputer Project of the Institute of Physical and Chemical Research (RIKEN). T.S. is financially supported by JSPS Grant-in-Aid for Scientific Research (B) (No. 10440073 and No. 11695029). ## Appendix A Appendix In SU(2), it is already known that the monopole action can be derived from the dual abelian higgs model on the fine $`a`$-lattice . We show how to derive the monopole action from the DGL theory on the lattice in SU(3). In DGL, $`\stackrel{}{A}_\mu (A_\mu ^3,A_\mu ^8)`$ are used as $`U(1)^2`$ photon fields, where $`A_\mu =A_\mu ^i\lambda ^i/2`$ ($`\lambda ^i`$ are Gell-Mann matrices). The magnetic charges $`\stackrel{}{m}`$ are distributed on the SU(3) root lattice: $`\stackrel{}{m}=g_m_{a=1}^3\xi _a\stackrel{}{ϵ}_a`$$`\xi _aZ`$$`g_m=4\pi /g`$ ($`g`$ is the SU(3) coupling constant). The DGL Lagrangian in the continuum limit is $$_{DGL}=\frac{1}{4}(_\mu \stackrel{}{C}_\mu _\nu \stackrel{}{C}_\mu )^2+\underset{a=1}{\overset{3}{}}(\|(_\mu +ig_m\stackrel{}{ϵ}_a\stackrel{}{C}_\mu )\varphi ^{(a)}\|^2+\lambda (\|\varphi ^{(a)}\|^2v^2)^2),$$ (A.1) where $`\stackrel{}{C}_\mu (C_\mu ^3,C_\mu ^8)`$ are dual $`U(1)^2`$ gauge fields and $`\varphi ^{(a)}=\rho ^{(a)}\mathrm{exp}(i\phi ^{(a)})`$ are monopole fields with a constraint $`_a\phi ^{(a)}=0`$. In the below, we use the notations of differential forms on the lattice . The lattice DGL action becomes $`S_{DGL}[\stackrel{}{C},\varphi ^{(a)}]`$ $`=`$ $`{\displaystyle \frac{1}{2g_m^2}}\|\|d\stackrel{}{C}\|\|^2\gamma {\displaystyle \underset{a,s,\mu }{}}(\varphi _s^{(a)}U_\mu ^{(a)}(s)\varphi _{s+\widehat{\mu }}^{(a)}+h.c)`$ (A.2) $`+\lambda {\displaystyle \underset{a,s}{}}(\varphi _s^{(a)}\varphi _s^{(a)}1)^2+{\displaystyle \underset{a,s}{}}\varphi _s^{(a)}\varphi _s^{(a)},`$ where $`U^{(a)}=\mathrm{exp}(i\stackrel{}{ϵ}_a\stackrel{}{C})`$. We modify the partition function of the DGL model using the Villain approximation $`e^{\alpha \mathrm{cos}\psi }_{lZ}e^{\frac{\alpha }{2}(\psi +2\pi l)^2}`$. In DGL, the summations of $`l^{(a)}`$ appear due to 2$`\pi `$ periodicity of $`\mathrm{cos}(d\phi ^{(a)}+\stackrel{}{ϵ}_a\stackrel{}{C})`$. Then there is a constraint $`_al^{(a)}=0`$, since $`_a(d\phi ^{(a)}+\stackrel{}{ϵ}_a\stackrel{}{C})=0`$. The partition function of the DGL theory is given by $`Z_{DGL}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\stackrel{}{C}{\displaystyle _\pi ^\pi }𝒟𝝋\delta ({\displaystyle \underset{a}{}}\phi ^{(a)}){\displaystyle _0^{\mathrm{}}}𝒟𝝆^2{\displaystyle \underset{l^{(a)}Z}{}}\delta _{\mathrm{\Sigma }_al^{(a)},0}`$ (A.3) $`\times \mathrm{exp}\{{\displaystyle \frac{1}{2g_m^2}}\|\|d\stackrel{}{C}\|\|^2\gamma {\displaystyle \underset{a,s,\mu }{}}\rho _s^{(a)}\rho _{s+\widehat{\mu }}^{(a)}(d\phi ^{(a)}+\stackrel{}{ϵ}_a\stackrel{}{C}+2\pi l^{(a)})_{s,\mu }^2`$ $`\lambda {\displaystyle \underset{a,s}{}}((\rho _s^{(a)})^21)^2{\displaystyle \underset{a,s}{}}(\rho _s^{(a)})^2\}.`$ Inserting $`1`$ $`=`$ $`\left\{{\displaystyle \underset{a,s,\mu }{}}(4\gamma \rho _s^{(a)}\rho _{s+\widehat{\mu }}^{(a)})^{1/2}\right\}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟𝑭\mathrm{exp}\{{\displaystyle \underset{a,s,\mu }{}}{\displaystyle \frac{1}{4\gamma \rho _s^{(a)}\rho _{s+\widehat{\mu }}^{(a)}}}`$ (A.4) $`\times (F_\mu ^{(a)}(s)2i\gamma \rho _s^{(a)}\rho _{s+\widehat{\mu }}^{(a)}(d\phi ^{(a)}+\stackrel{}{ϵ_a}\stackrel{}{C}+2\pi l^{(a)})_{s,\mu })^2\}`$ and using the relation in Eq. (5.7), we integrate out $`\varphi ,\phi ^{(a)}`$ and $`\stackrel{}{C}`$. The partition function is written by the variables $`j^{(a)}`$ in the r.h.s. of Eq. (5.7) and $`\rho ^{(a)}`$. From the integrals of $`\phi ^{(a)}`$, the relations $`\delta j^{(1)}=\delta j^{(2)}=\delta j^{(3)}`$ appear. Here monopole currents are introduced as follows: $$k^{(1)}j^{(2)}j^{(3)},k^{(2)}j^{(3)}j^{(1)},k^{(3)}j^{(1)}j^{(2)}.$$ (A.5) They satisfy the conservation law $`\delta k^{(a)}=0`$ and the constraint $`_ak^{(a)}=0`$. The partition function can be rewritten as follows: $$Z_{mon}[k^{(a)}]=\underset{k^{(a)}Z}{}\delta _{\delta k^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}\mathrm{exp}\{S_{mon}^G[k^{(a)}]S_{mon}^H[k^{(a)}]\},$$ (A.6) $$S_{mon}^G[k^{(a)}]=\frac{g_m^2}{4}\underset{a=1}{\overset{3}{}}(k^{(a)},\mathrm{\Delta }^1k^{(a)}),$$ (A.7) $`S_{mon}^H[k^{(a)}]`$ $`=`$ $`\mathrm{ln}\{{\displaystyle _0^{\mathrm{}}}𝒟𝝆^2\mathrm{exp}\{{\displaystyle \frac{1}{4\gamma }}{\displaystyle \underset{s,\mu }{}}{\displaystyle \frac{_a\rho _s^{(a)}\rho _{s+\widehat{\mu }}^{(a)}(k_\mu ^{(a)}(s))^2}{_{a<b}\rho _s^{(a)}\rho _{s+\widehat{\mu }}^{(b)}}}`$ (A.8) $`\lambda {\displaystyle \underset{a,s}{}}((\rho _s^{(a)})^21)^2{\displaystyle \underset{a,s}{}}(\rho _s^{(a)})^2\}\}.`$ The integration of $`\rho ^{(a)}`$ in Eq. (A.8) can’t be performed exactly. In the London limit ($`\lambda \mathrm{}`$), $`\rho ^{(a)}`$ are fixed to unity. $`S_{mon}^H[k^{(a)}]`$ becomes the quadratic self interaction: $$S_{mon}^H[k^{(a)}]=\frac{1}{12\gamma }\underset{a}{}k^{(a)}^2.$$ (A.9) The monopole action in Eq. (4.1) $`S_{mon}[k^{(a)}]`$ $`=`$ $`{\displaystyle \underset{a}{}}(k^{(a)},D_0k^{(a)}),D_0=\alpha \mathrm{\Delta }^1+\beta +\gamma \mathrm{\Delta }`$ (A.10) corresponds to the modified London limit of DGL: $`S_{DGL}[\stackrel{}{C},\phi ^{(a)}]`$ $`=`$ $`{\displaystyle \frac{1}{8\alpha }}d\stackrel{}{C}^2+{\displaystyle \frac{1}{12}}{\displaystyle \underset{a}{}}(X^{(a)},\{\beta +\gamma \mathrm{\Delta }\}^1X^{(a)}),`$ (A.11) $`X^{(a)}=d\phi ^{(a)}+\stackrel{}{ϵ}_a\stackrel{}{C}+2\pi l^{(a)}.`$ The Wilson loop is also calculated similarly. For simplicity, we consider the case of the London limit. Since the DGL theory is the dual theory for a color electric $`U(1)^2`$ charge, we consider the ’t Hooft loop: $`H_c(C)={\displaystyle \frac{1}{Z_{DGL}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝒟\stackrel{}{C}{\displaystyle _\pi ^\pi }𝒟𝝋\delta ({\displaystyle \underset{a}{}}\phi ^{(a)}){\displaystyle \underset{l^{(a)}Z}{}}\delta _{\mathrm{\Sigma }_al^{(a)},0}`$ (A.12) $`\times \mathrm{exp}\{{\displaystyle \frac{1}{2g_m^2}}d\stackrel{}{C}4\pi \stackrel{}{Q}_c{}_{}{}^{}S^2\gamma {\displaystyle \underset{a}{}}d\phi ^{(a)}+\stackrel{}{ϵ}_a\stackrel{}{C}+2\pi l^{(a)}^2\},`$ where $`S`$ is a surface with the fixed boundary $`C`$$`\delta S=j`$, where $`j`$ is the unit current on the loop $`C`$$`\stackrel{}{Q}_c`$ are color electric $`U(1)^2`$ charges of quarks with colors c (c=R,G,B): $$\stackrel{}{Q}_c=(Q_c^3,Q_c^8)=\{(\frac{1}{2},\frac{1}{2\sqrt{3}}),(\frac{1}{2},\frac{1}{2\sqrt{3}}),(0,\frac{1}{\sqrt{3}})\}.$$ (A.13) For c=R, we get $`H_R(C)={\displaystyle \frac{1}{Z_{mon}}}{\displaystyle \underset{k^{(a)}Z}{}}\delta _{\delta k^{(a)},0}\delta _{\mathrm{\Sigma }_ak^{(a)},0}\mathrm{exp}\{{\displaystyle \frac{g_m^2}{4}}{\displaystyle \underset{a}{}}(k^{(a)},\mathrm{\Delta }^1k^{(a)})`$ (A.14) $`{\displaystyle \frac{1}{12\gamma }}{\displaystyle \underset{a}{}}\|\|k^{(a)}\|\|^2+2\pi i({}_{}{}^{}S,\mathrm{\Delta }^1dk^{(3)}){\displaystyle \frac{8\pi ^2}{3g_m^2}}({}_{}{}^{}j,\mathrm{\Delta }^1{}_{}{}^{}j)\}.`$ The monopole Wilson loop operators are $$e^{2\pi i({}_{}{}^{}S,\mathrm{\Delta }^1dk^{(3)})},e^{2\pi i({}_{}{}^{}S,\mathrm{\Delta }^1dk^{(2)})},e^{2\pi i({}_{}{}^{}S,\mathrm{\Delta }^1dk^{(1)})}$$ (A.15) for c=R,G,B, respectively. When $`{}_{}{}^{}S`$ is a closed surface, $`({}_{}{}^{}S,\mathrm{\Delta }^1dk^{(a)})`$ in Eq. (A.15) is the 4-dimensional linking number between the closed surface $`{}_{}{}^{}S`$ and the closed loop $`k^{(a)}`$. Hence the monopole Wilson loop is independent of the choice of the surface $`S`$. This is very important. Let us consider the string representation of Eq. (A.12). Introducing the following $`U(1)^2`$ electric charges: $$𝒔_c=(s_c^{(1)},s_c^{(2)},s_c^{(3)})=\{(1,1,0),(1,0,1),(0,1,1)\}$$ (A.16) for c=R,G,B respectively , the string representation can be derived as in Section 5: $`H_c(C)={\displaystyle \frac{1}{Z_{str}}}{\displaystyle \underset{\sigma _c^{(a)}Z,\delta {}_{}{}^{}\sigma _{c}^{(a)}=j_c^{(a)}}{}}\delta _{\mathrm{\Sigma }_a\sigma _c^{(a)},0}\mathrm{exp}\{4\pi ^2\gamma \{{\displaystyle \underset{a}{}}(\sigma _c^{(a)},(\mathrm{\Delta }+m^2)^1\sigma _c^{(a)})`$ $`+{\displaystyle \frac{1}{m^2}}{\displaystyle \underset{a}{}}(j_c^{(a)},(\mathrm{\Delta }+m^2)^1j_c^{(a)})\}\},`$ (A.17) where $`j_c^{(a)}s_c^{(a)}j`$ are electric $`U(1)^2`$ currents and $`m^2=3\gamma g_m^2`$ is a mass of a dual gauge field in DGL. The two-forms $`\sigma _c^{(a)}s_c^{(a)}S+\sigma ^{(a)}`$ satisfy $`\delta {}_{}{}^{}\sigma _{c}^{(a)}=j_c^{(a)}`$ with the constraint $`_a\sigma _c^{(a)}=0`$. Thus we obtain the hadronic string model in the SU(3) lattice QCD.
warning/0002/cond-mat0002309.html	ar5iv	text	# Quenched Averaged Correlation Functions of the Random Magnets \[ ## Abstract It is shown that the ratios of the quenched averaged three and four-point correlation functions of the local energy density operator to the connected ones in the random-bond Ising model approach asymptotically to some $`universal`$ functions. We derive the explicit expressions of these universal functions. Moreover it is shown that the individual logarithmic operators have not any contribution to the connected correlation functions of the disordered Ising model. PACS: 05.70.jk, 11.25.Hf, 64.60.Ak \] Random systems represent the spatial inhomogenuity where scale invariance is only preserved on average but not for specific disorder realization. The understanding of the role played by quenched impurities of the nature of phase transition is one of the significant subjects in statistical physics and has attracted a great deal of attention . According to the Harris criterion , quenched randomness is a relevant perturbation at the second-order critical point for systems of dimension $`d`$, when its specific heat exponent $`\alpha `$, of the pure system is positive. Concerning the effect of randomness on the correlation functions, it is known that the presence of randomness induces a logarithmic factor to the correlation functions of pure system . Theoretical treatment of the quenched disordered systems is a non-trivial task in view of the fact that, one has to average the logarithm of the partition function over various realization of the disorder in the statistical ensemble and therefore find physical results. There are two standard methods to perform this averaging, the supersymmetry (SUSY) approach, and the well-known replica approach. Recently using the replica approach it has been shown by Cardy , that the logarithmic factor multiplying power law behavior are to be expected in the scaling behavior near non-mean field critical points. It is shown also that the results are valid for systems with short-range interactions and in an arbitrary number of dimensions. He concludes that in the limit of $`n0`$ of replicas the theory posses of a set of fields which are degenerate (they have the same scaling dimensions) and finds a pair of fields which form a Jordan cell structure for dilatation operator and derives logarithmic operator in such disordered systems. Cardy proves that the quenched disordered theory with $`Z=1`$ can be described by logarithmic conformal field theory as well. The logarithmic conformal field theories (LCFT) are extensions of conventional conformal field theories, which have emerged in recent years in a number of interesting physical problems of condensed matter physics , string theory , and nonlinear dynamical systems . The LCFT are characterized by the fact that their dilatation operator $`L_0`$ are not diagonalized and admit a Jordan cell structure. The non-trivial mixing between these operators leads to logarithmic singularities in their correlation functions. It has been shown that the correlator of two fields in such field theories, has a logarithmic singularity. $$<\psi (r)\psi (r^{^{}})>\|rr^{^{}}\|^{2\mathrm{\Delta }_\psi }\mathrm{log}\|rr^{^{}}\|+\mathrm{}$$ (1) In this direction we show that the quenched averaged connected correlation functions of local energy density field can be written in terms of ordinary scaling operators which can be constructed by the difference of energy operators in two different replicas. We write the connected 3 and 4-point correlation functions of energy density explicitly in terms of such ordinary operators. Furthermore we prove that the logarithmic operators have no contribution in the quenched averaged connected correlation functions of the local density operator. However, these operators play a considerable role on the disconnected ones and produce some logarithmic factors in the correlation functions. We calculate the various types of quenched averaged 3 and 4-point correlation functions of the local energy density and show that the ratios of these correlation functions to the connected ones have the specific universal asymptotic and write down these universal functions explicitly. 2-Connected Correlation Functions of Random Magnets We consider a quenched random ferromagnet, for instance an Ising model, with random-bond disorder. Let us describe this disordered system in the continuum limit by the following Hamiltonian, $$H=H^0+J(r)E(r)d^dr$$ (2) where $`H^0`$ is the Hamiltonian of the renormalization group at fixed point describing the pure Ising model. The field $`J(r)`$ is a quenched random variable coupled to the local energy density $`E(r)`$. When the coupling $`J(r)`$ is independent of $`x`$ and not random, the above Hamitonian describes the behavior of the statistical model near it’s critical point. For simplicity we assume that the random variable $`J(r)`$ is a gaussian variable which is characterized by its two moments $`<J(r)>=0`$ and $`<J(r)J(r^{})>=g\delta (rr^{})`$. The standard procedure of averaging over disorder is to introduce replicas, i.e., $`n`$ identical copies of the same model for which $$Z^n=Tr\mathrm{exp}\{\underset{a=1}{\overset{n}{}}H_a^0d^drJ(r)\underset{a=1}{\overset{n}{}}E_a(r)\}$$ (3) averaging over the disorder gives rise to the following effective replical Hamiltonaian. $$H_R=\underset{a=1}{\overset{n}{}}H_a^0g\underset{ab}{}E_a(r)E_b(r)d^dr$$ (4) we keep only the non-diagonal terms, since using the operator algebra of the pure system one can absorb the diagonal terms into $`H_a^0`$. The replicas are now coupled via the disorder. The scaling dimension of coupling $`g`$ is $`y_g=d2\mathrm{\Delta }_E`$ and is relevant at the pure fixed point if $`y_g>0`$. For small $`y_g`$ it is possible to use standard perturbation theory and find the possible random fixed point. It is noted by Cardy that the $`n`$ operator $`E_a`$ are degenerate at the pure fixed point and one can decompose them into irreducible representation of permutation group $`S_n`$. It has been shown that the combination $`E_t=_{a=1}^nE_a`$ is a singlet ( symmetric under the permutation of the replica group) and $`\stackrel{~}{E}_a=E_a\frac{1}{n}_{b=1}^nE_b`$ transforms according to an $`(n1)`$-dimensional representation of $`S_n`$. The fields $`\stackrel{~}{E}_a`$ satisfy the condition $`_{a=1}^n\stackrel{~}{E}_a=0`$. The important observation is that the fields $`E_t`$ and $`\stackrel{~}{E}_a`$ have the proper scaling dimensions close to $`n0`$ as $`\mathrm{\Delta }_{E_t}=\mathrm{\Delta }_{E_a}^{(0)}+\frac{1}{2}(1n)y_g+O(y_g^2)`$ and $`\mathrm{\Delta }_{\stackrel{~}{E}}=\mathrm{\Delta }_E^{(0)}+\frac{1}{2}y_g+O(y_g^2)`$ respectively. It is clear that the singlet field $`E_t`$ becomes degenerate with the $`(n1)`$ operators $`\stackrel{~}{E}_a`$. This is true to all orders . However they do not form the basis of the Jordan cell for the dilatation operator. To find the logarithmic pair according to we define the correlation function of $`<E_t(0)E_t(r)>=A_1`$ and $`<\stackrel{~}{E}_a(0)\stackrel{~}{E}_a(r)>=B_1`$ and find the following relations for $`A_1`$ and $`B_1`$. $`A_1`$ $`=`$ $`n(a(n1)b)nA(n)r^{2\mathrm{\Delta }_E(n)}`$ (5) $`B_1`$ $`=`$ $`(1{\displaystyle \frac{1}{n}})(ab)(1{\displaystyle \frac{1}{n}})B(n)r^{2\stackrel{~}{\mathrm{\Delta }}_E(n)}`$ (7) where $`a=<E_i(0)E_i(r)>`$ and $`b=<E_i(0)E_j(r)>`$ with $`ij`$. the above equations enable us to write the quenched averaged connected two-point correlation functions of energy density operator in terms of $`a`$ and $`b`$ in the limit of $`n0`$ as: $`\overline{<E(0)E(r)>}_c=ab`$ which is equal to $`B(0)r^{2\mathrm{\Delta }_E}`$ and it has a pure scaling behavior. However, the correlation functions $`a`$ and $`b`$ have the logarithmic singularities and behave as: $`<E_1(0)E_1(r)>`$ $`=`$ $`(A^{}(0)B^{}(0)+B(0)`$ (8) $``$ $`B(0){\displaystyle \frac{y_g}{2}}\mathrm{ln}r)r^{2\mathrm{\Delta }_E}`$ (10) $`<E_1(0)E_2(r)>`$ $`=`$ $`(A^{}(0)B^{}(0)A(0){\displaystyle \frac{y_g}{2}}\mathrm{ln}r)r^{2\mathrm{\Delta }_E}`$ (12) where $`A(0)=B(0)`$. The prime sign in the eq. (6) means differentiating with respect to $`n`$. This means that in the limit $`n0`$ the fields $`E_t`$ and $`E_a`$ form a basis of Jordan cell, i.e. their two point correlation functions behave as: $`<E_t(0)E_t(r)>=0`$, $`<E_t(0)E_a(r)>=a_1r^{2\mathrm{\Delta }_E}`$ and $`<E_a(0)E_b(r)>=(2a_1\mathrm{ln}r+D_{a,b})r^{2\mathrm{\Delta }}`$, where $`a_1`$ and $`D_{a,b}`$ are some constants. As noted by Cardy the ratio of quenched averaged two-point correlators of the energy density operator to the connected one has a universal $`r`$-dependence as: $$\frac{\overline{<E(0)E(r)>}}{\overline{<E(0)E(r)>}_c}\frac{\overline{<E(0)><E(r)>}}{\overline{<E(0)E(r)>}_c}\mathrm{ln}r$$ (13) To understand the structure of Jordan cell, we note that in 2D one can define the operator $`L_0`$ as $$L_0=\left(\begin{array}{cc}\mathrm{\Delta }_E\hfill & 0\hfill \\ 1\hfill & \mathrm{\Delta }_E\hfill \end{array}\right)$$ (14) so that, $`L_0E_t=\mathrm{\Delta }_EE_t`$ and $`L_0E_a=\mathrm{\Delta }_EE_a+E_t`$ in the limit of $`n0`$. Using this representation for $`L_0`$ one can show that the field $`E_t`$ with its logarithmic partner $`E_a`$ have the standard logarithmic correlation functions , (see the correlation above the eq.(8) ). We note that in 2D we have dealt with two-dimensional conformal field theory, relying heavily on the underlying Virasoro algebra. For an extention to $`D`$ dimensions one has to modify the representation of the Virasoro algebra to higher dimensions . We consider a doublet of fields (Jordan cell) $`\mathrm{\Phi }=\left(\begin{array}{c}E_t\hfill \\ E_a\hfill \end{array}\right)`$ and note that under D-dimensional conformal transformation $`𝐱𝐱^{}`$, we have, $`\mathrm{\Phi }(𝐱)𝚽^{}(𝐱^{})=𝐆^𝐓𝚽(𝐱)`$ where $`T`$ is a two dimensional matrix which has Jordan form and $`G=\frac{x^{}}{x}`$ is the Jacobian. For our particular case $`T`$ has the following Jordan form: $$T=\left(\begin{array}{cc}\frac{2\mathrm{\Delta }_E}{D}\hfill & 0\hfill \\ 1\hfill & \frac{2\mathrm{\Delta }_E}{D}\hfill \end{array}\right)$$ (15) and one can show that the two fields $`E_t`$ and $`E_a`$, transform as: $`E_t(𝐱^{})`$ $`=`$ $`G^{\frac{2\mathrm{\Delta }_E}{D}}E_t(𝐱)`$ (16) $`E_a(𝐱^{})`$ $`=`$ $`G^{\frac{2\mathrm{\Delta }_E}{D}}(\mathrm{ln}(G)E_t(x)+E_a(x))`$ (18) This expresses that the top-field $`E_t`$ always transforms as an ordinary scaling operator. It can be verified that the correlation functions of fields $`E_t`$ and $`E_a`$ have the standard D- dimensional logarithmic conformal field theory structure . Using the above results, it is evident that the dimension of field-difference $`E_aE_b`$ with $`ab`$ is $`\mathrm{\Delta }_E`$ and it transforms as an ordinary operator under the scaling transformation. The intresting observation is that the connected averaged correlation functions depends on the difference fields $`E_aE_b`$ only and therefore they behave as the ordinary correlation functions. For instance in the following we write the connected quenched averaged 2,3 and 4-point functions of local energy density in terms of the field-difference operators explicitly, $`\overline{<E(1)E(2)>}_c`$ $`=`$ $`{\displaystyle \frac{1}{2}}<(E_aE_b)_{(1)}(E_aE_b)_{(2)}>`$ (19) $`\overline{<E(1)E(2)E(3)>}_c`$ $`=`$ $`<(E_aE_b)_{(1)}`$ (22) $`(E_aE_c)_{(2)}(E_aE_b)_{(3)}>`$ $`\overline{<E(1)E(2)E(3)E(4)>}_c=`$ (23) (24) $`<(E_aE_b)_{(1)}(E_aE_c)_{(2)}(E_aE_d)_{(3)}(E_aE_b)_{(4)}>`$ (25) (26) $`{\displaystyle \frac{1}{2}}<(E_aE_b)_{(1)}(E_cE_d)_{(2)}(E_cE_d)_{(3)}(E_aE_b)_{(4)}>`$ (27) (28) $`{\displaystyle \frac{1}{4}}<(E_aE_b)_{(1)}(E_cE_d)_{(2)}(E_aE_b)_{(3)}(E_cE_d)_{(4)}>`$ (29) (30) $`{\displaystyle \frac{1}{4}}<(E_aE_b)_{(1)}(E_aE_b)_{(2)}(E_cE_d)_{(3)}(E_cE_d)_{(4)}>`$ (31) where the last equation has only 15 independent terms. To confirm this prediction also one can directly show that the quenched averaged connected correlation functions have a pure scaling behaviour which is determined by ordinary scaling operators and the logarithmic operators $`E_a`$ do not change its behavior. This can be verified directly for the quenched averaged connected 3-point correlation function of energy density. We are interested in deriving exactly the various 3-point quenched averaged functions as $`\overline{<E(1)E(2)E(3)>}`$, $`\overline{<E(1)E(2)><E(3)>}`$ and $`\overline{<E(1)><E(2)><E(3)>}`$, which can be written in terms of the replica correlation functions $`<E_1(1)E_1(2)E_1(3)>=a`$ $`<E_1(1)E_1(2)E_2(3)>=b`$ and $`<E_1(1)E_2(2)E_3(3)>=c`$, respectively. One can derive the correlation functions $`a`$, $`b`$ and $`c`$ by means of 3-point functions of $`E_t`$ and $`\stackrel{~}{E}_a`$ as follows: $`<E_t(1)E_t(2)E_t(3)>=na+3n(n1)b`$ (32) (33) $`n(n1)(n2)cnA_1`$ (34) $`<\stackrel{~}{E}_a(1)\stackrel{~}{E}_a(2)E_t(3)>=n_1a+(n_1^2(n1)`$ (37) $`4n_1^2+{\displaystyle \frac{1}{n^2}}(n1)^2+{\displaystyle \frac{2}{n^2}}(n1)(n2))b`$ $`+`$ $`({\displaystyle \frac{2}{n}}n_1(n1)(n2)+{\displaystyle \frac{1}{n^2}}(n2)^2(n1))c`$ (41) $`(1{\displaystyle \frac{1}{n}})B_1`$ and finally, $`<\stackrel{~}{E}_a(1)\stackrel{~}{E}_a(2)\stackrel{~}{E}_a(3)>=(n_1^2{\displaystyle \frac{n1}{n^3}})a`$ (42) (43) $`(3n_1^2{\displaystyle \frac{n1}{n}}{\displaystyle \frac{3}{n^3}}(n1)(n2)+{\displaystyle \frac{3}{n^2}}n_1(n1))b`$ (44) (45) $`+({\displaystyle \frac{3}{n^2}}n_1(n1)(n2){\displaystyle \frac{1}{n^3}}(n1)(n2)(n3))c`$ (46) (47) $`(1{\displaystyle \frac{1}{n}})(1{\displaystyle \frac{2}{n}})C_1`$ (48) where $`n_1=(1\frac{1}{n})`$ and $`A_1`$, $`B_1`$ and $`C_1`$ are pure scaling functions of variables $`r_{i,j}`$. To derive the above equations we use the replica symmetry and symmetry of the various types of 3-point correlation functions under interchanging of positions. We note that replica symmetry leads to have $`<\stackrel{~}{E}_a(1)E_t(2)E_t(3)>=0`$ and therefore, dose not give any new relationship between $`a`$, $`b`$ and $`c`$. Using the above equations, it can be found that the correlation functions $`a`$, $`b`$ and $`c`$ are as follows: $`a`$ $`=`$ $`{\displaystyle \frac{3nB_13nC_1+n^2C_1+A_13B_1+2C_1}{n^2}}`$ (49) $`b`$ $`=`$ $`{\displaystyle \frac{nB_1nC_1+A_13B_1+2C_1}{n^2}}`$ (51) $`c`$ $`=`$ $`{\displaystyle \frac{A_13B_1+2C_1}{n^2}}`$ (53) Using the above equations we can show that the connected quenched averaged 3-point function behaves as: $`\overline{<E(1)E(2)E(3)>}`$ $`=`$ $`2c+a3b`$ (54) $`=`$ $`C_1`$ (56) which is a scaling function and confirms the observation that the logarithmic operators (individually) have no role in the connected quenched averaged correlation functions. In addition one can derive the correlation functions $`<E_i(1)E_j(2)E_k(3)>`$ for given $`i`$, $`j`$ and $`k`$ in the limit of $`n0`$ and show that they have the following form: $`<E_i(1)E_j(2)E_k(3)>`$ $`=`$ $`[\alpha _{ijk}\beta _{ijk}D_1`$ (57) $`+`$ $`\gamma _{ijk}(4D_2D_1^2)]f(1,2,3)`$ (59) where $`f(1,2,3)=(r_{12}r_{13}r_{23})^{2\mathrm{\Delta }_E}`$ , $`D_1=\mathrm{ln}(r_{12}r_{13}r_{23})`$ and $`D_2=\mathrm{ln}r_{23}\mathrm{ln}r_{13}+\mathrm{ln}r_{13}\mathrm{ln}r_{12}+\mathrm{ln}r_{23}\mathrm{ln}r_{12}`$. It can also be shown that the ratio of various symmetrized 3-point functions to the connected one behaves asymptotically as a $`universal`$ function $$\frac{1}{3}(4D_2D_1^2).$$ (60) We generalize the above calculations to derive the various type of 4-point correlation functions and show that the ratio of the various disconnected to the connected one have the following universal asymptotic: $$\frac{1}{36}[O_1^36O_23O_312O_418O_5]$$ (61) where $`O_1=\mathrm{ln}(r_{12}r_{13}r_{14}r_{23}r_{24}r_{34})`$, $`O_2=(\mathrm{ln}r_{ij}\mathrm{ln}r_{kl}^2+\mathrm{})`$ with $`ijkl`$, $`O_3=(\mathrm{ln}r_{ij}\mathrm{ln}r_{ik}^2+\mathrm{})`$ with $`ijk`$, $`O_4=(\mathrm{ln}r_{ij}\mathrm{ln}r_{kl}\mathrm{ln}r_{lj}+\mathrm{})`$ with $`ijkl`$, and finally $`O_5=(\mathrm{ln}r_{ij}\mathrm{ln}r_{ik}\mathrm{ln}r_{il}+\mathrm{})`$ with $`ijkl`$ In summary, in this paper we have studied the correlation functions of disordered random magnets and obtain the various types of 3 and 4-point quenched averaged correlation functions. One can check directly that these different types of the 3 and 4-point correlation functions have the general property of a logarithmic conformal field theory that the logarithmic partner can be regarded as the formal derivative of the ordinary fields (top field) with respect to their conformal weight . In this case, one can consider the $`E_a`$ fields as the derivative of $`E_t`$ with respect to $`n`$ . We emphasise that the derivative with respect to scaling weight can be written in terms of the derivative with respect to $`n`$. These properties enable us to calculate any N-point correlation function containing the logarithmic field $`E_a`$ in terms of the correlation functions of the top-fields. The general expression of the correlation functions of the LCFT’s are given in ref. and here we determine the unknown constants in the logarithmic correlation functions in terms of details of the random-bond Ising model. It is noted that the formal derivations with respect to scaling dimensions can not predict the unknown constants in the quenched averaged correlation functions of the local energy density operators. The constant depends on the detail of the statistical model. We have shown that the individual logarithmic operators $`E_a`$ do not have any contribution to the quenched averaged connected correlation functions of the energy density. We also obtain that the connected correlation functions can be written in terms of the difference fields which transform as an ordinary scaling operator. However they will play a crucial role to the disconnected averaged correlation functions. Also we find that the ratio of the various types of 3 and 4-point quenched averaged correlation functions to the connected ones have a universal asymptotic behavior and give their explicit form. These predictions can also be investigated numerically. Our analysis are valid in all dimensions as long as the dimension is below the upper critical dimensions.To derive the above results we have used the replica symmetry. Any attempt towards the breaking of this symmetry will change completely the above picture and produces more than one logarithmic fields in the block and produces higher order logarithmic singularities . These results can be easily generalized to other problem such as polymer statistics, percolation and random phase sine-Gordon model etc. We would like to thank John Cardy for his useful comments and A. Aghamohammadi and R. Asgari, B. Davoudi and J. Davoudi for thier useful discussions. This paper is dedicated to Dr. A. M. Zaker, associate professor of Physics Department, who deceased last may 2000.
warning/0002/cs0002006.html	ar5iv	text	# Multiplicative Nonholonomic/Newton -like Algorithm ## 1 Introduction Suppose that $`N`$-dimensional stochastic variables $`\{X_i\|1iN\}`$ are observed. The independent component analysis (ICA) pursues a map $`XY`$, where each component of $`Y`$ becomes mutually independent. In this letter we restrict ourselves to the linear independent component analysis. There we want to find a linear transformation $`C:𝐗=(X_1,\mathrm{},X_N)^{}𝐘=(Y_1,\mathrm{},Y_N)^{}=C𝐗`$ which minimizes some cost function that measures the independence. Hereafter we denote by the upper subscript $``$ the transposition and by $``$ the complex conjugate. There can be many candidates for the cost function. For example the Kullback-Leibler information is a good measure for the independence. In this case the problem is translated to the minimization of $`_{i=1}^N𝑑y_iP_i(y_i)\mathrm{ln}P_i(y_i)`$, where $`P_i`$ is the probability density function of the $`i`$-th component. It is obvious that we must evaluate $`P_i`$’s to find the optimal solution. A robust estimation of the probability density functions is not an easy task and if it is possible it may be computationally expensive. An alternative idea is to make use of the cumulant of the fourth order, or the kurtosis\[A.Hyvärinen,1997\], which we will adopt in this letter. The fourth order cumulant vanishes for the normal distribution. So, this cost function is robust under the gaussian random noises. We will construct algorithms where a matrix, which specifies the linear transformation, is updated by the left-multiplication of a matrix $`D=\mathrm{e}^\mathrm{\Delta }`$. This expression implies that $`D`$ belongs to $`GL(N,R)`$ (more accurately, the component of $`GL(N,R)`$ connected to the unit element), which ensures the conservation of the rank. The specification of $`D`$ by the coordinate $`\mathrm{\Delta }`$ has many advantages since it has a compatibility with the homogeneous nature of the Lie group. There are variations for the form of the cost function. We will show our definitions in the following two sections, which are choosen to possess invariance under componentwise scaling. This invariance is crucial for a rigorous treatment of the convergence properties. Moreover, this invariance allows us to identify points in $`GL(N,R)`$ which is transformed to each other by the scaling. Then we can legitimately restrict the dynamics to a coset space which is introduced by this identification. Under these settings, we determine $`\mathrm{\Delta }`$ by using the Newton method for the second order expansion of the cost function with respect to $`\{\mathrm{\Delta }_{ij}\}`$. It is assumed that the diagonal elements of $`\mathrm{\Delta }`$ are zeros, which does not impose any restrictions. That is, a point can move toward any direction in this coset by a left-multiplication of $`\mathrm{e}^\mathrm{\Delta }`$. Thus it is not necesarry for our method to prewhiten the data. It is also not required that the optimal solution is the maximum or the minimum of the cost function. Indeed, the sole requirement is that the optimal point is a saddle point of the cost function since our method is in principle the Newton method. These are great advantages of our method. Our strategy is as follows. As an initial condition we set $`C_0`$. For $`t>0(t𝐍^+)`$, we introduce an $`N\times N`$ matrix $`\mathrm{\Delta }_t`$ and denote $`C_t`$ as $`C_t=\mathrm{e}^{\mathrm{\Delta }_t}C_{t1}`$. Next, we evaluate the cost function at $`C_t`$ by using the expansion around $`C_{t1}`$ with respect to the elements of $`\mathrm{\Delta }_t`$ up to the second order. Then $`\mathrm{\Delta }_t`$ is choosen as a saddle point of this second order expansion. We iteratively follow these procedures until we obtain a satisfactory solution. This letter is organized as follows. In Section 2 the main part of our algorithm is constructed, where the cost function is essentially identical to the sum of kurtoses. We adopt the square of the kurtoses for the cost function in Section 3. Explicit expressions for the optimal $`\mathrm{\Delta }`$ (up to the second order) are obtained both in Sections 2 and 3. Section 4 is a short section where we show how each updating step is combined to obtain the optimal $`C`$. In Section 5 the convergence property of our algorithm is discussed. Section 6 contains conclusions and discussions. ## 2 Multiplicative update algorithm ### 2.1 Expansion of the cost function Let us start by defining the cost function: $`f(C,X)={\displaystyle \underset{i}{}}f_i(C,X),`$ (2.1) where $`f_i`$’s are the fourth order moments of components divided by the square of their variances, $`f_i(C,X)={\displaystyle \frac{E((CX)_i^4)}{E((CX)_i^2)^2}}.`$ (2.2) In this letter we denote by $`E(A)`$ the expectation of $`A`$. Obviously the cost function $`f`$ coincides with the sum of kurtoses of all the components up to the constant. We set $`D=\mathrm{e}^\mathrm{\Delta }`$ and expand $`f(D,Y)`$ in terms of the elements of $`\mathrm{\Delta }`$. For example expansions term by term are evaluated as follows: $`E((DY)_i^4)`$ $`=`$ $`E(Y_i^4)+4{\displaystyle \underset{p}{}}(\mathrm{\Delta }_{ip}+({\displaystyle \frac{\mathrm{\Delta }^2}{2}})_{ip})E(Y_i^3Y_p)+6{\displaystyle \underset{p,q}{}}\mathrm{\Delta }_{ip}\mathrm{\Delta }_{iq}E(Y_i^2Y_pY_q)+O(\mathrm{\Delta }^3)`$ $`E((DY)_i^2)`$ $`=`$ $`E(Y_i^2)+2{\displaystyle \underset{p}{}}(\mathrm{\Delta }_{ip}+({\displaystyle \frac{\mathrm{\Delta }^2}{2}})_{ip})E(Y_iY_p)+{\displaystyle \underset{p,q}{}}\mathrm{\Delta }_{ip}\mathrm{\Delta }_{iq}E(Y_pY_q)+O(\mathrm{\Delta }^3).`$ (2.3) Hereafter we denote by $`O(\mathrm{\Delta }^k)`$ polynomials of matrix elements of $`\mathrm{\Delta }`$ which does not contain terms with degrees less than $`k`$. For brevity’s sake we introduce the following notations: $`\sigma _i^{(k)}=\|E(Y_i^k)\|^{1/k},`$ (2.4) $`R_{pi}^{(k)}={\displaystyle \frac{E(Y_i^kY_p)}{(\sigma _i^{(2)})^{k+1}}},`$ (2.5) $`U_{pq}^{(k,i)}={\displaystyle \frac{E(Y_i^kY_pY_q)}{(\sigma _i^{(2)})^{k+2}}},`$ (2.6) and $`\kappa _i=(\sigma _i^{(4)})^4/(\sigma _i^{(2)})^4.`$ (2.7) Using the quantities defined above we can show that the cost function is expanded as $`f_i(D,Y)`$ $`=`$ $`\left[\kappa _i+4\left[(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Delta }^2}{2}})R^{(3)}\right]_{ii}+6\left[\mathrm{\Delta }U^{(2,i)}\mathrm{\Delta }^{}\right]_{ii}+O(\mathrm{\Delta }^3)\right]`$ (2.8) $`\times \left[14\left[(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Delta }^2}{2}})R^{(1)}\right]_{ii}2\left[\mathrm{\Delta }U^{(0,i)}\mathrm{\Delta }^{}\right]_{ii}+12\left[\mathrm{\Delta }R^{(1)}\right]_{ii}^2+O(\mathrm{\Delta }^3)\right]`$ $`=`$ $`\kappa _i4\left[(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Delta }^2}{2}})(\kappa _iR^{(1)}R^{(3)})\right]_{ii}+2\left[\mathrm{\Delta }(3U^{(2,i)}\kappa _iU^{(0,i)})\mathrm{\Delta }^{}\right]_{ii}`$ $`+12\kappa _i\left[\mathrm{\Delta }R^{(1)}\right]_{ii}^216\left[\mathrm{\Delta }R^{(1)}\right]_{ii}\left[\mathrm{\Delta }R^{(3)}\right]_{ii}+O(\mathrm{\Delta }^3)`$ by straightforward calculations. Next, we evaluate partial derivatives of the cost function by the matrix elements of $`\mathrm{\Delta }`$. Partially differentiating (2.8), we get an expression, $`{\displaystyle \frac{f(\mathrm{e}^\mathrm{\Delta },Y)}{\mathrm{\Delta }_{kl}}}=4\left[KR^{(3)}\right]_{lk}2\left[(KR^{(3)})\mathrm{\Delta }+\mathrm{\Delta }(KR^{(3)})\right]_{lk}`$ $`+4\left[(3U^{(2,k)}\kappa _kU^{(0,k)})\mathrm{\Delta }^{}\right]_{lk}+24K_{lk}\left[\mathrm{\Delta }R^{(1)}\right]_{kk}16R_{lk}^{(1)}\left[\mathrm{\Delta }R^{(3)}\right]_{kk}16R_{lk}^{(3)}\left[\mathrm{\Delta }R^{(1)}\right]_{kk}`$ $`+O(\mathrm{\Delta }^2),`$ (2.9) where $`K`$ is an $`N\times N`$ matrix defined by $`K_{pq}=\kappa _qR_{pq}^{(1)}.`$ (2.10) We want to decide $`\mathrm{\Delta }`$ for which the partial derivative by $`\mathrm{\Delta }_{kl}(kl)`$ of the cost function vanish on condition that $`\mathrm{\Delta }_{ii}=0`$ for $`1iN`$. We neglect $`O(\mathrm{\Delta }^3)`$ terms in the cost function. Thus the right-hand side of (2.1) is regarded as a polynomial of $`\{\mathrm{\Delta }_{kl}\}`$ of at most first order and it is always possible in principle to determine $`\mathrm{\Delta }`$ which satifies the above condition. It is, at the same time, not easy to describe the problem in a form which is valid for arbitrary $`N`$. In the following subsection we will introduce a transparent and unified method for handling the partial derivatives of $`f`$. We leave this subsection by introducing $`N\times N`$ matrices $`V^{(i)}=3U^{(2,i)}\kappa _iU^{(0,i)}`$ (2.11) and $`Q=KR^{(3)}`$ (2.12) for later convenience. ### 2.2 Expression by tensor product and determination of $`\mathrm{\Delta }`$ The expression (2.1) is quite complicated and not convenient for our purpose, “ determine $`\mathrm{\Delta }`$, where all the partial derivatives vanish”. Fortunately by mapping the relations between elements of $`N\times N`$ matrices to those of $`N^2\times N^2`$ matrices, we can handle the problem transparently. Some preparations are needed. First, let us introduce a map $`\mathrm{cs}`$: $`\mathrm{Mat}(N,F)`$ $``$ $`F^{N^2}`$ $`A=\left(\begin{array}{cccc}A_{11}& A_{12}& \mathrm{}& A_{1N}\\ A_{21}& \multicolumn{3}{c}{}\\ \multicolumn{4}{c}{}\\ A_{N1}& \multicolumn{2}{c}{}& A_{NN}\end{array}\right)`$ $``$ $`\mathrm{cs}(A)=(A_{11}A_{21}\mathrm{}A_{N1}A_{12}A_{22}\mathrm{}A_{NN})^{},`$ (2.17) where $`F`$ is an unspecified field. We also introduce two useful operators $`T`$ and $`P`$. The “intertwiner” $`T`$ is an $`N^2\times N^2`$ matrix defined by $`\mathrm{cs}(A^{})=T\mathrm{cs}(A)\text{for }A\mathrm{Mat}(N,F).`$ (2.19) The projection operator $`P`$, $`P`$ $`=`$ $`\mathrm{diag}(p_1,\mathrm{},p_{N^2}),`$ (2.22) $`\{\begin{array}{cc}p_k=1\text{for}k=N(i1)+i,1iN\hfill & \\ p_k=0\text{otherwise},\hfill & \end{array}`$ is used to extract the “diagonal” elements of a matrix from its image by $`\mathrm{cs}`$. On this setting we can rewrite (2.1) as $`{\displaystyle \frac{f(\mathrm{e}^\mathrm{\Delta },Y)}{\mathrm{\Delta }_{kl}}}`$ $`=`$ $`[4\mathrm{c}\mathrm{s}(Q)2[I_NQ+T(I_NQ^{})T]\mathrm{cs}(\mathrm{\Delta })+4\left\{{\displaystyle \underset{i=1}{\overset{N}{}}}V^{(i)}\right\}\mathrm{cs}(\mathrm{\Delta }^{})`$ (2.23) $`+\{24(I_NK)P(IR^{(1)})^{}16(I_NR^{(1)})P(IR^{(3)})^{}`$ $`16(I_NR^{(3)})P(IR^{(1)})^{}\}\mathrm{cs}(\mathrm{\Delta }^{})]_{l+N(k1)},`$ where $`I_N`$ is the $`N\times N`$ unit matrix and $`{\displaystyle \underset{i=1}{\overset{N}{}}}V^{(i)}=\left(\begin{array}{ccccc}V^{(1)}\hfill & 0\hfill & \multicolumn{2}{c}{\mathrm{}\mathrm{}}& 0\hfill \\ 0\hfill & V^{(2)}\hfill & 0\hfill & \multicolumn{2}{c}{\mathrm{}\mathrm{}}\\ \multicolumn{5}{c}{}\\ \multicolumn{5}{c}{}\\ 0\hfill & \multicolumn{2}{c}{\mathrm{}\mathrm{}}& V^{(N1)}\hfill & 0\hfill \\ 0\hfill & 0\hfill & \multicolumn{2}{c}{\mathrm{}\mathrm{}}& V^{(N)}\hfill \end{array}\right).`$ (2.30) We make use of the following fact: For $`X\mathrm{Mat}(N,F)`$ $`T(I_NX)T=XI_N.`$ (2.31) See Appendix A for the proof of (2.31). Then (2.23) becomes $`{\displaystyle \frac{f(\mathrm{e}^\mathrm{\Delta },Y)}{\mathrm{\Delta }_{kl}}}=4[\mathrm{cs}(Q)]_{l+N(k1)}+\left[W\mathrm{cs}(\mathrm{\Delta })\right]_{l+N(k1)},`$ where $`W`$ $`=`$ $`2(I_NQ+Q^{}I_N)+4\left\{{\displaystyle \underset{i=1}{\overset{N}{}}}V^{(i)}\right\}T+[24(I_NK)P(IR^{(1)})^{}`$ $`16(I_NR^{(1)})P(IR^{(3)})^{}16(I_NR^{(3)})P(IR^{(1)})^{}]T.`$ Now let us determine $`\mathrm{\Delta }`$. Remember that we are going along the spirit of the Newton method. Thus we want to find $`\mathrm{\Delta }`$ which satisfies the conditions $`{\displaystyle \frac{f(\mathrm{e}^\mathrm{\Delta },Y)}{\mathrm{\Delta }_{kl}}}=0+O(\mathrm{\Delta }^2)\text{for }1k,lN,kl`$ (2.34) and $`\mathrm{\Delta }_{kk}=0\text{for}1kN.`$ (2.35) The conditions (2.35) make the problem rather complicated one. Fortunately, by using $`P`$ we can combine the conditions (2.34) and (2.35) into a matrix equation : $`\left[(I_{N^2}P)W(I_{N^2}P)+P\right]\mathrm{cs}(\mathrm{\Delta })4(I_{N^2}P)\mathrm{cs}(Q)=0.`$ (2.36) Immediately it follows that $`\mathrm{cs}(\mathrm{\Delta })=4\left[(I_{N^2}P)W(I_{N^2}P)+P\right]^1(I_{N^2}P)\mathrm{cs}(Q).`$ (2.37) Thus we have obtained $`\mathrm{\Delta }`$ which specify a saddle point of the expansion of $`f(C,Y)`$ up to the second order. Note that quantities in the right-hand side of (2.37) are easily estimated ones from the observed data. So, an updating is determined by (2.37) without any ambiguities. ## 3 Case $`\mathrm{I}\mathrm{I}`$: square of kurtosis Obviously, points where kurtosis vanishes do not play any special role for the cost function $`f`$ in Section 2. The optimal solution, however, contains components with zero kurtoses when the number of the sources is less than that of the observation channels. Thus, in this section we treat with a slightly different cost function, which is the sum, $`𝒇(C,X)={\displaystyle \underset{i}{}}𝒇_i(C,X),`$ (3.1) of the square of the kurtoses, $`𝒇_i(C,X)=\left[{\displaystyle \frac{E((CX)_i^4)}{E((CX)_i^2)^2}}3\right]^2.`$ (3.2) As in the last section, we want to know the saddle point $`D=\mathrm{e}^\mathrm{\Delta }`$ of the expansion of $`𝒇_i(D,Y)`$ in terms of $`\{\mathrm{\Delta }_{ij}\}`$ up to the second order. We do not describe details of the calculations in this section, which is carried out almost in the same way as in Section 2. First, the expansion of $`𝒇_i(D,Y)`$ is evaluated as $`𝒇_i(D,Y)`$ $`=`$ $`(\kappa _i3)^28\left[(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Delta }^2}{2}})(R^{(1)}\kappa _iR^{(3)})\right]_{ii}(\kappa _i3)`$ (3.3) $`+4\left[\mathrm{\Delta }(3U^{(2,i)}\kappa _iU^{(0,i)})\mathrm{\Delta }^{}\right]_{ii}(\kappa _i3)+16\left[\mathrm{\Delta }(R^{(1)}\kappa _iR^{(3)})\right]_{ii}^2`$ $`+24(\kappa _i3)\kappa _i\left[\mathrm{\Delta }R^{(1)}\right]_{ii}^232(\kappa _i3)\left[\mathrm{\Delta }R^{(1)}\right]_{ii}\left[\mathrm{\Delta }R^{(3)}\right]_{ii}+O(\mathrm{\Delta }^3).`$ Next, we introduce $`N\times N`$ matrices $`𝑲`$, $`\{𝑽^{(i)}\|1iN\}`$, $`𝑺`$, and $`𝑸`$ defined respectively by $`𝑲_{pq}=2R_{pq}^{(1)}(\kappa _q3)\kappa _q,`$ (3.4) $`𝑽^{(i)}=2(\kappa _i3)(3U^{(2,i)}\kappa _iU^{(0,i)}),`$ (3.5) $`𝑺=\mathrm{diag}(2(\kappa _i3)),`$ (3.7) and $`𝑸_{pq}=2(\kappa _q3)(R_{pq}^{(1)}\kappa _qR_{pq}^{(3)}).`$ (3.8) We also rewrite $`Q`$ in (2.12) by $`𝒒`$ in order to avoid confusions: $`𝒒_{pq}=(R_{pq}^{(1)}\kappa _qR_{pq}^{(3)}).`$ (3.9) Now we proceed to the expression by using the tensor product. We can show that the gradients of the cost function have the following expression: $`{\displaystyle \frac{𝒇(\mathrm{e}^\mathrm{\Delta },Y)}{\mathrm{\Delta }_{kl}}}=4[\mathrm{cs}(𝑸)]_{l+N(k1)}+\left[𝑾\mathrm{cs}(\mathrm{\Delta })\right]_{l+N(k1)}+O(\mathrm{\Delta }^2),`$ where $`𝑾`$ $`=`$ $`2(I_N𝑸+𝑸^{}I_N)+4\left\{{\displaystyle \underset{i=1}{\overset{N}{}}}𝑽^{(i)}\right\}T+[24(I_N𝑲)P(IR^{(1)})^{}`$ (3.11) $`+32(I_N𝒒)P(I_N𝒒)^{}16(I_NR^{(1)}𝑺)P(IR^{(3)})^{}`$ $`16(I_NR^{(3)}𝑺)P(IR^{(1)})^{}]T.`$ This is a completely analogous expression to (2.2). Thus the coordinate $`\mathrm{\Delta }`$ of the saddle point of the second order expansion is determined by $`\mathrm{cs}(\mathrm{\Delta })=4\left[(I_{N^2}P)𝑾(I_{N^2}P)+P\right]^1(I_{N^2}P)\mathrm{cs}(𝑸).`$ (3.12) In many cases obtained through the two cost functions in Section 2 and Section 3 are almost the same results. As implied at the beginning of this section, the main difference between these two lies in the points where the kurtosis of one of the components vanishes. These point indeed constitue saddle points of the cost function $`f`$, while it is impossible to capture them by the algorithm in Section 2. Thus, we must choose an appropriate method for individual problems having this differnce in mind. ## 4 Iteration of updating Now we have obtained the updating rules. It is not necessary to tune the learning rate. Apparently, (2.36) and (3.12) look complicated. They are, however, easily implemented by the numerical tools like MatLab. (The source codes will be available from our Web-site. ) Starting from $`C_0`$, $`C_i`$ for positive $`i`$ is determined by the left multiplication by $`\mathrm{e}^{\mathrm{\Delta }_i}`$, where $`\mathrm{\Delta }`$ is determined by setting $`Y=C_{i1}X`$, i.e, $`C_t=\mathrm{e}^{\mathrm{\Delta }_t}\mathrm{e}^{\mathrm{\Delta }_{t1}}\mathrm{e}^{\mathrm{\Delta }_{t2}}\mathrm{}\mathrm{e}^{\mathrm{\Delta }_1}C_0.`$ (4.1) If $`\mathrm{\Delta }`$ becomes saficiently small, we can stop the iteration and exit the process. ## 5 Second order convergence First, we will take over the notations in Section 2. The following discussion is, however, valid for the algorithm in Section 3 if we substitute the quantities $`f`$, $`W`$, and so on by their boldface counterparts. Let us start this section by introducing some additional notations. We set $`GGL(N,R)`$ (5.1) and $`KGL(1,R)^N.`$ (5.2) We also define the coset space $`K\backslash G`$ by introducing the equivalence relation $`g^{}g^1Kgg^{}`$ (5.3) to $`G`$. That is, $`K\backslash G\{Kg\|gG\}`$. Our method is understood as an orthodox adaptation of the Newton method to this coset space $`K\backslash G`$. Note that the cost function $`F()\stackrel{\mathrm{def}}{=}f(,Y)`$ on $`G`$ satisfies the relation $`F(g)=F(Kg).`$ (5.4) So $`F`$ is naturally considered as a function on $`K\backslash G`$. That is the reason of our choice for the cost function. Thus, the second-order convergence immediately follows if the the correction to the error with respect to the coordinating resulting from the multiplicative nature is properly evaluated. At time $`t`$, a point $`g`$ on $`K\backslash G`$ is specified by the coordinate $`X^{(t)}(g)𝔪`$ such that $`\mathrm{e}^{X^{(t)}(g)}C_tg,`$ (5.5) where $`𝔪`$ is the set of $`N\times N`$ matrices whose diagonal elements are zeros. Actually, this statement itself is not a thing of course, for which the proof will be given elsewhere. Define $`F_t`$, the representation of the cost function at $`t`$, by $`F_t(X)=F(\mathrm{e}^XC_t).`$ (5.6) Here we introduce an $`(N^2N)\times N^2`$ matrix $`\stackrel{~}{P}`$ by drawing out the $`i+N(i1)`$-th raws from the unit $`N^2\times N^2`$ matrix where $`i=N,N1,\mathrm{},2,1`$. We will denote by $`H^{(t)}`$ the Hessian, $`H_{kl}^{(t)}={\displaystyle \frac{^2F_t(X)}{(\stackrel{~}{P}\mathrm{cs}(X))_k(\stackrel{~}{P}\mathrm{cs}(X))_l}}`$ (5.7) Note that if we set $`h_t(X)=T\left((I_{N^2}P)W(I_{N^2}P)+P\right)\|_{C=\mathrm{e}^XC_t},`$ (5.8) the Hessian is written as $`H^{(t)}=\stackrel{~}{P}h_t\stackrel{~}{P}^{}.`$ (5.9) Suppose that at some neighborhood of the optimal solution $`g_{}`$, $`H^{(t)}(X)`$ is Lipschitz continuous for some $`t`$: $`H^{(t)}(X)H^{(t)}(X^{})LXX^{},`$ (5.10) where $`A`$ is the norm of a matrix $`A`$ as the Euclidian space, $`A^2=\mathrm{tr}(AA^{}).`$ (5.11) We set $`\beta =H^{(t)}(X^t(g_{}))^1.`$ (5.12) There exists a positive real number $`r`$, for which $`H^{(t)}(X^t(g))^1<2\beta \text{for}gB^{(t)}(g_{},r)\stackrel{\mathrm{def}}{=}\left\{g\right\|r>X^t(g)X^t(g_{})\}`$ (5.13) is satisfied. Then it is known that for all $`gB(g_{},\mathrm{min}(r,(2\beta L)^1))`$, $`X^t(C_{t+1})X^t(g_{})\beta LX^t(C_t)X^t(g_{})^2`$ (5.14) and $`X^t(C_{t+1})X^t(g_{}){\displaystyle \frac{1}{2}}X^t(C_t)X^t(g_{})`$ (5.15) are fulfilled. Thus the second order convergence in this norm is shown. Unfortunately, this norm is not invariant and is unnatural. (A natural metric on $`K\backslash G`$ is one which is invariant under the parallel transformation, which is induced by the action of elements in $`K\backslash G`$ from the right-hand side.) But, it suffices in practice. ## 6 Discussions ### 6.1 Nonholonomy? Our method is related to the nonholonomic method by Amari, Chen, and Chichocki\[Amari et al.,1997\]. In essence our method is a Newton approach to the same problem, the optimization without prewhitening. Let us set $`\mathrm{e}^z=\mathrm{e}^x\mathrm{e}^y`$ (6.1) for $`x,y𝔤l(N,R)`$. Then it is obvious that $`z`$ does not necessarily belongs to $`𝔪`$ even if $`x,y𝔪`$(, that is, $`z_{ii}`$’s do not always vanish when $`x_{ii}=y_{ii}=0`$ for $`1iN`$). This may be explained by using the concept of nonholonomy. The degree of freedom in each step, however, equals the dimension of the space $`K\backslash G`$ in our setting. The nonholonomic nature emerges when we go back to $`G=GL(N,𝑹)`$ again. There exist several studies\[M.Takeuchi,1994, S.Helgason,1978, S.Helgason,1962, S.Helgason,1984, T.Akuzawa & M.Wadati,1998\] which deal with cosets like $`K\backslash G`$ or the right coset $`G/K`$ when $`K`$ is a maximal compact subgroup of $`G`$. Unfortunately, what we are studying is the case where $`K`$ is not a maximal compact subgroup of $`G`$. So, for example it is necessary to show whether the coordinate (5.5) is justified or not. As mentioned above, further studies including this justification will appear elsewhere. ### 6.2 Global convergence We should carefully treat first few steps since this method also has a somewhat undesirable global convergent property inherent in the Newton method. Fortunately enough, there exist methods which can handle the earlier stage. For example, the nonholonomic gradient method\[Amari et al.,1997\] may be applicable. Another posiibility is to construct a nonholonomic fixed-point algorithm which uses the kernel method. These methods are suitable for capturing the optimal point which contains components with zero kurtoses. There we must, of course, use the method in Section 3. If it is not necessary to worry about these zero kurtosis components, there is little difference between the two methods described in Section 2 and Section 3. ### 6.3 Conclusions We have constructed a new algorithm for finding a optimal point in a matrix space, where we have introduced a new multiplicative updating method. The algorithm is in essence the Newton method on a coset. So it converges quite rapidly and it can capture the saddle point. Since it does not require prewhitening, it is not necessary to worry about the error resulting from the prewhitening. Indeed, it is possible to increase the kurtosis slightly for data preprocessed by the FastICA\[Hurri et al.,1998\]. ## appendix ## Appendix A proof of (2.31) > For $`BGL(N,F)`$ and $`1i,jN`$, > > $`[T(XY)T\mathrm{cs}(B)]_{i+N(j1)}`$ $`=`$ $`[(XY)T\mathrm{cs}(B)]_{j+N(i1)}`$ (A.1) > $`=X_{ip}Y_{jq}(B^{})_{qp}=(YB^{}X^{})_{ji}.`$ > On the other hand > > $`[(YX)\mathrm{cs}(B)]_{i+N(j1)}`$ $`=Y_{jp}X_{iq}B_{qp}=(YB^{}X^{})_{ji}.`$ (A.2) > This proves the statement since $`\mathrm{cs}`$ is bijective. $`\mathrm{}`$
warning/0002/math-ph0002017.html	ar5iv	text	# Abstract ## Abstract The fourth, missing example of an exactly solvable complex potential with $`𝒫𝒯`$ symmetry $`V(x)=[V(x)]^{}`$ defined on a bent contour and leading, at the real energies, to the Jacobi polynomial wave functions is found in a generalized Hulthén interaction. PACS 03.65.Ge, 03.65.Fd Quantum mechanics is often forced to formulate its predictions numerically. Its exactly solvable models are scarce. Only their subclass in one dimension is broader and, in this sense, privileged and exceptional. It involves the harmonic oscillator and Morse potentials (with bound states expressible in terms of the Laguerre polynomials) as well as several models solvable in terms of the polynomials of Jacobi (cf., e.g., review for more details). Within the framework of an alternative or modified quantum mechanics as proposed recently by Bessis and by Bender et al one works with the $`𝒫𝒯`$ symmetric complex potentials $`V(x)=[V(x)]^{}`$. The real bound state spectrum and even the solvability of several real potentials proves rather unexpectedly preserved after a complexification of this type. For example, within the Laguerre-related exactly solvable sub-family we may find both the complex harmonic oscillator and a complexified Morse interaction . The same high degree of analogy between the real and complex forces is also observed for the regular solvable models based on the use of the Jacobi polynomials . A word of warning against unlimited optimism comes from singular interactions. In particular, for the Laguerre-related and phenomenologically most appealing Coulombic interaction the only available $`𝒫𝒯`$ symmetrization proves merely partially solvable . In the present communication we intend to offer a partial remedy. We shall derive and describe a complete and exact solution for an appropriate $`𝒫𝒯`$ symmetric complexification of the singular phenomenological Hulthén potential which is known to mimick very well the shape of the Coulombic force in the vicinity of its singularity . In the first step let us recollect that in the spirit of the old Liouville’s paper the change of the (real) coordinates (say, $`r\xi `$) in Schrödinger equation $$\left[\frac{d^2}{dr^2}+W(r)\right]\chi (r)=\kappa ^2\chi (r)$$ (1) may sometimes mediate a transition between two different potentials. It is easy to show that once we forget about boundary conditions one simply has to demand the existence of the invertible function $`r=r(\xi )`$ and its few derivatives $`r^{}(\xi ),r^{\prime \prime }(\xi ),\mathrm{}`$ in order to get the explicit correspondence between the two bound state problems, viz., original eq. (1) and the new Schrödinger equation $$\left[\frac{d^2}{d\xi ^2}+V(\xi )\right]\mathrm{\Psi }(\xi )=E\mathrm{\Psi }(\xi )$$ (2) with the wave functions $$\mathrm{\Psi }(\xi )=\chi [r(\xi )]/\sqrt{r^{}(\xi )}$$ (3) generated by a new interaction at new energies, $$V(\xi )E=\left[r^{}(\xi )\right]^2\left\{W[r(\xi )]+\kappa ^2\right\}+\frac{3}{4}\left[\frac{r^{\prime \prime }(\xi )}{r^{}(\xi )}\right]^2\frac{1}{2}\left[\frac{r^{\prime \prime \prime }(\xi )}{r^{}(\xi )}\right].$$ (4) The mapping between the Morse and harmonic oscillators is one of the best known explicit illustrations of this rule since the necessary preservation of the correct physical boundary conditions is very straightforward to check there . An appropriate $`𝒫𝒯`$ symmetrized extension of this equivalence to the complex non-Hermitian cases is easy . In the alternative, Jacobi-polynomial context the Liouvillean changes of variables have been applied systematically to all the Hermitian models (cf. Figure 5.1 in the review or refs. for a more detailed illustration). A similar thorough study is still missing for the $`𝒫𝒯`$ symmetric models within the same subclass. In the present letter we shall try to fill the gap. For the sake of brevity we shall only restrict our attention to the $`𝒫𝒯`$ symmetric initial eq. (1) with the Pöschl-Teller potential studied and solved exactly in our recent preprint , $$W(r)=\frac{\beta ^21/4}{\mathrm{sinh}^2r}\frac{\alpha ^21/4}{\mathrm{cosh}^2r},r=xi\epsilon ,x(\mathrm{},\mathrm{})$$ (5) The normalizable $`𝒫𝒯`$ symmetric solutions $$\chi (r)=\mathrm{sinh}^{\tau \beta +1/2}r\mathrm{cosh}^{\sigma \alpha +1/2}rP_n^{(\tau \beta ,\sigma \alpha )}(\mathrm{cosh}2r)$$ are proportional to the Jacobi polynomials at all the negative energies $`\kappa ^2<0`$ such that $$\kappa =\kappa _n^{(\sigma ,\tau )}=\sigma \alpha \tau \beta 2n1>0.$$ These bound states are numbered by $`n=0,1,\mathrm{},n_{max}^{(\sigma ,\tau )}`$ and by the generalized parities $`\sigma =\pm 1`$ and $`\tau =\pm 1`$. We may note that our initial $`𝒫𝒯`$ symmetric model (1) remains manifestly regular provided only that its constant downward shift of the coordinates $`r=r_{(x)}=xi\epsilon `$ remains constrained to a finite interval, $`\epsilon (0,\pi /2)`$. In this sense our initial model (5) is closely similar to the shifted harmonic oscillator. At the same time, one still misses an analogue of a transition to its Morse-like final partner $`V(\xi )`$ in eq (2). In a key step of its present construction let us first pick up the following specific change of the axis of coordinates, $$\mathrm{sinh}r_{(x)}(\xi )=ie^{i\xi },\xi =viu.$$ (6) The main motivation of such a tentative assignment lies in the related shift and removal of the singularity (sitting at $`r=0`$) to infinity ($`u+\mathrm{}`$). In fact, one cannot proceed sufficiently easily in an opposite direction, i.e., from a choice of a realistic $`V(\xi )`$ to a re-constructed $`r(\xi )`$. This is due to the definition (4) containing the third derivatives and, hence, too complicated to solve. We shall see below that we are quite lucky with our purely trial and error choice of eq. (6). Firstly, we already clearly see that the real line of $`x`$ becomes mapped upon a manifestly $`𝒫𝒯`$ symmetric curve $`\xi =viu`$ in accordance with the compact and invertible trigonometric rules $$\mathrm{sinh}x\mathrm{cos}\epsilon =e^u\mathrm{sin}v,\mathrm{cosh}x\mathrm{sin}\epsilon =e^u\mathrm{cos}v,$$ i.e., in such a way that $$\begin{array}{c}v=\mathrm{arctan}\left(\frac{\mathrm{tanh}x}{\mathrm{tan}\epsilon }\right)=v_{(x)}(v_{(\mathrm{})},v_{(\mathrm{})})(\frac{\pi }{2}+\epsilon ,\frac{\pi }{2}\epsilon ),\\ u=u_{(x)}=\frac{1}{2}\mathrm{ln}\left(\mathrm{sinh}^2x+\mathrm{sin}^2\epsilon \right).\end{array}$$ This relationship is not too different from its Morse-harmonic predecessor studied in ref. . Our present path of $`\xi `$ is a very similar down-bent arch which starts in its left imaginary minus infinity, ends in its right imaginary minus infinity while its top lies at $`x=v=0`$ and $`u=u_{(0)}=\mathrm{ln}1/\mathrm{sin}\epsilon >0`$. The top may move towards the singularity in a way mimicked by the diminishing shift $`\epsilon 0`$. Indeed, although the singularity originally occurred at the finite value $`r0`$, it has now been removed upwards, i.e., in the direction of $`u+\mathrm{}`$. The first consequence of our particular change of variables (6) is that it does not change the asymptotics of the wave functions. As long as $`r^{}(\xi )=i\mathrm{tanh}r(\xi )`$ the transition from eq. (1) to (2) introduces just an inessential phase factor in $`\mathrm{\Psi }(\xi )`$. This implies that the normalizability (at a physical energy) as well as its violations (off the discrete spectrum) are both in a one-to-one correspondence. The explicit relation between the old and new energies and couplings is not too complicated. Patient computations reveal its closed form. With a bit of luck, the solution proves non-numerical. The new form of the potential and of its binding energies is derived by the mere insertion in eq. (4), $$V(\xi )=\frac{A}{(1e^{2i\xi })^2}+\frac{B}{1e^{2i\xi }},E=\kappa ^2.$$ (7) At the imaginary $`\xi `$ and vanishing $`A=0`$ this interaction coincides with the Hulthén potential. In the new formula one has to notice the positivity of the energies. It is extremely interesting since the potential itself is asymptotically vanishing. One may immediately recollect that a similar paradox has already been observed in a few other $`𝒫𝒯`$ symmetric models where even an asymptotic decrease of the potential to minus infinity did not destroy a lower bound of the spectrum. The exact solvability of our modified Hulthén potential is not yet guaranteed at all. A critical point is that the new couplings depend on the old energies and, hence, on the discrete quantum numbers $`n`$, $`\sigma `$ and $`\tau `$ in principle. This could induce an undesirable state-dependence into our new potential. Vice versa, the closed solvability of the constraint which forbids this state-dependence will be equivalent to the solvability at last. Fortunately, a smooth removal of the obstacle is possible by a transfer of the state-dependence (i.e., of the $`n`$, $`\sigma `$ and $`\tau `$dependence) in $$A=A(\alpha )=1\alpha ^2,C(=A+B)=\kappa ^2\beta ^2$$ from $`C`$ to $`\beta `$. To this end, employing the known explicit form of $`\kappa `$ we may re-write $$C=C(\sigma ,\tau ,n)=(\sigma \alpha +2n+1)(\sigma \alpha +2n+1+2\tau \beta ).$$ (8) This formula is linear in $`\tau \beta `$ and, hence, its easy inversion defines the desirable state-dependent quantity $`\beta =\beta (\sigma ,\tau ,n)`$ as an elementary function of the constant $`C`$. The whole new energy spectrum acquires the closed form $$E=E(\sigma ,\tau ,n)=A+B+\frac{1}{4}\left[\sigma \alpha +2n+1\frac{A+B}{\sigma \alpha +2n+1}\right]^2.$$ (9) Our construction is complete. The range of the quantum numbers $`n,\sigma `$ and $`\tau `$ remains the same as above. In the light of our new result we may now split the whole family of the exactly solvable $`𝒫𝒯`$ symmetric models in the two distinct categories. The first one “lives” on the real line and may be represented or illustrated not only by the popular Laguerre-solvable harmonic oscillator but also by our initial Pöschl-Teller Jacobi-solvable force (5). The second category requires a narrow arch-shaped path of integration which all lies confined within a vertical strip. It contains at last both the Laguerre and Jacobi solutions. The former ones may be represented by the complex Morse model of ref. . Our present new Hulthén example offers its Jacobi solvable counterpart. The scheme becomes, in a way, complete. The less formal difference between the two categories may be also sought in their immediate physical relevance. Applications of the former class may be facilitated by a limiting transition which is able to return them back on the usual real line. In contrast, the second category may rather find its most useful place in the methodical considerations concerning, e.g., field theories and parity breaking . Within the quantum mechanics itself an alternative approach to the second category might also parallel studies of the “smoothed” square wells in the non-Hermitian setting. In the conclusion let us recollect that the $`𝒫𝒯`$ symmetry of a Hamiltonian replaces and, in a way, generalizes its usual hermiticity. This is the main reason why there existed a space for a new solvable model among the singular interactions. At the same time, an exactly solvable example with another, intermediate (i.e., hyperbola-shaped) arc of coordinates remains still to be discovered. Indeed, this type of a deformed contour has only been encountered in the quasi-solvable (i.e., partially numerical) model of ref. ) and in the general unsolvable forces studied by several authors by means of the perturbative , numerical and WKB approximation techniques. ## Acknowledgement Partially supported by the grant Nr. A 1048004 of the Grant Agency of the Academy of Sciences of the Czech Republic.
warning/0002/cond-mat0002348.html	ar5iv	text	# Charge and Spin Order in La₂₋ₓ⁢Srₓ⁢NiO₄ (x=0.275 and 1/3) \[ ## Abstract We report polarized and unpolarized neutron diffraction measurements on $`\mathrm{La}_{2\mathrm{x}}\mathrm{Sr}_\mathrm{x}\mathrm{NiO}_4`$ (x=0.275 and 1/3). The data for the spin ordered states are consistent with a collinear spin model in which spins with S=1 are in the NiO<sub>2</sub> plane and are uniformly rotated from the charge and spin stripe direction. The deviation angle is larger for x=1/3 than for x=0.275. Furthermore, for x=1/3 the ordered spins reorient below 50 K, which suggests a further lock-in of the doped holes and their magnetic interactions with the S=1 spins. \] Since the discovery of incommensurate (IC) magnetic fluctuations in superconducting cuprates, magnetism in doped antiferromagnets has attracted much attention. Doping holes into the antiferromagnetic (AFM) insulator $`\mathrm{La}_2\mathrm{CuO}_4`$ rapidly weakens the static magnetic correlations and leads to a metallic, superconducting phase at moderate hole concentrations. The discovery of static IC spin and charge peaks in $`\mathrm{La}_{1.6\mathrm{x}}\mathrm{Nd}_{0.4}\mathrm{Sr}_\mathrm{x}\mathrm{CuO}_4`$ led to a model of charge stripes acting as antiphase domain walls for the intervening AFM regions. $`\mathrm{La}_2\mathrm{NiO}_4`$ is a Mott insulator with $`T_N650`$ K and spins aligned with the shortest orthorhombic lattice direction. Doping of divalent Sr<sup>2+</sup> ions into the trivalent La<sup>3+</sup> sites in $`\mathrm{La}_{2\mathrm{x}}\mathrm{Sr}_\mathrm{x}\mathrm{NiO}_4`$ (LSNO(x)) rapidly suppresses the ordering but the compound remains insulating until $`x1`$ where it becomes metallic. For $`x>0.1`$, IC charge and spin superlattice reflections appear, and can be characterized by $`Q_{spin}=(1\pm ϵ,0,0)`$ and $`Q_{charge}=(2ϵ,0,1)`$ with $`ϵx`$. For $`x1/3`$, there is a commensurability effect which favors $`ϵ=1/3`$. Considering that the static spin structure should be a basis for understanding the nature of spin and charge stripes in the doped nickelates and cuprates, it is surprising that the local spin structure, evolving upon doping, has not been studied. In this letter we report polarized and unpolarized neutron diffraction measurements on LSNO(x=0.275 and x=1/3) and present a collinear spin model to explain our data in which spins are rotated uniformly by an angle $`\theta `$ from the stripe direction. For x=0.275, the angle $`\theta `$ is $`27(7)^o`$ below $`T_N=155(5)`$ K. For x=1/3, $`\theta =40(3)^o`$ for 50 K $`<T<T_N=200(10)`$ K and below 50 K another phase transition occurs involving reorientation of spins to $`\theta =53(3)^o`$. We argue that the spin reorientation is due to a further localization of the doped holes on the lattice, which is consistent with recent resistivity measurements that showed a delocalization of charge stripes by an electic field (although in the lowest temperature phase the charge stripes are more robust.) Our polarized neutron data also provide unequivocal evidence that upon cooling the charge order precedes the spin order in these materials. Charge orders at 200(10) K for x=0.275, and at 240(10) K for x=1/3. The LSNO(x) crystals with x=0.275 and x=1/3 used in this work were grown by the floating-zone method at Kyoto University and at Bell Laboratories, respectively. The crystal structure of both materials remains tetragonal with the I4/mmm symmetry over the range of temperature, 10 K $`<T<`$ 320 K, with $`a_t=3.82(1)`$ Å and $`c=12.64(5)`$ Å at 11 K for x=0.275 and with $`a_t=3.83(1)`$ Å and $`c=12.69(5)`$ Å at 11 K for x=1/3. For convenience, however, we use orthorhombic lattice units with $`a_o=\sqrt{2}a_t`$ which is rotated by 45<sup>o</sup> with respect to the Ni-O bonds in the NiO<sub>2</sub> planes. Neutron scattering measurements were performed on the cold neutron triple axis spectrometer SPINS at NIST. The incident neutron energy was $`E_i=5`$ meV and higher order contamination was eliminated by a cryostat-cooled Be filter before the sample. The incident neutrons were polarized in one spin state by a forward transmission polarizer. The spin state of the scattered neutrons from the sample was then analyzed with a rear flipper and polarizer combination. A collimator was placed right after each polarizer to eliminate the other spin state. The samples were mounted in such a way that the (h0l) reciprocal plane becomes the scattering plane. A guide field was applied vertically along the (010)-axis. Collimations were guide-40-40-open for the LSNO(x=0.275) measurements and guide-40-20-open for the LSNO(x=1/3) measurements. The polarizing efficiency was 0.89(1) for the former and 0.90(1) for the latter. In this geometry, the spin component perpendicular to the wavevector transfer $`\stackrel{}{Q}`$ which is parallel with the guide field together with the nuclear structural component, contribute to non-spin-flip (NSF) channel whereas the spin component normal to $`\stackrel{}{Q}`$ and to the guide field contributes to the spin-flip (SF) channel. Since spins in the nickelates are coplanar in ab-plane, the NSF and SF neutron scattering cross sections for $`Q=(h0l)`$, $`\sigma _{NSF}(Q)`$ and $`\sigma _{SF}(Q)`$ respectively, can be written as $`\sigma _{NSF}(Q)`$ $`=`$ $`\sigma _N(Q)+\sigma _M^b(Q)`$ (1) $`\sigma _{SF}(Q)`$ $`=`$ $`\sigma _M^a(Q)\left(1{\displaystyle \frac{(ha^{})^2}{(ha^{})^2+(lc^{})^2}}\right)`$ (2) where $`\sigma _N`$ is the structural scattering cross section, $`\sigma _M^a`$ and $`\sigma _M^b`$ are the a- and b-component of the magnetic scattering cross section, respectively. Fig. 1. Polarized elastic neutron scattering data along (h05) from LSNO(x=0.275) ((a) and (b)), and from LSNO(x=1/3) ((c) and (d)). Fig. 1 shows the results of polarized neutron diffraction on LSNO(x=0.275 and 1/3). For x=0.275, at 11K the data exhibit two types of superlattice peaks: the first harmonic at (0.7,0,5) separated by $`(ϵ,0,0)`$ from the (105) AFM Bragg reflection in pure $`\mathrm{La}_2\mathrm{NiO}_4`$ and the second harmonic at (0.6,0,5) separated by $`(2ϵ,0,1)`$ from the (004) nuclear Bragg reflection with $`ϵ=0.3x`$. The first harmonic has both NSF and SF components whereas the second harmonic has a NSF component only. These provide direct and unambiguous evidence that the first harmonic is magnetic (spin peak) whereas the second harmonic is non-magnetic (charge peak) in origin, and suggest the formation of charge stripes acting as AFM domain walls (see Fig. 2). As shown in Fig. 3 (a), the charge peak gradually increases below $`T_c`$ and then get enhanced in a weakly first order at $`T_s`$ when spins order. Inplane charge correlation length, $`\xi _c`$, also increases at $`T_s`$. In the charge and spin ordered phase the charge peak is broader than the spin peak, indicating that $`\xi _c`$ is shorter than inplane spin correlation length, $`\xi _s`$: $`\xi _c/\xi _s100`$ Å/300 Å $`=1/3`$ (see the inset of Fig. 3 (a)). It has been argued that disorders in the stripes are mostly due to non-topological elastic deformations along the stripes and the decrease of the correlation lengths is inversely proportional to a power of the periodicity, which can explain why $`\xi _c`$ is smaller than $`\xi _s`$ even though the charge and spin correlations are inexorably connected. For comparison, $`\xi _c/\xi _s=1/4`$ for $`\mathrm{La}_{1.6\mathrm{x}}\mathrm{Nd}_{0.4}\mathrm{Sr}_\mathrm{x}\mathrm{CuO}_4`$. As shown in Fig. 1 (a) and (b), upon heating the (0.7,0,5) and (0.6,0,5) superlattice peaks shift toward $`(\frac{2}{3},0,5)`$ until the (0.7,0,5) peak vanishes at 155(5) K and the (0.6,0,5) peak at 200(10) K. In the intermediate temperature range $`(155K<T<200K)`$, as shown in Fig. 1 (b), the charge order exists without the spin order. For x=1/3, at 11 K the first and second harmonics with $`ϵ=x`$ come together at the same wave vector as shown in Fig. 1 (c). Unlike the case for x=0.275, the position of the superlattice peak for x=1/3 is independent of $`T`$, implying the stability of commensurability. At 220 K, the (2/3,0,5) peak has a NSF component only, indicating that the x=1/3 system also has an intermediate $`T`$ phase with charge order only. Fig. 2. Possible spin structures in a NiO<sub>2</sub> plane for LSNO(x=1/3). (a) Spins are perpendicular to the propagation vector along the (100)-axis. (b) Spins are uniformly rotated by an angle $`\theta `$ from the (010)-axis. Dashed lines indicate the magnetic unit cell. Besides identifying the origins of the scattering, polarized neutron diffraction data also allow us to determine the spin structure in the material. Firstly, it is to be noted that the spin peaks in the spin ordered states of both samples are SF in nature. If spins are aligned along the (010)-direction (see Fig. 2 (a)), $`\sigma _M^a=0`$ and Eq. (1) would yield zero SF scattering. There are two possible scenarios to produce nonzero SF scattering. One possibility is two magnetic domains with a propagation vector along (100)-axis: in one domain spins are perpendicular to and in the other parallel to the propagation vector. The population of the domains must be unequal to explain our data, which seems unlikely. The other possibility is one single magnetic domain with $`\sigma _M^a0`$. For quantitative studies of the spin structure, we have studied the $`T`$-dependence of the NSF and SF scattering at various superlattice reflections from both samples. Fig. 3 (b) shows some of the results for x=1/3. Upon cooling at 240 K NSF scattering at the (2/3,0,5) reflection develops without SF scattering, indicating charge ordering. Below around 200 K, SF as well as NSF scattering increases due to spin ordering in the usual second order fashion. At 50 K, however, the trend changes: the NSF scattering decreases whereas the SF increases. This behavior can be more clearly seen in the ratio of SF to NSF scattering shown in Fig. 3 (c). The sharp increase in the ratios below 50K indicates that the phase transition at 50 K involves reorientation of spins. For x=0.275, the SF to NSF ratios of the spin peak in the spin-ordered phase were smaller than those for the corresponding reflections in the x=1/3 system and no spin reorientation transition was observed down to 10K. Fig. 3. (a) T-dependence of unpolarized elastic neutron scattering intensities of the charge peak and the spin peak. (b) NSF (open symbols) and SF (filled symbols) scattering intensities at various superlattice reflections from LSNO(x=1/3) as a function of T. (c) T-dependence of the ratio of SF to NSF scattering intensities after backgrounds were determined above the transition temperature and subtracted. Correction for finite polarizing efficiency was also made. The symbols in Fig. 4 summarize the measured $`\sigma _{SF}/\sigma _{NSF}`$ for LSNO(x=0.275) and for the two spin-ordered phases in LSNO(x=1/3). The data show a general trend: for a given value of h in each phase, $`\sigma _{SF}/\sigma _{NSF}`$ increases as $`l`$ increases with an exception at the (2/3,0,5) reflection for x=1/3. The deviation of the (2/3,0,5) reflection is due to the fact that the (2/3,0,5) peak has a charge component as well as spin component. Our unpolarized neutron scattering data along (0.6,0,$`0l5)`$ on the x=0.275 sample have the strongest intensity at $`l=5`$ but negligible at other $`l`$, indicating the structure factor due to charge ordering is strong at $`(0.6,0,5)`$ but not at other $`l`$. We expect the same holds for x=1/3. Now consider the collinear spin structure shown in Fig. 2 (b) in which spins are rotated by $`\theta `$ about the c-axis. Since the spins have a b-component, SF scattering would be non-zero and the ratio of SF to NSF scattering would be determined by the angle $`\theta `$. The lines in Fig. 4 are the calculated $`\sigma _{SF}/\sigma _{NSF}`$ for the spin structure shown in Fig. 2 (b): $$\frac{\sigma _{SF}}{\sigma _{NSF}}=\mathrm{tan}^2\theta \left(1\frac{(ha^{})^2}{(ha^{})^2+(lc^{})^2}\right).$$ (3) Fig. 4. Experimental (symbols) and the model calculated (lines) ratios of SF to NSF scattering cross sections as a function of the rotation angle $`\theta `$. From the comparison to the data, we conclude that $`\theta =27^o\pm 7^o`$ for x=0.275, $`40^o\pm 3^o`$ for phase II and $`52.5^o\pm 2.5^o`$ for phase III in the x=1/3 sample. We can also estimate the charge contribution in the (2/3,0,5) reflection to be $`\frac{\sigma _N}{\sigma _M^b}=0.20(6)`$. For x=0.275 (see Fig. 1(a)) $`\frac{\sigma _N}{\sigma _M^b}=\frac{\sigma _{NSF}(0.6,0,5)}{\sigma _{NSF}(0.7,0,5)}=0.12(1)`$. Unpolarized neutron diffraction data along the $`l`$-direction also contain information on spin structure. As shown in Fig. 5, the peaks are much sharper in $`l`$ for x=0.275 than for x=1/3, which indicates that the correlations in x=0.275 are nearly three-dimensional whereas those in x=1/3 are quasi-two-dimensional. Another obvious difference between them is that the even $`l`$ to old $`l`$ peak intensity ratio, $`\frac{I(evenl)}{I(oddl)}`$, is much larger for x=0.275 than for x=1/3. There are other subtle features to be noted. For x=0.275, the odd $`l`$ peak weakens as $`l`$ increases whereas the even $`l`$ peak is strongest at $`l=2`$. For x=1/3, at 180K (phase II) the odd $`l`$ peak also weakens as $`l`$ increases. However at 15K (phase III) the $`l=3`$ peak becomes strongest. If the spin structure of Fig. 2 (a) is displaced by $`(\frac{a}{2}n,\frac{1}{2},\frac{1}{2})`$ in neighboring NiO<sub>2</sub> planes, the magnetic neutron scattering cross section would become $`\sigma (Q)\|F(Q)\|^2(1+()^l\mathrm{cos}n\pi h)`$ where $`F`$ is the magnetic form factor of the Ni<sup>2+</sup>. Therefore, for a given $`h`$, the $`l`$-dependence of odd or even $`l`$ peaks would just follow $`\|F(Q)\|^2`$. This $`\sigma (Q)`$ cannot explain those subtle features. On the other hand, the spin structure of Fig. 2 (b) would give $`\sigma (Q)`$ $``$ $`\|F(Q)\|^2(1+()^l\mathrm{cos}n\pi h)`$ (5) $`\times \left(1\mathrm{sin}^2\theta {\displaystyle \frac{(ha^{})^2}{(ha^{})^2+(lc^{})^2}}\right).`$ With the $`\theta `$’s which are obtained from the polarization analysis, Eq. 3 reproduces the $`l`$-dependence remarkably well for the case of x=0.275 with three-dimensional magnetic correlations, as shown as shaded bars in Fig. 5. Fig. 5. $`l`$-dependence of (a) (0.7,0,$`l`$) scan form LSNO(x=0.275) at 11 K, and (2/3,0,$`l`$) scan from LSNO(x=1/3) (b) at 180 K and (c) 15 K. The shaded bars are obtained by the model which is described in the text. In the case of x=1/3 where the magnetic order is quasi-two dimensional, the long tails of the peaks make it difficult to extract the integrated intensities to compare with the model. Nevertheless, the model explains the subtle features such as the (2/3,0,3) peak being considerably stronger than the other reflections in phase III. Spin reorientations upon cooling have been observed in some isostructural materials such as pure $`\mathrm{La}_2\mathrm{CoO}_4`$ and $`\mathrm{Nd}_2\mathrm{CuO}_4`$. Those reorientations involve spin flip (by $`\pi /2`$ or $`\pi `$) which are due to either a structural phase transition (in $`\mathrm{La}_2\mathrm{CoO}_4`$) or strong coupling between the rare-earth moments and $`Cu^{2+}`$ moments (in $`\mathrm{Nd}_2\mathrm{CuO}_4`$). For LSNO(x=1/3), however, there is no structural phase transition down to 10 K within our experimental uncertainty. What causes then the reorientation of spins in LSNO(x=1/3) below 50 K ? The answer might be a further localization of the holes on the lattice and their interactions with the surrounding S$`=`$1 $`Ni^{2+}`$ moments. Among the two charge ordered states, below and above 50 K, the holes should be more strongly pinned on the lattice in the lower $`T`$ phase. The further localization of the holes might either induce a subtle and yet undetected local crystal distortion around the surrounding S$`=`$1 spins or magnetically order themselves, either of which would in turn reorient the S$`=`$1 spins. This idea is consistent with recent resistivity measurements. In the resistivity measurements, a moderate inplane electric field ($`V_{cir}100`$ V) decreases the resistivity by up to 5 orders of magnitude in phase II. However, although the charge-ordered state in phase III becomes suppressed by the electric field, phase III survives until $`V_{cir}900`$ V, which indicates the stronger localization of the holes in phase III than in phase II. Understanding the microscopic origin of the further localization in LSNO(x=1/3) and its effects, such as the spin reorientation, as well as the evolution of spin structure upon doping in LSNO(x) requires more theoretical and experimental studies. The authors thank G. Aeppli, J.M. Tranquada, G. Shirane, Y.S. Lee, and S. Wakimoto for helpful discussions. Work at SPINS is based upon activities supported by the National Science Foundation under Agreement No. DMR-9423101.
warning/0002/cond-mat0002428.html	ar5iv	text	# High-field AFMR in single-crystalline La0.95Sr0.05MnO3: Experimental evidence for the existence of a canted magnetic structure ## Abstract High-field antiferromagnetic-resonance (AFMR) spectra were obtained in the frequency range 60 GHz $`<\nu <`$ 700 GHz and for magnetic fields up to 8 T in twin-free single crystals of La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub>. At low temperatures two antiferromagnetic modes were detected, which reveal different excitation conditions and magnetic field dependencies. No splitting of these modes was observed for any orientation of the static magnetic field excluding the phase-separation scenario for this composition. Instead, the full data set including the anisotropic magnetization can be well described using a two-sublattice model of a canted antiferromagnetic structure. The idea of phase separation in the manganite perovskites (R<sub>1-x</sub>M<sub>x</sub>MnO<sub>3</sub>, R=La, Pr …, M=Ca, Sr …) is one of the most controversially discussed topics concerning the electronic properties of these compounds. After the pioneering works of Jonker and van Santen and of Wollan and Koehler, de Gennes developed a model in which the purely antiferromagnetic and insulating LaMnO<sub>3</sub> on increasing doping passes through a canted (CAF) ground state and arrives at a purely ferromagnetic and metallic state at high doping level ($`x0.2`$). This phase diagram was calculated using competing superexchange (SE) and double exchange (DE) interactions. However, in recent years a number of theoretical models predicted that the CAF structure becomes unstable against electronic phase separation into ferromagnetic (FM) and antiferromagnetic (AFM) regions. As discussed by Yunoki et al., the tendency to the phase separation seems to be an intrinsic property of the double-exchange model. A number of experimental data including neutron scattering and NMR pointed toward the existence of the phase separation in different types of manganites. For discussion of the recent results see Refs.. It has to be pointed out, that the experimental observation of electronic phase separation is rather difficult. Already more than forty years ago, Wollan and Koehler stated that, on the basis of neutron diffraction experiments, it is impossible to decide whether the structure of the doped manganites is homogeneously canted or inhomogeneously mixed FM and AFM. Only a few experimental methods can distinguish between inhomogeneous and homogeneous magnetic phases because the technique has to be sensitive to the local magnetic structure of the sample. In addition, a sample quality appears to be of major importance for these experiments. Antiferromagnetic resonance (AFMR) seems to be an excellent tool for the solution of the phase-separation problem. The main parameters of the resonance lines, like position, excitation conditions, behavior in magnetic field etc., sensitively depend on the local environment of the magnetic moments. Recently, using this method, we have investigated the concentration dependence of AFMR-lines in low-doped La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> without external magnetic field. The results were explained within the frame of a two-sublattice model, which strongly supported the existence of a canted magnetic structure. Most La<sub>1-x</sub>Sr<sub>x</sub>MnO<sub>3</sub> crystals of this series were twinned. However, the samples with 5% Sr concentration were identified as untwinned single crystals. This fact allowed the unambiguous determination of the excitation conditions of AFMR lines and to carry out detailed investigations in static magnetic fields. In this paper we present, in addition to the results of the magnetic-field experiments, anisotropic magnetization curves and compare the observed data with the predictions of a two-sublattice model. The possible explanations within phase separation models are also discussed. La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub> single crystals were grown by a floating zone method with radiation heating. X-ray powder diffraction measurements showed that the crystals were single-phase. Four-circle X-ray analysis showed the twin-free structure of the crystal. The temperature dependence of the dc-resistivity of these samples has been published previously and agrees well with literature data. Plane-parallel plates of size approximately 8$`\times `$8$`\times `$1 mm<sup>3</sup> were used for optical measurements. The magnetic measurements were carried out on small pieces of the same crystals. The magnetization curves of La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub> were measured using a SQUID magnetometer in fields up to 6.5 T. The transmission spectra in the frequency range 40 GHz $`\nu `$ 700 GHz were recorded using a quasioptical technique utilizing backward-wave oscillators as coherent light sources. Combining this method with a superconducting split-coil magnet equipped with optical windows allows to carry out transmission experiments in fields up to 8T. The data were obtained in the frequency-sweep mode at constant magnetic field. However, in some cases field-sweep measurements were performed because this procedure enhances the accuracy of determination of the resonance frequency. The frequency-dependent transmission spectra were analyzed using the Fresnel optical formulas for a transmission coefficient of a plane-parallel plate. The relative transparency of the sample in the frequency range investigated resulted in the observation of interference patterns in the spectra. The observation of these interferences allowed the calculation of the optical parameters of the sample without measuring the phase shift of the transmitted signal. The dispersion of the magnetic permeability was taken into account assuming a harmonic oscillator model for the complex magnetic permeability: $$\mu ^{}(\nu )=\mu _1+i\mu _2=1+\mathrm{\Delta }\mu \nu _0^2/(\nu _0^2\nu ^2+i\nu g)$$ (1) where $`\nu _0`$, $`\mathrm{\Delta }\mu `$ and $`g`$ are eigenfrequency, mode strength and width of the resonance respectively. The dielectric parameters of the sample $`(n^{}=n+ik)`$ were assumed to behave regular in the vicinity of the resonance frequency. Hence, the frequency-sweep measurements allowed to obtain absolute values of the parameters of AFMR lines. Fig. 1 shows the low-temperature magnetization curves of single-crystalline La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub> for different orientations of the magnetic field. The magnetization along the c-axis shows a spontaneous magnetization and therefore is identified as the direction of the weak ferromagnetic moment. The magnetization along the crystallographic b-axis reveals the weakest field dependence and thus resembles the data of a simple antiferromagnet along the easy (antiferromagnetic) axis. In order to understand the magnetization measurement and the submillimeter spectra quantitatively, a two-sublattice model has been adopted. The two-sublattice model is a widely used approximation to describe the properties of magnetically ordered materials . It was originally applied to the spin-wave spectrum of manganites by de Gennes . For a realistic description, additional contributions to the free energy, i.e. the single ion anisotropy $`D_x\mathrm{\Sigma }_iS_{xi}^2+D_z\mathrm{\Sigma }_iS_{zi}^2`$ and the Dzyaloshinsky-Moria (D-M) antisymmetric exchange interactions $`\mathrm{\Sigma }_{i,j}𝐝_{ij}[S_iS_j]`$ have to be taken into account. In the classical approximation the free energy at T=0 is given by: $$\begin{array}{c}F(𝐦,𝐥)=\frac{1}{2}A𝐦^2B\|𝐦\|+\frac{1}{2}K_x(m_x^2+l_x^2)+\hfill \\ +\frac{1}{2}K_z(m_z^2+l_z^2)d(m_zl_ym_yl_z)M_0\mathrm{𝐦𝐇}\hfill \end{array}$$ (2) In Eq. (2) the (x, y, z) axes of the coordinate system are directed along the crystallographic axes (a, b, c) of the sample (a=5.547Å, b=5.666Å, c=7.725Å). The first and the second terms of Eq. (2) describe antiferromagnetic and ferromagnetic (double) exchange, the third and fourth terms give the single ion anisotropy, the fifth term describes the D-M exchange while the last term takes into account effects of an external magnetic field. In Eq.(2), $`𝐦`$ and $`𝐥`$ are dimensionless ferro- and antiferromagnetic vectors, which are defined as $`𝐦=(𝐌_1+𝐌_2)/2M_0`$, $`𝐥=(𝐌_1𝐌_2)/2M_0`$ and satisfy the conditions $`\mathrm{𝐦𝐥}=0`$, $`𝐦^2+𝐥^2=1`$ since the sublattices $`𝐌_1`$ and $`𝐌_2`$ are assumed to be saturated at $`T=0`$. The parameter $`B`$ describes the DE interaction. $`K_{x,z}>0`$ are anisotropy constants stabilizing the $`A_yF_z`$ configuration in pure LaMnO<sub>3</sub>. $`d`$ is the interlayer antisymmetric exchange constant. $`M_0=0.95M_0(Mn^{3+})+0.05M_0(Mn^{4+})=3.95\mu _\text{B}`$, is the saturation magnetization of the sublattices. The equilibrium arrangement of the sublattices has been obtained minimizing the free energy given by Eq. (2). The frequencies of the resonance modes were calculated in the limit of small perturbations from the equations of motion $`d𝐌_i/dt=\gamma [𝐌_i\times F/𝐌_i],`$ $`(i=1,2)`$, where $`\gamma `$ is the gyromagnetic ratio. The solutions of Eq.(2) for the H$`\|`$$`\|`$c were published previously . The full set of solution for the field orientation along the a and b-axes is too lengthy and will be published in full form elsewhere. However, in order to understand the experimental data qualitatively some approximations can easily be made. Assuming $`(B,d,K_x,K_z,M_0H)A`$, the approximate solution for the magnetization can be written as: $$M_xM_0m_x=\chi _{}H_x(B+d)/d\text{ , }𝐇=(H_x,0,0)$$ (3) $$M_yM_0m_y=\chi _{rot}H_y\text{ , }𝐇=(0,H_y,0)$$ (4) $$M_zM_0m_z=M_z^0+\chi _{}H_z\text{ , }𝐇=(0,0,H_z)$$ (5) where $`M_z^0M_s=M_0(B+d)/(A+K_z)`$ is the spontaneous magnetic moment along the c-axis, $`\chi _{}=M_0^2/(A+K_z)`$ and $`\chi _{rot}=M_s^2/K_z`$ are the transverse and rotational susceptibilities respectively. The analysis of Eq. (5) shows that the z-axis exhibits weak ferromagnetism as the magnetization is nonzero in the absence of an external magnetic field. The magnetization along the y-axis (Eq. 4) is determined by the small rotational susceptibility and disappears in the pure antiferromagnetic case $`(B=d=0)`$. The low-field susceptibility along the x-direction (Eq. 3) is enhanced compared to the z-axis by the factor $`(B+d)/d`$. Qualitatively similar behavior of the magnetization is observed in Fig. 1. The solid lines in Fig. 1 were calculated using the exact expressions based on Eq. (2) and describe the experimental data reasonably well. A small static moment along the y-axis appears to be strongly angle dependent and possibly is due to some residual influence of the spontaneous moment along the z-axis. The absolute values of the parameters of the model were obtained by simultaneously fitting the magnetization curves and the values of the resonance frequencies in the absence of magnetic field. Despite the relatively large number of parameters in Eq. (2) $`(A,B,d,K_x,K_z)`$, the requirement of a simultaneous fit allows the unambiguous determination of the parameters: $`A=4.6710^7`$ erg/g, $`B=7.410^6`$ erg/g, $`K_z=3.3310^6`$ erg/g, $`K_x=3.4210^6`$ erg/g, $`d=2.110^6`$ erg/g, and $`M_0=92.14`$ emu/g. From the values of $`A`$ and $`K_x`$ the interlayer exchange ($`J_2=0.37meV`$) and the single-ion anisotropy ($`C=0.11meV`$) constants can be calculated which are in good agreement with neutron scattering data for La<sub>0.95</sub>Ca<sub>0.05</sub>MnO<sub>3</sub> and for La<sub>0.95</sub>Sr<sub>0.06</sub>MnO<sub>3</sub>. Fig. 2 shows the transmission spectra of the La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub> single crystal at low temperatures. The solid lines were calculated using the Fresnel equations and Eq.(1) as described above. The parameters of the magnetic mode were obtained by fitting the transmission spectra. In addition, the frequency position and the appearance of the modes were also examined using field sweeps at fixed frequencies. From the data shown in Fig. 2 two peculiarities of the observed AFMR modes immediately become clear: i) the observed lines have unique excitation conditions: $`\stackrel{~}{h}`$ c -axis for the high-frequency mode and $`\stackrel{~}{h}`$ b -axis for the low-frequency mode, and ii) no splitting of the AFMR lines is observed in finite magnetic fields for any geometry of the experiment. Both conclusions are characteristic properties of a canted antiferromagnetic structure and follow naturally from the solution of Eq. (2). The magnetic field dependencies of the resonance frequencies of both AFMR lines are shown in Fig. 3. The solid lines in Fig. 3 were calculated on the basis of the two-lattice model, discussed above. However, the parameters of the model were already fixed by fitting the magnetization curves and absolute values of the AFMR frequencies in the absence of magnetic field. Having this in mind, the theoretical curves describe the experimental data reasonably well. The most important feature of Fig. 3 is the softening of the FM-mode for B$``$b. This softening represents a common property of magnetic resonance in antiferromagnets and is followed by the field-induced rearrangement of the magnetic structure (spin-flop) at a critical value of magnetic field. The softening of the FM-mode at low fields is in good agreement with the model calculations. However, the behavior for higher fields (B$``$7T) significantly deviates from the model predictions. These deviations are most probably due to the extreme sensitivity of the data with respect to the exact orientation of the static magnetic field and the neglect of the higher-order terms in Eq. (2). The angular dependence of a critical behavior in a canted antiferromagnet has been calculated in details by Hagedorn and Gyorgy . These calculations show that already a misalignment of the magnetic field as low as one degree strongly suppress the softening of the FM-line in the vicinity of the critical field. Most probably similar effects explain the deviations observed in Fig. 3. Within the presented model it is also possible to calculate the absolute intensities of the AFMR modes. In the absence of static field the solution of Eq. (2) gives $`\mathrm{\Delta }\mu _{xx}=0.0080`$ and $`\mathrm{\Delta }\mu _{zz}=0.0136`$, using the parameters obtained above. These values are in good agreement with the experimental values $`\mathrm{\Delta }\mu _{xx}=0.012\pm 0.003`$ and $`\mathrm{\Delta }\mu _{zz}=0.0120\pm 0.0010`$. Finally, we discuss the possible explanation of the above-presented data within the concept of phase separation in the ferromagnetic droplets in an antiferromagnetic matrix. Already the magnetization data (Fig. 1) impose a set of constraints on the possible configuration of the phases. E.g. ferromagnetic moments have to be parallel to the c-axis and the b-axis has to be the antiferromagnetic easy axis. However, the most important consequences of a possible electronic phase separation follow for the properties of the magnetic resonances: \- the antiferromagnetic phase would reveal an AFMR mode with a resonance frequency similar to the frequencies in pure LaMnO<sub>3</sub> ($`\nu `$18 cm<sup>-1</sup>). \- this AFMR mode should split into two modes in the presence of magnetic field, as was observed by Mitsudo et al. in pure LaMnO<sub>3</sub>. \- a ferromagnetic line arising from the ferromagnetic droplets has to be observed. The frequency of this mode is expected at substantially lower frequencies ($`\nu 10GHz`$) as observed by Lofland et al. in La<sub>0.9</sub>Sr<sub>0.1</sub>MnO<sub>3</sub>. The resonance frequency of this mode should increase roughly linear with external magnetic field up to frequencies 150-250 GHz for B=7T. None of these properties could be detected in the present experiment. Instead, the observed picture can be well described using the canted magnetic structure. In conclusion, twin-free single crystals of La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub> were grown by the floating-zone method. The low-temperature magnetization was measured along the principal crystallographic directions. High field AFMR spectra of this sample were investigated in the frequency range 60-700 GHz and for magnetic fields up to B=8T. Two AFMR lines having different excitation conditions were detected at low temperatures. The softening of the low-frequency mode was observed for the orientation of the static field B$`\|`$$`\|`$b and is explained by approaching to the spin-flop transition. The full data set, obtained for the La<sub>0.95</sub>Sr<sub>0.05</sub>MnO<sub>3</sub> single crystal, can be easily explained as arising from a canted magnetic structure and clearly contradicts the concept of phase separation into ferro- and antiferromagnetic regions. This work was supported in part by BMBF (13N6917/0 - EKM), by DFG (Pi 372/1-1), by RFBR (99-02-16849), and by INTAS (97-30850).
warning/0002/cond-mat0002207.html	ar5iv	text	# Dynamical Origin of Decoherence in Clasically Chaotic Systems ## Abstract The decay of overlap between a wave packet evolved with a Hamiltonian $``$ and the same state evolved with $`+\mathrm{\Sigma }`$ serves as a measure of the decoherence time$`\tau _\varphi `$. Recent experimental and analytical evidence on classically chaotic systems suggest that, under certain conditions, $`\tau _\varphi `$ depends on $``$ but not on $`\mathrm{\Sigma }`$. By solving numerically a Hamiltonian model we find evidence of that property provided that the system shows a Wigner-Dyson spectrum (which defines quantum chaos) and the perturbation exceeds a crytical value defined by the parametric correlations of the spectra. , and The existence of chaos in classical mechanics is manifested in the evolution of a state as an extreme sensitivity to the initial conditions. Quantum mechanics, on the opposite, does not show this sensitivity.. This has raised several problems in a dynamical definition of quantum chaos. In particular, numerical and experimental studies show that time reversal can be achieved with great accuracy. Therefore, the search for a quantum definition of chaos, lead to investigate the spectral properties of quantum systems whose classical equivalent is chaotic. Quantum chaos appears as the regime in which the properties of the eigenstates follow the predictions of the Random Matrix Theory (RMT). In particular the normalized spacing between energy levels $`s=(\epsilon _{i+1}\epsilon _i)/\mathrm{\Delta }\epsilon `$, with $`\mathrm{\Delta }\epsilon `$ the mean level spacing, should have a probability distribution given by the Wigner Dyson distribution $`P_{WD}^O(s)=(\pi s/2)\mathrm{exp}(\pi s^2/4)`$ for an orthogonal ensemble and $`P_{WD}^U(s)=(32s^2/\pi ^2)\mathrm{exp}(4s^2/\pi )`$ for the unitary ensemble. An infinite set of interacting spins is an example of a many-body system which is chaotic in its classical version (lattice gas) and hence it is expected to present the quantum signatures of chaos in the spectrum. The dynamics of this particular system can be studied by Nuclear Magnetic Resonance (NMR). Surprisingly, the “diffusive” dynamics of a local excitation $`\mathrm{exp}[\mathrm{i}t_R]\|0>`$ can be reversed , generating a Polarization Echo $`M`$ at time 2t$`_R.`$ In this case $``$ is the many-body Hamiltonian of a network of spins with dipolar interaction. To accomplish this, the transformation $`(+\mathrm{\Sigma })`$ at time $`t_R`$ is performed with standard NMR techniques . This transformation is possible due to the anisotropic nature of the dipolar interaction. The perturbation $`\mathrm{\Sigma }`$ is a non-invertible component of the Hamiltonian. In some systems, the only contribution to $`\mathrm{\Sigma }`$ is proportional to the inverse of the radio frequency power and hence can be made arbitrarily small. In one body systems the Polarization Echo (i.e. magnitude of the excitation recovered at time $`2t_R)`$ can then be written exactly as: $$M(t)=\|0\left\|\mathrm{exp}[\mathrm{i}(+\mathrm{\Sigma })t/\mathrm{}]\mathrm{exp}[\mathrm{i}t/\mathrm{}]\right\|0\|^2,$$ (1) where $`\|0`$ is the initial wave function, $``$ the unperturbed Hamiltonian and $`\mathrm{\Sigma }`$ the perturbation which can be associated with an environmental disturbance. Then the magnitude of experimental interest is the overlap between the same initial wave function evolved with the two different Hamiltonians, $``$ and $`(+\mathrm{\Sigma })`$. We should note that the second evolution can be seen as an “imperfect time reversal” of the wave function!. Consistently, more than 10 years ago Peres proposed that dynamical signatures of quantum chaos should be searched on the sensitivity to perturbations in the Hamiltonian. Actually, for a classically chaotic system a perturbation in the initial conditions is equivalent to a perturbation in the Hamiltonian. The experiments show that $`M`$ decays rapidly with $`t_R`$ with Gaussian law indicating a progressive failure in rebuilding the original state. We can define a decoherence time $`\tau _\varphi `$ from this failure, as the width of the Gaussian. This is found to be roughly independent of $`\mathrm{\Sigma }`$ and it extrapolates to a finite value when $`\mathrm{\Sigma }0.`$ Using a semiclassical one-body analytical approach in classically chaotic systems characterized by a Lyapunov exponent $`\lambda `$, Pastawski and Jalabert have shown that there is a regime where the attenuation of $`M`$ it is independent of the perturbation $`\mathrm{\Sigma }`$ and becoming $`1/\tau _\varphi =\lambda .`$ This non-perturbative result is valid for long times and as long as $`\mathrm{\Sigma }`$ does not change the Hamiltonian nature. Our general aim is to find numerical evidence of this regime where $`M(t)`$ is independent of $`\mathrm{\Sigma }`$ considering the simplest Hamiltonian that could model spin diffusion. In this work, we study one-body Hamiltonian in quasi-1D systems with $`N`$ states which we called the Star’s necklace model. More specifically, we use a tight binding model of a ring-shaped lattice with on-site disorder, hopping matrix elements $`V`$, and a magnetic flux $`\mathrm{\Phi }`$ perpendicular to the plane of the ring (see inset in Fig.1). Let us discuss the general features through one representative class, each star has $`20`$ sites and there are $`L=35`$ beads in the necklace which makes $`N=700`$. In our case, perturbation acts only between two star beads: $`\mathrm{\Sigma }(\delta \mathrm{\Phi })=V\mathrm{exp}[\mathrm{i2}\pi \mathrm{\Phi }/\mathrm{\Phi }_o](\mathrm{exp}[\mathrm{i2}\pi \delta \mathrm{\Phi }/\mathrm{\Phi }_o]1)\|1L\|+c.c.`$ Bra and ket states contain orbitals in the star. The localized wave packet with energy $`0\left\|\right\|00`$ moves along the string and contains only $`3N/4`$ states. Anderson disorder is $`W=3V`$ which gives a $`\tau _{\mathrm{imp}}3\mathrm{}/V`$ and $`\mathrm{\Phi }=0.1\mathrm{\Phi }_0.`$ We verify that for $`\delta \mathrm{\Phi }=0`$ the dynamics of the system follows a diffusive law, and that for all $`\delta \mathrm{\Phi }`$ the statistics of the eigenvalues correspond to those predicted by RMT (see Fig.1). We interpret these fact as a signature of chaos. However, when studying the parametric correlations of the energy spectrum as a function of $`\delta \mathrm{\Phi }`$ the correlation function is definite positive. A critical value $`\delta \mathrm{\Phi }_c0.1`$ $`\mathrm{\Phi }_0`$ can be extrapolated from the strong decay. For small $`t`$ the decay of $`M(t)`$ is Gaussian like with a characteristic time scaling with $`\mathrm{\Sigma }`$exactly as could be predicted by perturbation theory. After a certain time of the order of the collision time $`\tau _{\mathrm{imp}}`$ it becomes a stretched exponential with a characteristic time $`\tau _\varphi `$ independent of $`\mathrm{\Sigma },`$ $$M(t)\mathrm{exp}(t/\tau _\varphi )^\nu +M_{\mathrm{}},$$ (2) with $`\nu 0.87.`$ and $`\tau _\varphi 18\mathrm{}/V`$ (see Fig. 2). The constant $`M_{\mathrm{}}`$ arises from finite size effects. Nonetheless, if the perturbation does not exceed the critical threshold consistent with that of the correlation function, $`M_{\mathrm{}}`$ increases with decreasing $`\mathrm{\Sigma }.`$ On the other hand, if the perturbation is large $`\delta \mathrm{\Phi }>\delta \mathrm{\Phi }_c,`$ one gets $`M_{\mathrm{}}1/N.`$ (see inset of Fig. 2). The fact that $`\mathrm{\Sigma }_c`$ goes to zero when $`N`$ goes to infinity, together with results of Fig. 2 for a finite system, could be a signature of the existence of a nontrivial thermodynamic limit $`\underset{\mathrm{\Sigma }0}{lim}\underset{N\mathrm{}}{lim}M_{\mathrm{}}=0`$ different from the non-thermodynamic one $`\underset{N\mathrm{}}{lim}\underset{\mathrm{\Sigma }0}{lim}M_{\mathrm{}}=1`$. Preliminary results on the variation of $`\tau _\varphi `$ with the amount of disorder and system size are consistent with this hypothesis. In order to present in a graphical way the physical phenomena of decoherence, we calculated the weight of the wave functions evolved with the two different Hamiltonians and the overlap between them as a function of the layer. The results are shown in Fig. 3. It can be seen that the evolution of the probabilities described by the perturbed and unperturbed Hamiltonians are not significantly different, i.e. they would give the same coarse grained values. However, the perturbation causes spatial fluctuations in the phases which produce an integral overlap that decays to zero. To sum up, our numerical calculations of the polarization echo $`M`$ in a simple chaotic model indicate that there is a regime of the perturbation where the decoherence time depends only on the perturbation. This basic feature is also found in experiments and in other theoretical models. Some differences in the details remain, such as the value in the exponent $`\nu .`$ According to preliminary numerical evidence $`\nu `$ could be related to the particular topology induced by the matrix elements in the Hamiltonian model. Different systems such as disordered cylinders, maximally connected Hamiltonians (RMT) and chaotic stadiums should be studied in order to characterize this dependence.
warning/0002/cond-mat0002254.html	ar5iv	text	# Ferromagnetism and Canted Spin Phase in AlAs/GaMnAs Single Quantum Wells: Monte Carlo Simulation ## I Introduction During the last decade, due to the advances in the control of materials growth, and also in the techniques of characterization, a new interest arose in the study of the magnetic order in layered materials. This area is not restricted to magnetism in metals, but it also includes the study of magnetic semiconductor pseudo-binary alloys like A<sub>1-x</sub>M<sub>x</sub>B, where M stands for a magnetic ion. These alloys are called Diluted Magnetic Semiconductors (DMS). Recently some groups succeeded in producing homogeneous samples of Ga<sub>1-x</sub>Mn<sub>x</sub>As alloys with $`x`$ up to $`7\%`$ using low temperature ($`200300^o`$ C) MBE techniques. Mn is a transition metal having its $`3d`$ level partially filled with five electrons, in such a way that it carries a magnetic moment of $`5\mathrm{}/2`$, according to the Hund’s rule. In the insulating phase, as in II-VI DMS, two Mn<sup>2+</sup> ions occupying the nearest neighbor positions are assumed to interact with each other via a super-exchange mechanism, resulting in an anti-ferromagnetic ordering of their magnetic moments. In the fcc alloys, these interactions are frustrated, establishing the possibility of settling a spin-glass phase at low temperature. A double exchange mechanism which might stabilize a ferromagnetic coupling between the Mn ions in III-V DMS has been suggested by Akai, but has not been confirmed in the EPR experiments performed by Szczytko et al, who did not observe trace of neutral Mn, concluding that the double exchange mechanism is not effective. The possibility of having a DMS based on GaAs opens a wide range of potential applications such as integrated magnetooptoelectronic devices. Besides its practical importance, this kind of DMS introduces an interesting problem: an Mn impurity in GaAs is an acceptor (it binds one hole), and at the same time it carries a localized magnetic moment. In the Ga<sub>1-x</sub>Mn<sub>x</sub>As alloy Mn is, in fact, a strong $`p`$ dopant, the free hole concentration reaching even $`10^{2021}cm^3`$. At small Mn concentrations, the alloy is a paramagnetic insulator. As $`x`$ increases it becomes ferromagnetic, going through a non-metal-to-metal transition for higher concentrations, and keeping its ferromagnetic phase. For $`x`$ above $`7\%`$, the alloy becomes a ferromagnetic insulator. In the metallic phase, depending on the value of $`x`$, the temperature of the ferromagnetic transition is observed in the range of 30-100 K, the highest values observed in DMS. The ferromagnetic order in the metallic phase is understood, at present, as resulting from the indirect exchange between the Mn<sup>2+</sup> ions due to the spin polarization of the hole gas. The aim of this work is to study the magnetic order resulting of the indirect exchange between magnetic moments in a GaAlAs$`/`$Ga<sub>1-x</sub> Mn<sub>x</sub>As quantum well. A confinement-adapted Ruderman-Kittel-Kasuya-Yosida (RKKY) mechanism is believed to be the most important interaction in such systems, if sufficiently strong doping is provided, as it is the case in metallic samples. It leads to an indirect exchange coupling between Mn<sup>2+</sup> ions, mediated by carriers (holes), which come from the same Mn<sup>2+</sup> ions. This article is organized as follows. In Sec. II we present the calculation of the RKKY exchange for a confined Fermi gas in a semiconductor heterostructure. For the sake of relating our results with other previous ones, we explicitly separate our calculations as intra-subband and inter-subband contributions. We emphasize that, in the quantum limit, i.e., when only the first subband is occupied, the intra-subband exchange is factorized into a purely 2-D RKKY exchange times a form factor determined by the architecture of the confining structure. In Sec. III a Monte Carlo simulation is performed to determine the resulting magnetic phases and the relevant properties. Our calculations reveal that a ferromagnetic order may occur in a single DMS quantum well only beyond a minimum width of the magnetic layer, otherwise the sample is paramagnetic. This is in keeping with recent experiments, and is a consequence of the need of a certain number of magnetic neighbors before a ferromagnetic phase settles in. Depending on the well width and on the effective two-dimensional carrier concentration, a canted phase can occur, with a sizeable net low-temperature magnetization, $`<S>/S_{max}`$, and a well behaved Edwards-Anderson order parameter $`q`$. The origin of the canted spin phase is investigated by analyzing the parallel and the perpendicular magnetizations, and the spin-spin correlation function components. The magnetic susceptibilities are calculated in the existing phases. Finally, in Sec. IV we summarize the results obtained, and comment about the expected magnetic order in the structures analyzed. ## II RKKY interaction in a quantum well The RKKY interaction between localized magnetic moments imbedded in a Fermi gas is a well understood problem, since its early developments almost fifty years ago. However, the new areas for experimental research brought into evidence some theoretical problems concerning the RKKY interaction in low dimensional systems, so far unexplored, such as purely 2-D and 1-D arrangements of magnetic moments, interaction between magnetic layers, magnetic moments in inhomogeneous electron gas, etc. The indirect exchange between localized magnetic moments in a quantum well mediated by a Fermi gas has been addressed several times. Basically, it deals with a confined electron (or hole) gas, therefore a quasi-two-dimensional system, being locally polarized by magnetic moments distributed in a layer. To our knowledge, Korenblit and Shender, were the first to obtain a closed expression to the equivalent of the RKKY interaction in the limiting situation of a purely 2-D electron gas, although Kittel obtained a numerical solution to this question earlier. Larsen derived an expression for a general dimensionality, reproducing Korenblit and Shender results for d=2. A detailed calculation to obtain a closed expression in 2-D was shown by Béal-Monod. Gummich and da Cunha Lima studied the indirect exchange between magnetic impurities in a doped GaAs/AlAs quantum well in the diluted regime, obtaining a ferromagnetic interaction. Finally, another expression for a generic dimensionality has been derived by Aristov. The approximation of the Fermi gas in a quantum well by a purely 2-D system is seldom a reasonable choice. Helman and Baltensperger treated the question of the polarization of an inhomogeneous electron gas in several circumstances, emphasizing the roles of the confined and extended states. The specific case of a DMS quantum well was addressed recently by Dietl et al, but they assumed that the magnetic moments are spread all over the region allowed to the carriers, and in that case the inter-subband contributions to the Curie-Weiss temperature cancel out in a mean field approximation. The interaction potential between a Fermi gas and a set of localized magnetic moments at positions $`\stackrel{}{R}_i`$ is well described by the Hund-type exchange potential: $$H_{\text{ex}}=I\underset{i}{}\stackrel{}{S}_i\stackrel{}{s}(\stackrel{}{r})\delta (\stackrel{}{r}\stackrel{}{R}_i),$$ (1) where $`\stackrel{}{S}_i`$ is the spin of the magnetic moment at position $`\stackrel{}{R}_i`$, which will be treated as a classical variable, and $`\stackrel{}{s}(\stackrel{}{r})`$ is the spin operator of the fermion at $`\stackrel{}{r}`$. $`I`$ is the $`spd`$ interaction. If $`\widehat{\psi }_\sigma (\stackrel{}{r})`$ and $`\widehat{\psi }_\sigma ^{}(\stackrel{}{r})`$ describe the fermion field operator for spin $`\sigma `$, then $$s^z(\stackrel{}{r})=\frac{1}{2}(\widehat{\psi }_{}^{}(\stackrel{}{r})\widehat{\psi }_{}(\stackrel{}{r})\widehat{\psi }_{}^{}(\stackrel{}{r})\widehat{\psi }_{}(\stackrel{}{r})),$$ (2) $$s^+(\stackrel{}{r})=\widehat{\psi }_{}^{}(\stackrel{}{r})\widehat{\psi }_{}(\stackrel{}{r}),$$ (3) $$s^{}(\stackrel{}{r})=\widehat{\psi }_{}^{}(\stackrel{}{r})\widehat{\psi }_{}(\stackrel{}{r}),$$ (4) with the usual definitions of $`s^+=s_x+is_y`$, and $`s^{}=s_xis_y`$. Instead of free fermions in a 3-D space, the electrons and holes in a ssemiconductor heterostructure are confined in the growth direction, assumed to be the z-axis, due to the mismatch of the conduction and valence band edges. Since they are free particles in the plane perpendicular to that growth direction, i.e., in the plane parallel to the semiconductor interfaces, their field operator are given by: $$\widehat{\psi }_\sigma (\stackrel{}{r})=\frac{1}{\sqrt{A}}\underset{n,\stackrel{}{k}}{}e^{i\stackrel{}{k}.\stackrel{}{R}}\varphi _n(z)\eta c_{n,\stackrel{}{k},\sigma },$$ (5) where $`A`$ is the normalization area, $`\stackrel{}{k}`$ is a wavevector in the plane ($`x,y`$), $`\eta `$ is the spin tensor, $`\varphi _n(z)`$ is the envelope function which describes the motion of the fermion in the $`z`$-direction, and $`c_{n,\stackrel{}{k},\sigma }`$ is the fermion annihilation operator for the state ($`n,\stackrel{}{k},\sigma `$). $`\stackrel{}{R}`$ represents a vector in the 2-D coordinates plane ($`x,y`$). The usual RKKY perturbation calculation up to second order leads to the correction on the ground state energy of the system formed by the set of (classical) localized moments and the Fermi gas : $$\delta E^{(2)}=\delta E_a^{\left(2\right)}+\delta E_b^{\left(2\right)},$$ (6) where $$\delta E_a^{\left(2\right)}=\left(\frac{I}{2N}\right)^2\underset{i}{}\underset{n,n^{}}{}\varphi _n(z_i)^2\varphi _n^{}(z_i)^2S_i(S_i+1)\underset{\stackrel{}{q}}{}\chi ^{n,n^{}}(\stackrel{}{q}),$$ (7) and $$\delta E_b^{\left(2\right)}=\left(\frac{I}{2N}\right)^2\underset{j}{}\underset{ij}{}\underset{n,n^{}}{}\underset{\stackrel{}{q}}{}2\text{Re}\left[\varphi _n^{}(z_i)\varphi _n^{}(z_i)\varphi _n^{}^{}(z_j)\varphi _n(z_j)e^{i\stackrel{}{q}.(\stackrel{}{R}_i\stackrel{}{R}_j)}\right]\chi ^{n,n^{}}(\stackrel{}{q})\stackrel{}{S}_i\stackrel{}{S}_j.$$ (8) Eqs. (7) and (8) are, respectively, the self-energy term and the RKKY exchange in the form they assume for confined fermions. The coordinates ($`\stackrel{}{R}_i,z_i`$) describe the position of the impurity $`i`$ in the plane ($`x,y`$), and in the growth direction; $`\stackrel{}{q}`$ is a two-dimensional wavevector. $`\chi ^{n,n^{}}(\stackrel{}{q})`$ is the equivalent to the Lindhard function: $$\chi ^{n,n^{}}(\stackrel{}{q})=\underset{\stackrel{}{k}}{}\frac{\theta (E_Fϵ_{n,\stackrel{}{k}})\theta (E_Fϵ_{n^{},\stackrel{}{k}+\stackrel{}{q}})}{ϵ_{n^{},\stackrel{}{k}+\stackrel{}{q}}ϵ_{n,\stackrel{}{k}}}.$$ (9) Eqs. (7) and (8) are used to define the exchange Hamiltonian: $$H_{ex}=\underset{i,j}{}J_{ij}\stackrel{}{S}_i\stackrel{}{S}_j.$$ (10) For $`ij`$, $$J_{ij}=\left(\frac{I}{2A}\right)^2\underset{n,n^{}}{}\underset{\stackrel{}{q}}{}2\text{ Re}\left[\varphi _n^{}(z_i)\varphi _n^{}(z_i)\varphi _n^{}^{}(z_j)\varphi _n(z_j)e^{i\stackrel{}{q}.(\stackrel{}{R}_i\stackrel{}{R}_j)}\right]\chi ^{n,n^{}}(\stackrel{}{q}).$$ (11) ### A Intra-subband terms To our knowledge, complete calculations of Eq. (8) have only been performed for intra-subband transitions, using different approaches. For the sake of completeness, we will show how this is achieved in our treatment. The contribution of a subband $`n`$ to the exchange reads: $$J_{ij}^{(n)}=\left(\frac{I}{2A}\right)^2\varphi _n(z_i)^2\varphi _n(z_j)^2\underset{\stackrel{}{q}}{}2\mathrm{cos}(\stackrel{}{q}\stackrel{}{R}_{ij})\chi ^{n,n}(\stackrel{}{q}).$$ (12) We observe that, in the so-called quantum limit, when only the first subband ($`n=0`$) is occupied, the difference between Eq. (12) and the indirect exchange mediated by a 2D electron gas comes from the non-uniform charge density in the confining direction $`z`$. Actually, in that case, Eq. (12) factorizes into a form factor $`^{ij}`$ and a purely 2-D exchange: $$J_{ij}^{(0)}=^{ij}J_{ij}^{(2D)}.$$ (13) where $$^{ij}=\varphi _0(z_i)^2\varphi _0(z_j)^2$$ (14) and $$J_{ij}^{2D}=\left(\frac{I}{2A}\right)^2\underset{\stackrel{}{q}}{}2\mathrm{cos}(\stackrel{}{q}\stackrel{}{R}_{ij})\chi ^{n,n}(\stackrel{}{q}).$$ (15) It is easy to show, by using the dimensionless variables $`x=kR_{ij}`$ , and $`y=qR_{ij}`$, that the Fourier transform of the modified Lindhard function, appearing in the summation in $`\stackrel{}{q}`$ at the rhs of Eq. (12), becomes: $$\chi ^n(R_{ij})=\frac{4m_t^{}A^2}{\pi ^3\mathrm{}^2R_{ij}^2}_0^{\mathrm{}}𝑑yyJ_0(y)_0^{k_F^{(n)}.R_{ij}}𝑑xx_0^{\pi /2}𝑑\varphi \frac{1}{y^24x^2\mathrm{cos}^2\varphi },$$ (16) where $`\chi ^n(R_{ij})=_\stackrel{}{q}2\mathrm{cos}(\stackrel{}{q}\stackrel{}{R}_{ij})\chi ^{n,n}(\stackrel{}{q})`$. The transversal effective mass, $`m_t^{}`$, is assumed as isotropic in the plane parallel to the interfaces. As usual, $`k_F^{(n)}=`$ $`\sqrt{2m_t^{}(E_Fϵ_n)}/\mathrm{}`$. Performing the $`\varphi `$ integral and changing variables again ($`y/2xy`$), $$\chi ^n(R_{ij})=\frac{2m_t^{}A^2}{\pi ^2\mathrm{}^2R_{ij}^2}_0^{k_F^{(n)}.R_{ij}}𝑑xx_1^{\mathrm{}}𝑑yJ_0(2xy)\frac{1}{\sqrt{y^21}}.$$ (17) The integral on $`y`$ is straightforward: $$\chi ^n(R_{ij})=\frac{m_t^{}A^2}{\pi \mathrm{}^2R_{ij}^2}_0^{k_F^{(n)}.R_{ij}}𝑑xxJ_0(x)N_0(x).$$ (18) After performing the integral on $`x`$, Eq. (18) results in: $$\chi ^n(R_{ij})=\frac{m_t^{}A^2}{\pi \mathrm{}^2}k_F^{(n)2}[J_0(k_F^{(n)}R_{ij})N_0(k_F^{(n)}R_{ij})+J_1(k_F^{(n)}R_{ij})N_1(k_F^{(n)}R_{ij})].$$ (19) This expression for the real space Lindhard function has been derived in a different context, by several authors. The final expression for the intra-subband exchange becomes: $`J_{ij}^{(n)}`$ $`=`$ $`\left({\displaystyle \frac{I}{2}}\right)^2{\displaystyle \frac{m_t^{}}{\pi \mathrm{}^2}}k_F^{(n)2}\varphi _n(z_i)^2\varphi _n(z_j)^2\times `$ (21) $`[J_0(k_F^{(n)}R_{ij})N_0(k_F^{(n)}R_{ij})+J_1(k_F^{(n)}R_{ij})N_1(k_F^{(n)}R_{ij})].`$ ### B Inter-subband terms The contribution of the inter-subband terms cannot be expressed easily in a closed form. Starting over from Eq. (11), and using the same approach as in Ref. , we arrive to: $$J_{ij}^{(n,n^{})}=\left(\frac{I}{2}\right)^2\frac{1}{\pi }\text{Re}\left[\varphi _n^{}^{}(z_i)\varphi _n(z_i)\varphi _n^{}(z_j)\varphi _n^{}(z_j)\right]_0^{\mathrm{}}𝑑qqF_{n,n^{}}(q)J_0(qR_{ij}),$$ (22) where we used $$F_{n,n^{}}(q)=\frac{m_t^{}}{\mathrm{}^2}\frac{1}{\pi ^2}d^2k\frac{q^2+\mathrm{\Delta }_{n^{},n}}{(q^2+\mathrm{\Delta }_{n^{},n})^2(2\stackrel{}{k}\stackrel{}{q})^2}\theta (ϵ_n^{}E_F),$$ (23) and $`\mathrm{\Delta }_{n^{},n}=2m_t^{}(E_n^{}E_n)/\mathrm{}^2`$. The integral in Eq. (23) is, then, straightforward: $$F_{n,n^{}}(q)=\frac{m_t^{}}{2\pi \mathrm{}^2}(1\frac{\mathrm{\Delta }_{n^{},n}}{q^2})[1\sqrt{1(\frac{2k_F^{(n)}q}{q^2+\mathrm{\Delta }_{n^{},n}})^2}\theta (q^2+\mathrm{\Delta }_{n^{},n}2qk_F^{(n)})]\theta (ϵ_n^{}E_F).$$ (24) ## III Monte Carlo simulation: Magnetic ordering In order to determine the possible magnetic order in GaAs:Mn quantum wells, we have performed extensive Monte Carlo simulations. Classical spins $`𝐒_i`$, randomly distributed on the cation sites with concentration $`x`$, are assumed to interact through the RKKY exchange Hamiltonian defined by Eq. (10). In the present work we have focused our attention on metallic single quantum wells with the Mn concentration $`x=5\%`$, and we have neglected possible (anti-)ferromagnetic interaction between the nearest neighbors and the next nearest neighbors pairs. The RKKY exchange interaction derived in Sec. II is assumed to be effective within a cutoff radius which we have taken as $`R_c=4a`$, $`R_c=2a`$, and $`R_c=a`$, where $`a`$ is the fcc lattice parameter of GaAs. This makes the smallest value assumed for $`R_c`$ nearly equal to the hole mean free path estimated from bulk transport measurements. The highest value, $`R_c=2.2`$nm, amounts to 3-4 values of that mean free path. The consequences of the cutoff radius on the results will be discussed below. The calculation is performed in a finite box, whose axes are parallel to the directions. Its dimensions are $`L_x=L_y`$, and $`L_z=Na/2`$, and $`N`$ is the number of DMS monolayers (ML) in the barrier. Periodic boundary conditions are imposed in the $`(x,y)`$ plane. The lateral dimensions are adjusted in such a way that the total number $`N_s`$ of spins is about 4400, for all samples with different $`L_z`$. The initial spin orientations are randomly assigned. At a given temperature, the energy of the system due to RKKY interaction is calculated, and the equilibrium state for a given temperature is sought by changing the individual spin orientation according to the Metropolis algorithm. A slow cooling stepwise process is accomplished making sure that the thermal equilibrium is reached at every temperature. The resulting spin configuration is taken as the starting configuration for the next step at a lower temperature. For every temperature, the average magnetization $`<M>`$ and the Edwards-Anderson (EA) order parameter $`q`$ are calculated. The latter is defined as $$q=\frac{1}{N}\underset{i=1}{\overset{N}{}}\left(\underset{\alpha }{}\left\|\frac{1}{t}\underset{t^{}=t_0}{\overset{t_0+t}{}}S_{i\alpha }(t^{})\right\|^2\right)^{1/2},$$ (25) where $`\alpha =x`$, $`y`$ and $`z`$. In order to avoid spurious results in obtaining the average over a large time interval $`t`$, a summation on $`t^{}`$ is performed starting from a time $`t_0`$, when the system already reached the thermal equilibrium. In our calculations we used the value $`N_0\beta =1.2`$ eV, taken from Ref. , for which we obtain transition temperatures in good agreement with the experimental data. Recently, the value $`N_0\beta =0.9`$ eV has been obtained theoretically, confirming the result of Ref.. The earlier estimate $`\|N_0\beta \|=3.3`$ eV is probably too high. Monte Carlo calculations have been performed in sixteen samples, as shown in Table I, the well widths varying from 25 Å to 100 Å (9 ML to 35 ML). In sample #01, for a well width of 50 Å and assuming $`R_c=8`$ ML, we tested the effect of a small hole concentration, making $`p`$ just 1% of $`x`$. This amounts to $`p1.1\times 10^{19}cm^3`$. We found that the spins can be arranged in a ferromagnetic phase even at this carrier concentration. For that sample the calculation gives a transition temperature near 27 K. In Fig. 1 the normalized average magnetization is shown for samples #02 to #07 with different well widths, but with a fixed carrier concentration of 10% of $`x`$, i.e., $`p1.1\times 10^{20}cm^3`$, and the same cutoff $`R_c=8`$ ML. For a very thin well (25 Å), we found that the sample is paramagnetic. The EA order parameter for these samples is shown in Fig. 2. It can be observed that, for sample #02, there is no phase transition. All the other curves in Fig. 2 (samples #03 to #07) are characteristic of an ordered phase. Raising the well width (starting from 35 Å), the number of interacting neighbors increases, and the sample shows successively a ferromagnetic phase (samples #03, #04, and #06) and a canted spin arrangement (samples #5 and #7). The value chosen for the cutoff radius is larger than the first zero of the $`J_{ij}^{\text{RKKY}}`$, so antiferromagnetic interactions are turned on. With these choices of $`p`$ and $`R_c`$, depending on the well width, the antiferromagnetic interactions can settle a fraction of the Mn magnetic moments antiparallel. This is the origin of the canted spins phase. The final average magnetizations in Fig. 1 at $`T=0`$ K are only a fraction of the maximum magnetization, around 60% for sample #5, and 70% for sample #7. The transition temperatures were found in the range of 40 to 50 K. No spin glass phase with vanishing magnetization was found. The existence of the canted phase requires a careful analysis of a possible dependence on the cooling process initial conditions. In order to clarify that point, we performed two additional simulations on sample #5, starting at $`T=25`$ K, where the sample shows already a significant partial alignment of spins, and we proceeded with the slow cooling down process. In the first simulation, we assumed a starting configuration in which all spins are aligned perpendicularly to the interface. In the second, the spins alignment is made parallel to the interfaces. This choice of a rather low starting temperature for cooling is necessary, otherwise the thermal excitation would immediately randomize the initial configuration. The results are shown in the inset of Fig. 1. We observe that the appearance of the canted phase does not depend on the choice of the starting configuration, and the three simulations converge, within statistical fluctuations, to the same value of the magnetization at every temperature step. In Fig. 3 the magnetization as a function of temperature is shown for samples #08 to #12, with a higher carrier concentration, $`p=0.25x`$, but keeping the same cutoff radius of 8 ML. The EA order parameter for these samples shown in Fig. 4 as a function of temperature, gives evidence for the existence of ordered phases. The Fermi wave number increases with carrier concentration, what decreases the in-plane distance corresponding to the first zero of the RKKY interaction. In consequence, in some samples, the magnetic moments order in the canted spin phase (samples #10 to #12). In other samples, however, the ferromagnetic interaction prevails and the total magnetization is reached (samples #08 and #09). Notice that the canted phase appears here already for $`L=50\AA `$, while the phase is still ferromagnetic for that width when $`p=0.1x`$. The transition temperatures were estimated to lie between 30 K and 50 K. In Fig. 5 we explored the effect of the cutoff radius on the spin ordering with $`p=0.25x`$. A cutoff radius $`R_c=2`$ ML was used in samples #13 ($`L=60\AA `$) and #14 ($`L=80\AA `$), while $`R_c=4`$ ML was used in samples #15 ($`L=60\AA `$) and #16 ($`L=80\AA `$). The former is too small, resulting in the fact that no net magnetization is allowed. The respective EA order parameters (Fig. 6) are typical of a paramagnetic phases for samples #13 and #14. The choice of a larger $`R_c`$ ordered spins in a ferromagnetic phase, with transition temperatures calculated to be $`35`$ K for the sample #15, and $`80`$ K for sample #16. The cutoff radius of 4 ML is smaller than the first zero of $`J_{ij}^{\text{RKKY}}`$. In this situation the canted spin arrangement is not allowed, and the sample can be either paramagnetic or ferromagnetic. Finally in Fig. 7 the magnetic susceptibility for sample #07 is presented, calculated from the equilibrium magnetization fluctuations. Notice that the peak in the in the susceptibility indicates the same $`T_c`$ as estimated from the magnetizations and EA order parameter curves (Figs. 1 and 2 respectively). ## IV Discussions and Final Comments The results of the Monte Carlo simulations indicate that, besides the choice of the range of the interaction (the cutoff radius $`R_c`$), two parameters are determining in the magnetic ordering in these heterostructures: the magnetic layer width $`L`$, and the carrier concentration $`p`$. It is observed from resistivity measurements, that the holes have a small mean free path in these materials. The criteria for choosing the cutoff in a Monte Carlo simulation must take into account natural scales of the interaction. These scales are the transport mean free path (since the RKKY interaction is based on the very existence of free carriers), and also the spin coherence length. We tested different $`R_c`$’s, in the range of the mean free path estimated from resistivity data. The influence of this parameter on the magnetic order is simple. If $`R_c`$ is smaller than the first zero of $`J_{ij}^{\text{RKKY}}`$ (a proper choice for the case in which the transport mean free path or the spin coherence length are small), there are two possibilities of magnetic order: ferromagnetic or paramagnetic. For larger $`R_c`$, on the other hand, corresponding to the cases where both the transport mean free path and the coherence length are large, a canted magnetization may be observed. In all the explored samples, no spin glass phase was found. This is presumably due to the fact that while the spin frustration exists, as witnessed by the occurrence of canted spin phases, it is not strong enough to produce a spin glass phase, as in canonical metallic spin glasses. A spin glass phase in these DMS structures would probably require a much higher carrier concentration. In what concerns the influence of the width of the quantum well, we conclude that, for thin layers, the number of interacting ions is small within the cutoff radius, and the sample is paramagnetic. When $`L`$ becomes larger, the number of interacting ions increases and a collective magnetic ordering may be observed. The fact that the appearance of a magnetic order occurs only above a minimum thickness of the magnetic layer has already been observed experimentally. Since the RKKY interaction oscillates with the argument ($`k_\text{F}R`$), which depends on the carrier concentration, raising $`p`$ produces a change on $`k_\text{F}`$, increasing the number of oscillations of $`J_{ij}^{\text{RKKY}}`$. Therefore, antiferromagnetic interactions can be turned on, resulting in all kind of couplings. In this situation, ferromagnetic and antiferromagnetic interactions compete in establishing the magnetic order, which depends on the other sample characteristics, resulting in a total or in a partial alignment of the Mn magnetic moments. The occurrence of partial magetization (about 40%) has also been observed in samples of (In,Mn)As/(Ga,Al)Sb. To conclude, we believe that the RKKY mechanism explains the high transition temperatures experimentally observed in Ga<sub>1-x</sub>Mn<sub>x</sub>As heterostructures, at least in the metallic phase. Additionally, it explains, as it becomes clear after these Monte Carlo simulations, the occurrence of samples showing a partial magnetization at low temperatures. The possibility of having a ferromagnetic phase in samples with a low Mn concentration, i.e., in the ferromagnetic insulator Ga<sub>1-x</sub>Mn<sub>x</sub>As, remains to be explained. ###### Acknowledgements. This work was supported by CENAPAD-SP (Centro Nacional de Processamento de Alto Desempenho em São Paulo) UNICAMP/FINEP-MCT, CAPES, CNPq and FAPERJ in Brazil, and by the PAST grant from Ministère de l’Éducation Nationale, de l’Enseignement Supérieure et de la Recherche (France).
warning/0002/gr-qc0002076.html	ar5iv	text	# Locating Boosted Kerr and Schwarzschild Apparent Horizons. ## I Introduction Apparent horizon locators play an integral role in the application of black hole excision techniques in the computational evolution of black hole spacetimes. Excision techniques delete the regions of spacetime that contain the curvature singularity from the computational domain. Assuming cosmic censorship, these curvature singularities are expected to be contained within an event horizon. The event horizon is a causal boundary whose interior does not causally affect the exterior spacetime; as a result it is possible to excise a region within the event horizon, thereby excising the black hole’s curvature singularity. In our approach to computationally solving the Einstein field equations we focus on the use of Cauchy techniques, in which a 3+1 splitting of spacetime into a foliation of spacelike hypersurfaces, $`\mathrm{\Sigma }`$, parametrized in time, is the basis for an evolution in time. The result of this splitting is a system of elliptic and hyperbolic partial differential equations in the 3-metric, $`\gamma _{ij}`$, and extrinsic curvature, $`K_{ij}`$. These are the four constraint equations and 12 first-order-in-time evolution equations. The Cauchy approach starts with an initial spacelike slice with $`\gamma _{ij}`$ and $`K_{ij}`$ set by solving an initial value problem (the elliptic constraint equations). One then uses the evolution equations to evolve to the next spacelike slice obtaining $`\gamma _{ij}`$ and $`K_{ij}`$ at the next time (See York for a detailed discussion). In the evolution of black hole spacetimes in this manner we do not have a complete history of the entire spacetime and hence do not have a knowledge of the location of the event horizon. Since the event horizon is a global object that depends on geometric information for all time (or at least until the black hole becomes quiescent) we cannot use it to determine an inner excision boundary in our Cauchy evolution. There is, however, an alternative, and that is to use the apparent horizon surface which is a local object, locatable (if it exists) with $`\gamma _{ij}`$ and $`K_{ij}`$ at one time. The apparent horizon is the outermost marginally trapped surface. It is a closed spacelike 2-surface whose future-directed outgoing null normals have zero divergence. The apparent horizon is slicing dependent and may not necessarily exist even though an event horizon does. An example of this is given by Wald and Iyer through nonspherically-symmetric slicings for the Schwarzschild spacetime. Provided a non-pathological slicing is chosen the apparent horizon or any trapped surface within it may be used for excising the black hole singularity. These surfaces define a local causal structure that distinguishes instantaneously escaping null rays from those that are certain to collapse. This distinction makes their treatment very amenable to computational black hole excision techniques. Since these surfaces can be determined with geometric information at one instant of time, they are used in practice as an inner boundary in Cauchy evolutions. With this purpose in mind, we developed a 3D apparent horizon locator that utilizes $`\gamma _{ij}`$ and $`K_{ij}`$ on a given spacelike slice of spacetime and locates an apparent horizon. Once the apparent horizons are located, a region contained within the apparent horizon is excised. Thus the method is really a trapped-surface excision. There has been a variety of work done on apparent horizon location in spherical symmetry, axisymmetry and 3D. We focus solely on the 3d locators. These can be classified into those that use finite difference methods, and those that utilize pseudo-spectral schemes. Further, one can classify each of these finders in terms of those that use flow methods versus those that directly solve the apparent horizon equation either via a minimization scheme or Newton’s method for root finding. One of the first published 3d apparent horizon locators was developed by Nakamura, Kojima and Oohara. Their method expands the apparent horizon shape function, $`r=\rho (\theta ,\varphi )`$ in spherical harmonics to some maximum $`l=l_{max}`$: $`\rho (\theta ,\varphi )`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{l_{max}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}a_{lm}Y_{lm}(\theta ,\varphi ).`$ (1) With this expansion Nakamura et al. evaluate the apparent horizon equation and solve for the coefficients $`a_{lm}`$ via a “direct” functional iteration scheme. Kemball and Bishop reimplemented this approach and made modifications that led to improved convergence and stability behaviour. Anninos et al. and Baumgarte et al. implement similar methods that involve an expansion of $`\rho (\theta ,\varphi )`$; the primary differences being that they expand in terms of symmetric trace-free tensors and use Powell’s method for minimization of the square of the apparent horizon equation, Eq.(2) below (which is related to the expansion of the outgoing null normals). Thornburg gives a very good treatise on the use of finite differencing to solve the apparent horizon equation using spherical coordinates $`(r,\theta ,\varphi )`$ via Newton’s method. He discusses in general how algebraic Jacobians may be applied in a full 3D context. His implementation for horizon finding however is axisymmetric; his full 3-d finder suffers from problems with the z-axis ($`\theta =0,\pi `$). Our method for finding horizons uses closely related concepts except that we finite difference in cartesian coordinates, eliminating any potential z-axis problems. Another class of apparent horizon locators casts the elliptic apparent horizon equation into a parabolic one as suggested by Tod, who suggested the use of flow methods in locating apparent horizons. Bernstein implemented Tod’s algorithm in axisymmetry using finite differences, but encountered problems with differencing on a sphere in spherical coordinates in the general case. The advantage to flow-methods is that one can start with an arbitrary initial guess and flow towards the apparent horizon(s). In some implementations it is possible to find multiple apparent horizon surfaces starting from a single initial guess surface (i.e., there is a topology change in the course of location of the apparent horizon). Pasch uses a level-set method to locate multiple apparent horizons in 3d. He demonstrates his method utilizing time-symmetric conformally flat initial data for multiple black holes. An hybrid flow/level-set-like method utilizing our approach to evaluating the outgoing null expansions via Cartesian finite differences has been implemented by Shoemaker et al.. That method flows towards the apparent horizon(s) from an arbitrary initial guess allowing for topology changes. Gundlach has implemented a “fast flow” method for finding apparent horizons. In the sections that follow we give a brief discussion of the algorithm used and relegate the details to the appendices. The model spacetime in which all of the results are presented is discussed in section III. In section IV we discuss tests of the algorithm and demonstrate that for distorted apparent horizons for boosted Kerr black holes the algorithm fairs well. ## II Boundary Value Problem Approach On a particular 3D spacelike hypersurface, $`\mathrm{\Sigma }`$, from our foliation of spacetime we are given the 3-metric, $`\gamma _{ij}`$ and the extrinsic curvature, $`K_{ij}`$. Let $`𝒮`$ be a closed 2-surface in $`\mathrm{\Sigma }`$. At any point $`p`$ on $`𝒮`$ we can define a spacelike normal, $`s^a`$, to $`𝒮`$, and a time-like normal, $`n^a`$, to $`\mathrm{\Sigma }`$. From these we can construct the outgoing null normal, $`k^a`$, at $`p`$. If the divergence $`_ak^a`$($`_a`$ is the covariant derivative compatible with the spacetime metric, $`g_{ab}`$) is zero everywhere on $`𝒮`$, then $`𝒮`$ is a marginally trapped surface (MTS). The apparent horizon is the outermost such MTS. The expansion of the outgoing null normals, $`_ak^a=0`$, can be rewritten as an equation entirely in terms of quantities in $`\mathrm{\Sigma }`$: $`D_is^i+K_{ij}s^is^jK`$ $`=`$ $`0.`$ (2) Here $`D_i`$ is the covariant derivative compatible with the 3-metric $`\gamma _{ij}`$ and $`K`$ is the trace of the extrinsic curvature. The apparent horizon equation is an elliptic partial differential equation on $`𝒮`$ (for a function of coordinates on $`𝒮`$). This can be made apparent by noting that an MTS is a closed 2-surface; spherical coordinates are a natural set of coordinates for $`𝒮`$. The location of $`𝒮`$ can then be written as the radial distance from the origin of the coordinate system, $`r=\rho (\theta ,\varphi )`$. In general one can generate a foliation of such closed spacelike 2-surfaces based on the distance from the MTS. This is given by $`\phi `$ $`=`$ $`r\rho (\theta ,\varphi ),`$ (3) where the $`\phi =0`$ level surface is the MTS. From $`\phi `$ we define a spacelike vector field, the normal $`s^i`$ $`=`$ $`\gamma ^{ij}_j\phi /\sqrt{\gamma ^{kl}_k\phi _l\phi },`$ (4) at every point on these level surfaces. Substituting this into Eq.(2) results in a second order elliptic partial differential equation on $`𝒮`$, $`F[\rho ]`$ $`=`$ $`\gamma ^{ab}_a_b\phi +\gamma _{,a}^{ab}_b\phi {\displaystyle \frac{1}{2}}\omega ^1\gamma ^{ab}\gamma _a^{cd}_b\phi _c\phi _d\phi `$ (7) $`\omega ^1\gamma ^{ab}\gamma ^{cd}_b\phi _a_c\phi _d\phi +\mathrm{\Gamma }_{ab}^a\gamma ^{bc}_c\phi `$ $`+\omega ^{\frac{1}{2}}K_{ab}\gamma ^{ac}\gamma ^{bd}_c\phi _d\phi \omega ^{\frac{1}{2}}K=0,`$ where $`\omega =\gamma ^{cd}_c\phi _d\phi `$ and $`\mathrm{\Gamma }_{bc}^a`$ is the connection coefficient associated with the 3-metric $`\gamma _{ab}`$. Our approach involves casting Eq.(7) as a boundary value problem on $`𝒮`$. As stated, points on $`𝒮`$ are parametrized in spherical coordinates $`(\theta [0,\pi ],\varphi [0,2\pi ))`$. $`𝒮`$ is discretized into an uniform mesh, $`\widehat{𝒮}`$, of $`N_\theta \times N_\varphi `$ points where $`N_\theta =N_\varphi =N_s`$. The domain on $`\widehat{𝒮}`$ is $`(0\theta \pi ;0\varphi <2\pi )`$ where at the poles, $`\theta =0,\pi `$ all $`N_\varphi `$ points are identified as one. The $`\varphi =2\pi `$ branch cut is identified with the $`\varphi =0`$ line. The boundary conditions simply are periodicity at $`\varphi =2\pi `$ and $`\varphi `$ identification at $`\theta =0,\pi `$. These boundary conditions are key to avoiding the coordinate singularities at the poles in combination with using Cartesian coordinates to discretize partial derivatives on $`\widehat{𝒮}`$. We treat $`\phi `$ as a function of Cartesian coordinates $`x,y,z`$ and center on each mesh point of $`\widehat{𝒮}`$ a 3-d Cartesian difference stencil of 27 points. Using the form Eq.(3) we interpolate values of $`\phi (x,y,z)`$ onto each of the 26 stencil points surrounding each $`\widehat{𝒮}`$ stencil point. (See the appendix for more details.) Using this difference stencil we can evaluate first, second and mixed derivatives of $`\phi (x,y,z)`$ as required by the discretized version of Eq.(7). At every point on $`\widehat{𝒮}`$ we then construct the residual $`\widehat{F}[\widehat{\rho }]`$ on $`\widehat{𝒮}`$ (Note that the discrete version of a continuum quantity, $`T`$, is denoted by, $`\widehat{T}`$). The problem at hand is to solve for a $`\rho `$ that yields $`F[\rho ]=0`$. Since $`F[\rho ]`$ is a nonlinear operator (as shown in Eq.(7)), we use Newton’s method to solve for $`\rho `$. Given an initial guess surface, $`\rho =\rho _0`$, we wish to find a $`\delta \rho `$ (the change in the surface) that leads to $`F[\rho _0+\delta \rho ]=0`$ or, to lowest order, $$F[\rho _0+\delta \rho ]=F[\rho _0]+\frac{F[\rho ]}{\rho }\|_{\rho =\rho _0}\delta \rho +O(\delta \rho ^2)=0.$$ (8) The Jacobian of $`F[\rho ]`$ is defined to be $$J\frac{F}{\rho }.$$ (9) In the discretized case, $`\widehat{J}`$ is an $`N\times N`$ matrix, where $`N`$ is the total number of points used in the discretization. To obtain a $`\delta \rho `$ that leads to $`F[\rho +\delta \rho ]=0`$ we must solve, $$J\delta \rho =F[\rho ]$$ (10) for $`\delta \rho `$. Computationally our tasks are to first evaluate the discrete form of the Jacobian matrix, $`\widehat{J}`$ and second to solve the discrete form of Eq.(10). We numerically compute the Jacobian matrix by perturbing the surface pointwise and examining the effect of of the perturbation on the residual, $`\widehat{F}`$. Let $`\overline{\mu }`$ denote “independent” points in the computational mesh, $`\widehat{𝒮}`$. By independent we mean the unique points (points modulo boundary identification) on $`\widehat{𝒮}`$. In particular from the identifications made earlier, there are $`N_s^22N_s+2`$ independent points in $`\widehat{𝒮}`$. $`N_s=N_\theta =N_\varphi `$ points at each of the poles are treated as one point. $`\overline{\mu }=1`$ represents the $`\theta =0`$ point for all the $`N_\varphi `$ points($`0\varphi <2\pi `$) and $`\overline{\mu }=N_s^22N_s+2`$ is the $`\theta =\pi `$ point. Eq.(10) then becomes a linear system of equations where $`\widehat{J}`$ is a $`(N_s^22N_s+2)\times (N_s^22N_s+2)`$ matrix and $`\widehat{F}`$ and $`\delta \widehat{\rho }`$ are vectors of length $`N_s^22N_s+2`$. The $`\overline{\mu }\overline{\nu }`$ component of the Jacobian is then computed by perturbing $`\rho `$ at the $`\overline{\nu }`$-th point and computing the change in the residual, $`\widehat{F}`$ at the $`\overline{\mu }`$-th point. Using a first order forward difference approximation we have, $$\widehat{J}_{\overline{\mu }\overline{\nu }}=\frac{1}{ϵ}\left\{\widehat{F}_{\overline{\mu }}[\widehat{\rho }_{\overline{\mu }}+ϵ]\widehat{F}_{\overline{\mu }}[\widehat{\rho }_{\overline{\nu }}]\right\},$$ (11) where $`ϵ`$ is the amount by which we perturb the surface. We define $`ϵ`$ to be the perturbation parameter. The process for generating the components then involves numerically evaluating $`\widehat{F}`$ in only a small neighborhood of the $`\nu `$-th point since $`\widehat{F}`$ has a domain of dependence dependent on the finite difference operators used. In this case the operators are (finite difference) derivative operators convolved with interpolation operators. The solution of Eq.(10) is achieved via Newton’s method. We solve for the change $`\delta \widehat{\rho }`$ that leads to a new surface $`r=\widehat{\rho }_{}(\theta ,\varphi )`$ that yields $`\widehat{F}[\widehat{\rho }_{}]0`$ up to $`O(h^2)`$, where $`h`$ is the Cartesian stencil spacing, proportional to $`\delta \theta `$. Newton’s method then involves updating $`\widehat{\rho }(\theta ,\varphi )`$ by $`\delta \widehat{\rho }`$. If $`F`$ were a linear operator then one iteration would result in a $`\delta \widehat{\rho }`$ that leads to a solution. Since $`F`$ is nonlinear we have to iterate until the $`L_2`$-norm of the residual, $`\widehat{F}_2`$, is driven down to a chosen stopping criterion. We discuss the implementation details in the appendices and discuss further the properties of the Jacobian and Newton’s method in the results section. We now turn our attention to the model spacetime in which we shall conduct our numerical experiments. ## III 3+1 Splitting of the Kerr-Schild metric In the rest of the paper we focus on tests of the algorithm based on boosted Schwarzschild and Kerr black holes. The particular form that we use is given by the Kerr-Schild metric: $$g_{\mu \nu }=\eta _{\mu \nu }+2Hl_\mu l_\nu ,$$ (12) where $`l_\mu `$ is an ingoing null vector (i.e: $`g^{\mu \nu }l_\mu l_\nu =\eta ^{\mu \nu }l_\mu l_\nu =0`$), $`H`$ is a scalar function of the spacetime coordinates and $`\eta _{\mu \nu }`$ is the Minkowski spacetime metric. This metric describes the Kerr and Schwarzschild spacetimes. We note that under a Lorentz transformation the spacetime metric is form invariant. By definition such a transformation takes $`\eta _{\mu \nu }\eta _{\mu \nu }`$ and $`l^\mu `$ and $`H`$ are transformed to a new null vector and left unchanged (though evaluated at the new coordinate labels for the same event) respectively. This property makes our analysis easier since a 3+1 decomposition of Eq.(12) has the same form as a 3+1 decomposition of the boosted metric. As we shall see, we only need to specify $`H`$, $`l_\mu `$ and their spacetime derivatives in order to obtain the 3-metric and extrinsic curvature on $`\mathrm{\Sigma }`$. For a vacuum spacetime, $`l^\mu `$ is geodesic and in the Kerr and Schwarzschild spacetimes is the tangent to geodesics of ingoing photons. The null nature of Eq.(12) leads to a slicing of these spacetimes that is well behaved at the horizon. That is, spacelike slices penetrate the horizon and hit the black hole singularity. This is a desirable property for black hole excision in computational applications and this metric has shown itself to be a good choice for the study of single and multiple black hole evolutions with exicision. For the Kerr spacetime, $`H`$ and $`l_\mu `$ are given by $$H=\frac{Mr^3}{r^4+a^2z^2}$$ (13) and $`l_\mu =(1,{\displaystyle \frac{rx+ay}{r^2+a^2}},{\displaystyle \frac{ryax}{r^2+a^2}},{\displaystyle \frac{z}{r}}),`$ (14) where $`r`$ is given by $$\frac{x^2+y^2}{r^2+a^2}+\frac{z^2}{r^2}=1,$$ (15) or $$r^2=\frac{1}{2}\left(\rho ^2a^2\right)+\sqrt{\frac{1}{4}\left(\rho ^2a^2\right)^2+a^2z^2}.$$ (16) $`M`$ is the mass of the Kerr black hole and $`a=J/M`$ is the angular momentum of the black hole and $`\rho =\sqrt{x^2+y^2+z^2}`$. In the $`a0`$ limit we get the Schwarzschild metric in ingoing Eddington-Finkelstein coordinates where $$H=\frac{M}{r},$$ (17) $$l_\mu =(1,\frac{x}{r},\frac{y}{r},\frac{z}{r})$$ (18) and $`r=\sqrt{x^2+y^2+z^2}`$. In a spacelike slice of either Kerr or Schwarzschild spacetimes the apparent horizon is known to coincide with the intersection of the event horizon with that slice. In the Kerr spacetime then the apparent horizon is a surface of radius $`r=r_+`$: $`r_+`$ $`=`$ $`M+\sqrt{M^2a^2}`$ (19) and area $`A`$ $`=`$ $`4\pi (r_+^2+a^2).`$ (20) In the more general nonstationary case the apparent horizon and event horizon will not coincide. We use the properties of the Kerr and Schwarzschild spacetimes to test out our method for finding horizons. To get the spacetime metric for a boosted black hole consider $`\overline{𝒪}`$ to be the rest frame of the black hole, with coordinates $`(\overline{t},\overline{x}^i)`$. Let $`𝒪`$ be another stationary frame with coordinates $`(t,x^i)`$ such that $`𝒪`$ is related to $`\overline{𝒪}`$ via a Lorentz boost along the $`\widehat{𝐯}=(\widehat{v_x},\widehat{v_y},\widehat{v_z})`$ direction: in the $`𝒪`$ frame the black hole moves in the $`\widehat{𝐯}`$ direction with boost velocity, $`v`$ ($`\delta _{ij}\widehat{v}^i\widehat{v}^j=1`$). As usual, we define $`\gamma =1/\sqrt{1v^2}`$. $`H(x_{\overline{\mu }})`$ and $`l_{\overline{\mu }}`$ (bar denoting $`\overline{𝒪}`$ frame) now transform as $`H(x_\mu )=H(\mathrm{\Lambda }_\mu ^{\overline{\nu }}x_{\overline{\nu }})`$ (21) and $`l_\mu =\mathrm{\Lambda }_\mu ^{\overline{\nu }}l_{\overline{\nu }}(\mathrm{\Lambda }_\gamma ^{\overline{\sigma }}x_{\overline{\sigma }}).`$ (22) These preserve the form of (12). ### A 3+1 Decomposition The standard ADM 3+1 form of the spacetime metric is given by $$ds^2=\alpha ^2dt^2+\gamma _{ij}\left(dx^i+\beta ^idt\right)\left(dx^j+\beta ^jdt\right)$$ (23) If we compare (12) to (23) and use the property that $`l^\mu l_\mu =0`$, we find that the lapse is given by $`\alpha `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+2Hl_t^2}}},`$ (24) and the shift is given by $$\beta _i=2Hl_tl_i$$ (25) or $$\beta ^i=2Hl_t\delta ^{ij}l_j/(1+2Hl_t^2).$$ (26) The 3-metric is given by $$\gamma _{ij}=\eta _{ij}+2Hl_il_j.$$ (27) as expected and the extrinsic curvature is determined from $`K_{ij}`$ $`=`$ $`_t\gamma _{ij}/2\alpha +D_i\beta _j+D_j\beta _i`$ (28) $`=`$ $`_t(Hl_il_j)/\alpha +2\left(D_i(Hl_tl_j)+D_j(Hl_tl_i)\right)`$ (29) and $`\gamma ^{ij}`$ $`=`$ $`\delta ^{ij}2H\delta ^{il}\delta ^{jk}l_ll_k/(1+2Hl_t^2).`$ (30) Note that $`det\gamma _{ij}`$ $`=`$ $`1+2Hl_t^2.`$ (31) To obtain the 3-metric and extrinsic curvature we need to specify $`H`$, $`l_\mu `$, $`_\mu H`$ and $`_\mu l_\nu `$ and substitute these into Eq.(27) and Eq.(29). In order to evaluate $`_\mu l_\nu `$ and $`_\mu H`$ a specific choice of spacetime has to be made. For example for the Kerr spacetime we take the expressions for $`H`$ and $`l_\mu `$ from Eq.(14) and Eq.(13), compute their derivatives and substitute. This gives us a “Kerr-Schild” slice of the Kerr spacetime. ## IV Results In the presentation that follows we first conduct a series of basic tests of the algorithm using metric and extrinsic curvature data just presented. In the second part of this section we set up a 3-dimensional Cartesian grid, $`\widehat{\mathrm{\Sigma }}`$, of $`n^3`$ points, on which we define a coordinate system where the black hole (either Kerr or Schwarzschild) is placed at the origin. Using Eq.(27,29) again we generate $`\gamma _{ij}`$ and $`K_{ij}`$ on $`\widehat{\mathrm{\Sigma }}`$ everywhere but the region that contains the curvature singularity (for Schwarzschild at $`x^2+y^2+z^2=0`$ and Kerr $`\rho =\sqrt{x^2+y^2+z^2}a`$). With the data on $`\widehat{\mathrm{\Sigma }}`$, we start the apparent horizon locator with an initial guess surface which is a 2-sphere of radius $`r_0`$. The horizon locator surface mesh sizes used in this calculation are $`N_s=33,41,49,65,81,97`$. The stopping criterion for the Newton iterations was determined empirically to be $`10^9`$. The cases we present are (1) $`v=0`$, $`a=0`$ (unboosted Schwarzschild), (2) $`v=0`$, $`a0`$ (unboosted Kerr), (3) $`v0`$, $`a=0`$ (boosted Schwarzschild) and (4) $`v0`$, $`a0`$ (boosted Kerr). ### A Tests with Eddington-Finkelstein metric data First, with $`\widehat{v}_x=0`$, $`\widehat{v}_y=0`$, $`\widehat{v}_z=0`$ and $`a=0`$ (unboosted black hole in ingoing Eddington-Finkelstein coordinates) we show some basic tests of the apparent horizon locator. Most importantly we show that solutions obtained with our locator are $`O(h^2)`$. With these data all components of $`\gamma _{ij}`$ and $`K_{ij}`$ are non-zero. The latter property makes this a good initial model problem to work with, because the computation is fully exercised in an analytically tractable situation. As stated earlier the apparent horizon is expected to be located at $`r=2M`$. These tests are conducted with data specified analytically where required. #### 1 Residual Evaluation and Second Order Convergence We place the black hole at the origin of the computational domain($`x=0,y=0,z=0`$). In spherical symmetry for this metric the apparent horizon equation becomes the algebraic equation, $$F(r)=\frac{12M/r}{r\sqrt{1+2M/r}}=0.$$ (32) A plot of $`F(r)`$ is shown in FIG.(1). At $`r=2M`$ we have $`F=0`$. A useful test of the evaluation of the expansion of the outgoing normals $`F(r)`$ is to see if indeed the residual $`\widehat{F}[\widehat{\rho }]`$ is correctly evaluated to $`O(h^2)`$ as $$\widehat{F}=F+e_2h^2+\mathrm{},$$ (33) where $`e_2h^2`$ is the leading order truncation error term. Given that the exact value is known for $`\widehat{F}[\widehat{\rho }]`$ we can approximate the leading order truncation error. We carry out a convergence test by evaluating $`\widehat{F}[\widehat{\rho }]`$ on a 2-sphere of $`r=2M`$ for a series of mesh sizes, $`N_s=17,25,33,49,65,96,129`$. We examined the behavior of $`\mathrm{log}\widehat{F}_2`$ (where this is the $`L_2`$ norm) versus $`\mathrm{log}N_s`$, where $`N_s`$ is the number of mesh points on one side of the $`N_s\times N_s`$ mesh. At $`r=2M`$ $`\widehat{F}e_2h^2`$ and so the $`L_2`$-norm, $`\widehat{F}_2e_2_2h^2`$. Since $`h1/N_s`$ we expect that if the residual is $`O(h^2)`$ then the slope of a plot of $`\mathrm{log}\widehat{F}_2`$ versus $`\mathrm{log}N_s`$ should be $`2.0`$, which we validated via a least squares fit. A closely related test is to also evaluate $`\mathrm{log}\widehat{\rho }\rho _2`$ versus $`\mathrm{log}N_s`$, where $`\widehat{\rho }`$ is the numerical solution from the apparent horizon locator and $`\rho `$ is the exact horizon location. FIG.(3) shows the result. From a least squares fit to a straight line the slope is found to be $`2.1`$ which validates our solution as $`O(h^2)`$. #### 2 Jacobian For the same 2D mesh discussed we generate the Jacobian matrix for a single Newton step. FIG.(2) shows the structure of the matrix for a $`33\times 33`$ run. There are $`1025`$ independent points on $`\widehat{S}`$ and hence $`J_{\overline{\mu }\overline{\nu }}`$ is a $`1025\times 1025`$ matrix. The dots in the figure are non-zero Jacobian entries. There are seven bands in this matrix with 2 additional ones in the vicinity of the poles at $`\overline{\mu }=1`$ and $`\overline{\mu }=1025`$. The structure reflects the domain of dependence of the finite difference operators used in the evaluation of $`\widehat{F}`$. Here it comes from a combination of interpolations and Cartesian finite differencing. Near $`\overline{\mu }=1`$ and $`\overline{\mu }=1025`$ the additional bands come from our special choice of interpolation stencils at the poles, as discussed in the appendix. The structure reflects the fact that a perturbation at a single mesh point affects the residual in a small neighborhood around it so we can optimize the generation of the Jacobian to $`O(N)`$ by evaluating $`\widehat{F}[\rho +\delta \rho ]`$ only in a small neighborhood of the perturbed point. The Jacobian generation was found to be order $`O(N^p)`$ where $`p=1.08`$ and $`N`$ is the total number of independent points on the 2D computational mesh. A matrix $`A`$ is defined to be diagonally dominant if its elements, $`A_{ij}`$, satisfy $$\underset{\stackrel{j=1}{ji}}{\overset{n}{}}A_{ij}A_{ii}\text{for all i}.$$ (34) We found that the Jacobian is not diagonally dominant since the inequality in Eq.(34) is not satisfied for all $`i`$ and $`j`$. This is of interest since for some iterative solution techniques (Gauss-Seidel and SOR for example) a sufficient condition for the solution of a linear system, $`Ax=b`$, is that the matrix, $`A`$, be diagonally dominant. In our case we concluded from early experiments that indeed such simple iterative solvers did not converge for this problem. The Jacobian matrix is not symmetric but it is well-conditioned for the spacetimes that we have considered. For a $`33\times 33`$ run the Jacobian has a condition number of $`\kappa `$ of about $`10^4`$ to $`10^5`$ where, $$\kappa =AA^1,$$ (35) (In their definition of the condition number, Dongarra et al. use the $`L_1`$-norm.) The condition number tells us how close the matrix $`A`$ is to being singular. A very large condition number or a reciprocal condition number close to machine epsilon tells us that $`A`$ is singular. An identity matrix has a condition number of $`1`$. To estimate $`\kappa `$ we used the LINPACK library routine, DGECO. #### 3 Solution of the Linear system As stated before to locate the apparent horizon using our technique we have to obtain a solution, $`\delta \widehat{\rho }`$, to the linear system $`\widehat{J}\delta \widehat{\rho }=\widehat{F}[\widehat{\rho }]`$ (36) which is the discrete form of Eq.(10). Since the properties of the matrix do not allow us to use the standard iterative methods such as Jacobi and Gauss-Seidel methods, We use a modified conjugate-gradient method due to Kershaw (The standard form of the Conjugate gradient method will not work since $`J`$ is not symmetric.) Kershaw’s method, termed the Incomplete LU-Conjugate gradient method (ILUCG), can solve any linear system, $`Ax=b`$, with $`A`$ being any nonsingular, sparse matrix. The method involves preconditioning the matrix via an incomplete LU decomposition. This method has worked quite well for our purposes. In principle other schemes for solving the resulting system can be used. #### 4 Solution for the apparent horizon location Using the ILUCG method to solve for $`\delta \widehat{\rho }`$ we demonstrate the Newton solver’s ability to sucessfully locate apparent horizons in the Eddington-Finkelstein metric data. The apparent horizon in this data is a 2-sphere of radius $`2M`$. Using a 2-sphere of radius $`r_0`$, centered on the origin as the initial surface we carried out a series of runs for $`r_0=\mathrm{0.5..3.0}`$ with a $`33\times 33`$ mesh. Table (I) shows the radius, $`r_0`$, of the initial starting 2-sphere, the number of Newton iterations taken if it converged, the final residual value and the solution error. The stopping criterion is that the norm of the change in the solution, $`\delta \widehat{\rho }_2`$ be less than $`10^{10}`$. We see that for all the cases, provided the solver managed to drive $`\delta \rho _2`$ below $`10^{10}`$ the final percentage error remains fixed; what differed in each case was the number of iterations taken and rate of convergence. Once the solver has driven $`\widehat{\rho }`$ into the vicinity of the solution the Newton convergence is quadratic. This happened in the 6th iteration in the $`r_0=0.5`$ run. For the $`r_0=2.0`$ run the solver took four iterations to converge down to the stopping criterion. In another series of runs with $`N_s=33`$, $`M=1`$ and $`r_0=2.5`$ the perturbation parameter, $`ϵ`$, was varied from $`10^1`$ to $`10^6`$. An optimum value of $`ϵ`$ for this metric data was found to be about $`10^4`$ to $`10^6`$. In general the stopping criterion need not be as stringent as we have set it. In numerical spacetimes where the metric data will have truncation error associated with them the truncation error of $`F`$ is expected to be much larger than our test stopping criterion. In that case a larger stopping criterion should to be chosen to avoid wasting computational effort. The optimum value is dependent on the error in the metric data. In the results we present in the paper however we drive the residuals down as far as possible. On the other hand, a perturbation parameter $`ϵ`$ must be chosen such that $`\widehat{F}[\widehat{\rho }+ϵ]\widehat{F}[\widehat{\rho }]`$ is sufficiently large that this expression is not dominated by truncation error. #### 5 Numerical Metric data Tests described so far used data analytically computed at each point as needed. Since the ultimate goal is to incorporate this apparent horizon location algorithm into an evolution code it is useful to gauge the performance of the algorithm with numerical metric data and with the data structures expected in the real application, where, for instance, part of the domain is excised from consideration. Thus we set up the same Eddington-Finkelstein data on a 3D Cartesian grid of $`n^3`$ points, with a region of this grid excised to emulate the situation in an evolution code where the interior of the black hole is excluded. The apparent horizon surface which is embedded in this 3D Cartesian grid typically does not lie on Cartesian grid points and as a result an interpolation tool is required. If the surface mesh, during the course of the Newton iterations, overlaps the excised region then extrapolation is required. We make use of an interpolator/extrapolator written by S. Klasky to obtain the 3-metric, extrinsic curvature and the spatial derivatives of the 3-metric at any point. The use of an interpolator brings in truncation errors associated with the interpolation/extrapolation operations. In the following we show that even with extrapolation errors the solver works quite well in locating apparent horizon surfaces. We set up an uniform 3D Cartesian grid, $`\widehat{\mathrm{\Sigma }}`$, of size $`n^3`$. On this grid we excise a region interior to a sphere of radius $`R_m`$ centered at $`(x_m,y_m,z_m)`$ so that the metric data is defined for $`r>R_m`$ and undefined for $`r<R_m`$ where r is the Cartesian distance in Kerr-Schild coordinates from the excision center. In the following discussion on radial and offset apparent horizon locations we take $`n=65`$ for the Cartesian grid (with $`h=1/8`$) and $`N_s=33`$ for the surface mesh. The stopping criterion used in the horizon finder is $`\beta =10^4`$. That is, if $`\delta \rho _2<\beta `$ then the Newton iterations are stopped. The perturbation parameter, $`ϵ`$, is taken to be $`10^4`$. The interpolator tool is used to fourth order. The initial guess surface used is a sphere of radius $`r_0=2.1M`$ centered at the origin of the Cartesian grid. With these parameters we carry out two set of tests. The first is a radial test of the horizon locator with the use of the interpolator and the second is an offset test. These tests examine the effect of extrapolation of metric data on the residual, $`\widehat{F}`$, and solution to the apparent horizon equation. #### 6 Apparent horizon location (Radial tests with excision) In the first test case we center the black hole at $`(0,0,0)`$. The masked region is also centered at $`(0,0,0)`$. We carry out a series of tests with the excision radius, $`R_m`$, varying from $`1.5`$ to $`2.6`$. Thus the apparent horizon is in the defined region ($`R_m<2M`$) for some of the tests, and for others it is inside the excised region ($`R_m>2M`$). This provides evidence of the effect of extrapolations on the residual of the apparent horizon equation, $`\widehat{F}`$, and the error in its solution. FIG.(4) illustrates the behaviour of the $`L_2`$-norm of the residual, $`\widehat{F}_2`$, as a function of $`R_m`$. FIG.(5) shows the percentage relative error of the solution of the apparent horizon equation as a function of $`R_m`$. The percentage relative error is calculated from the exact solution for $`r=\rho (\theta ,\varphi )=\overline{\rho }=2M`$. With this exact solution, we calculate the percentage relative error, $`e\rho \overline{\rho }/\overline{\rho }\times 100\%`$. For $`R_m<2M`$ the derivative interpolator/extrapolator uses interpolation for regions near the apparent horizon location ($`r=2M`$), while for $`R_m>2M`$ it uses extrapolation. As $`R_m`$ increases further the errors due to extrapolation increase, as expected. This can be seen in FIG.(4) where $`\widehat{F}_2`$ increases quickly for $`R_m2.4`$, as does the error shown in FIG.(5). At $`R_m=2.5M`$, the solver could not bring $`\delta \rho _2`$ down to below $`10^4`$, and so failed to meet the stopping criterion. This can be understood in terms of the Cauchy-Schwarz inequality , Since $`J\delta \rho =\widehat{F}`$ we have that $$\delta \rho \frac{\widehat{F}}{J}.$$ (37) At $`R_m=2.2`$ where $`\widehat{F}10^3`$ and $`\delta \rho 10^5`$, we have from Eq.(37) that $`J10^2`$. Therefore at $`R_m=2.5M`$ we expect with $`\widehat{F}10^2`$ that $`\delta \rho 10^4`$. By relaxing the criterion past $`R_m=2.5`$ we can still obtain a solution. Past $`R_m=2.6`$ the convergence progressively worsens. For example, at $`R_m=2.9`$, $`\widehat{F}`$ could not be brought below $`10^3`$, and the solution error is 5%. The amount of error sustained from interpolation of the metric data is dependent on the resolution of the Cartesian grid and the behaviour of the functions being interpolated. If the gradients of $`\gamma _{ij}`$ and $`K_{ij}`$ are large near the horizon then a larger interpolation error is sustained. This in turn leads to a larger truncation error in $`F`$. In the numerical evolution of black hole spacetimes with excision then buffer zones may not be necessary for the location of apparent horizons. However, for other reasons buffer zones might be necessary. #### 7 Locating Offset apparent horizons We examine behaviour of the locator with the derivative interpolator for a black hole offset so that it overlaps the excised region. This is important in tracking moving black holes. The center of the masked region is at $`(0,0,0)`$ and the black hole of radius 2M is centered at ($`\delta /\sqrt{3},\delta /\sqrt{3},\delta /\sqrt{3}`$), so that the radial distance between the mask center and the hole is $`\delta `$. With a grid spacing of $`h=1/8`$, an offset of $`\delta =1`$ corresponds to approximately 8 grid zones. FIG.(6) shows the percentage relative error in the apparent horizon location as a function of the offset $`\delta `$. As the graph illustrates, up to $`\delta =0.7`$ the percentage relative error is below one percent. (At $`\delta =0.7`$ the percentage error is 0.6%.) From $`\delta =0.7`$ onwards, however the solver becomes sensitive to initial conditions and extrapolation errors and quickly ceases to converge. At $`\delta =0.7`$, about 5-6 grid points offset, we are still able to find horizons. Generally in explicit time-evolution codes the CFL condition restricts the black hole motion from one time slice to another, to be less than one zone ($`\delta <h`$ or about $`\delta 0.1`$ in our test case). Hence we expect, based on the results for our model spacetime as shown in FIG.(6), that in such an evolutionary scheme with a similar resolution we will be able to track black hole apparent horizons to less than 0.1%. ### B Apparent Horizons in Boosted Kerr data In this section we now focus on apparent horizon location for boosted Kerr and Schwarzschild black holes. For the data that follow we excise a 2-sphere of radius $`r>a`$ centered about the origin from the computational domain to avoid the ring singularity structure of the Kerr black hole. Using the interpolator tool in conjunction with the apparent horizon locator, we locate horizons for various values of the angular momentum parameter, $`a`$. The horizon locator begins with a trial surface which is a 2-sphere of radius $`r_0=2M`$. The locator was run for $`a=0.0,0.1,0.2,\mathrm{},0.9`$ at $`N_s=33`$. FIG.(7) shows a cross-section of the horizon in the $`xz`$ plane as a function of $`a`$. The apparent horizon is seen to have the shape of an oblate spheroid. In the runs used to generate these data we used $`ϵ=10^5`$ and a stopping criterion that ensured that the $`l_2`$-norm of $`\widehat{F}`$ on the computational mesh was less than $`10^{11}`$. FIG.(8) shows the $`l_2`$-norm of the error in the solution, $`(\widehat{r}r_+)/r_+_2`$, versus mesh size. This set of runs was carried out with $`a=0.9`$ and $`N_s=17,25,33,49,65`$. Where $`r_+`$ is given by (19) and the $`\widehat{r}`$ is computed from $$\widehat{r}\left\{\frac{1}{2}\left(\widehat{\rho }^2a^2\right)+\left[\frac{1}{4}\left(\widehat{\rho }^2a^2\right)^2+a^2z^2\right]^{1/2}\right\}^{1/2}$$ (38) where $`\widehat{\rho }=\sqrt{x^2+y^2+z^2}`$ is the solution from the apparent horizon locator. From a least squares fit the slope is found to be $`2.1`$ and the solution is $`O(h^2)`$. The area of the event horizon in the Kerr spacetime is given by $`𝒜`$ $`=`$ $`4\pi \left(r_+^2+a^2\right)`$ (39) Let $`\widehat{𝒜}`$ be the computed apparent horizon area. FIG.(9) shows the percentage errors $`(\widehat{𝒜}𝒜)/𝒜\times 100\%_2`$ versus $`a`$ for various resolutions $`N_s=17,\mathrm{},65`$. The area of the apparent horizon is computed via a technique which projects the 3-metric, $`\gamma _{ij}`$, onto the 2-surface to obtain an area element $`\sqrt{{}_{}{}^{(2)}\gamma }d\theta d\varphi `$, and then computes the surface integral. FIG.(9) shows the percentage errors in the area for increasing resolution. We now consider Schwarzschild and Kerr black holes boosted in the yz-direction. That is, we look at $`\widehat{v}_x=0,\widehat{v}_y=1/\sqrt{2},\widehat{v}_z=1/\sqrt{2}`$ and $`a=0,0.9`$. In each of these cases we locate apparent horizons for $`v=0,0.1,\mathrm{},0.9`$. From $`v=0`$ to $`v=0.8`$ we started with a two-sphere of radius $`2M`$ and found an apparent horizon with outgoing expansions driven down to $`10^{12}`$. For $`v>0.8`$ we had difficulty driving the expansions down. As a result we utilized the solution at $`v=0.8`$ as an initial guess and were subsequently able to find horizons by stepping every $`0.25`$ from $`v=0.8`$ to $`v=0.9`$. We used $`ϵ=10^5`$ again for these runs. At $`v=0`$ the initial guess is the apparent horizon and there within the six Newton iterations the expansions were driven down around $`10^{12}`$. The first Newton iteration took the expansions down around $`10^6`$. For $`v=0.5`$ starting from an initial guess of a sphere of radius $`\rho =2M`$ it took four Newton iterations to drive the expansions down around $`10^6`$ and nine Newton iterations to get down to $`10^{12}`$. Typically in a numerical time-evolution of such a spacetime we would not need to drive the expansions down to this level. If we are utilizing a surface within the apparent horizon as an excision boundary then we need only to drive the expansions down far enough to be certain an apparent horizon is present. FIG.(10) shows the yz-cross-section of the apparent horizon for various boost velocities compared against an unboosted black hole apparent horizon cross-section. We find that the apparent horizon is Lorentz contracted in the yz-direction in the boosted coordinates. We have considered a slice of such a boosted spacetime in which the event horizon appears Lorentz contracted in the resulting coordinates. We know that in these spacetimes the apparent horizon should coincide with the event horizon and we find that this is indeed the case. First, the area of the apparent horizon coincides with the area of the event horizon which is invariant under a boost. FIG.(11) shows the error in the apparent horizon area as a function of $`v`$ for various resolutions. We find that with increasing resolution the error in the area converges towards zero. This demonstrates that the area of the apparent horizon is invariant under a Lorentz boost. This is coupled to an interesting property of the Kerr-Schild type of metrics that $`r=r_+`$ remains fixed for Kerr and Schwarzschild black holes. This is illustrated in FIG.(12) where we show the error in the radial coordinate $`r=2M`$ on the apparent horizon for various boost velocities. In this case the black hole is boosted in the $`xyz`$-direction for generality. That is, $`\widehat{v}_x=1/\sqrt{3},\widehat{v}_y=1/\sqrt{3},\widehat{v}_z=1/\sqrt{3}`$ and $`a=0`$. Here $`r`$ is computed from the boosted coordinates. We find that $`r`$ converges towards $`2M`$ for increasing resolution satisfying yet another property of the boosted Kerr-Schild spacetime. FIG.(13) shows surface plots of the apparent horizon for $`v=0,0.3,0.6`$ and $`v=0.9`$ displayed in Kerr-Schild Cartesian coordinates. Note how distorted the apparent horizon gets with increasing boost velocities. As seen in the figures for the $`yz`$-boosts the boosted apparent horizons in this case are always contained within the apparent horizon for $`v=0`$. That is, the boost contracts the apparent horizon in the boost-direction. Again for a boost velocity of $`v=0.5`$ it took the solver eight Newton steps to drive the expansions down around $`10^{12}`$. The stopping criterion used in this run was $`10^{12}`$ and the final expansions are $`10^{13}`$. On average it took four Newton steps to drive the expansions down to $`10^6`$ and nine Newton steps to $`10^{12}`$ starting from an expansion of $`0.1`$. In the case of a boosted Kerr black hole with $`a=0.9`$ the results are again very similar to those of the Schwarzschild black hole. Note that now with $`a=0.9`$ and $`v0.9`$ we get even more distorted apparent horizons. These results show that this algorithm for finding apparent horizons does quite well with such large distortions. In addition the cost of finding these surfaces increases by only two additional Newton steps. FIG.(14) shows the $`yz`$-cross-sections for the apparent horizon found for $`a=0.9`$ as a function of $`v`$. Again the boosted apparent horizon is contained within the unboosted one and Lorentz contracted. FIG.(15) shows the error in the area for the same data. With $`a=0.9`$ we expect that the area should be $`36`$. The graph shows that for increasing resolution the error in the area tends towards zero. Hence the area remains fixed with increased boost velocity as is expected. Similarly, $`r`$ computed from the boosted coordinates remains fixed at $`r_+`$ as is shown by FIG.(16). The apparent horizons found here were obtained with $`\widehat{v}_x=1/\sqrt{3},\widehat{v}_y=1/\sqrt{3},\widehat{v}_z=1/\sqrt{3}`$ and $`a=0.9`$. That is, the boost was in the $`xyz`$-direction with magnitude $`v`$. Again we find that $`r`$ on the apparent horizon converges to $`r_+1.4`$ with increasing resolution for all boost velocities. At a resolution of $`33\times 33`$ we have an percentage error of $`8\%`$ and $`1\%`$ at $`65\times 65`$. FIG.(17) shows surface plots of the apparent horizon for the boosted Kerr black hole for $`v=0,0.3,0.6`$ and $`0.9`$. Note how distorted the final apparent horizon surface is. Our algorithm required one more Newton step to drive the expansions down to $`10^{13}`$ for $`v=0`$ compared to $`v=0.9`$. It took six Newton iterations to drive the resolution from about $`0.2`$ at the inital step to $`10^6`$ for both boost velocities. Hence, this algorithm has the advantage that given sufficient resolution on the computational mesh, the work done does not drastically increase for increasing distortions. ## V Discussion We have demonstrated in this paper that our method based on finite difference techniques is viable for locating very distorted boosted Kerr black hole apparent horizons. We have shown that the located horizons obey the expected analytical rule, of invariance of the area of the event horizon, in cases corresponding to at-rest or boosted single black holes, where the apparent horizon is known to coincide with the event horizon. We have additionally given a number of computational tests demonstrating the behavior of the tracker on interpolated or extrapolated data which is realistically similar to that from evolutions. In other contexts algorithm has been thoroughly tested with the canonical set of test problems such as the two and three black hole initial data sets and additionally in an evolution code tracking the apparent horizon for a Schwarzschild black hole in geodesic slicing and demonstrated to be capable of tracking apparent horizons in boosted Schwarzschild data. Those tests and the tests given here show its viability as a method for locating black hole apparent horizons and using them for black hole excision. Since black hole excision is essential for long-term evolutions of single or multiple black hole spacetimes. It is very useful to have efficient apparent horizon locators that can locate apparent horizons “fast” relative to the time taken for an evolution time step. Our algorithm is dominated primarily by computations of the Jacobian matrix in the use of Newton’s method. These operations are optimized such that they scale as $`𝒪(N)`$ approximately where $`N`$ is the total number of points on the two-dimensional mesh used for the solution. With a $`33\times 33`$ mesh we find that each Newton iteration takes on average $`20`$ Origin 2000 CPU seconds. This is independent of the distortion of the apparent horizon. However the number of Newton iteration steps is determined by the “distance” of the initial starting surface from the final solution. During the course of an evolution it is expected that the apparent horizons over several timeslices will be “close” enough to each other that two to three Newton iterations will be sufficient to locate the horizon at low accuracy with the expansion of the outgoing null rays on its surface being at the level of $`10^5`$ or $`10^6`$. Obtaining a better accuracy requires more Newton iterations and the number of timesteps taken depends also on the accuracy of the background metric data. Typically in our model problems eight iterations will take us below $`10^{10}`$. One of the drawbacks of a Newton’s method for finding apparent horizons is its sensitivity to the initial guess. An initial guess outside of the radius of convergence will not lead to a solution. Additionally Newton methods are known to be sensitive to high frequency components in the solution. This is demonstrated in axisymmetry by Thornburg. Sensitivity to the initial guess can be easily handled by combining the Newton method algorithm with apparent horizon trackers that are based on flow methods. The flow finder is used to obtain an initial guess for the Newton method which then converges on the solution very quickly. The efficiency of our boundary-value method can be compared to the efficiency of other approaches (variations of flow mothods) due to Tod, as developed by Shoemaker et al. and fast flow methods developed by Gundlach. The flow method is based on a parabolic partial differential equation whose rate of convergence to the solution slows as it approaches it. Typically for a $`33\times 33`$ run the flow method takes on the order of thousands of seconds to converge down to expansions of $`10^4`$. The advantage of this method however is its ability to find multiple apparent horizons from an arbitrary initial guess. Combined with the Newton finder this will result in a robust apparent horizon finding scheme. We can also compare the effort to spectral decomposition methods. We concentrate on a method similar to that of Nakamura et al. in which the equation for the apparent horizon surface is written: $`\rho (\theta ,\varphi )`$ $`=`$ $`{\displaystyle \underset{l=0}{\overset{l_{max}}{}}}{\displaystyle \underset{m=l}{\overset{l}{}}}a_{lm}Y_{lm}(\theta ,\varphi ).`$ (40) We do not have access to an apparent horizon finder based on pseudo-spectral methods but we will analytically compute the coefficients for the case of a boosted Schwarschild black hole; this will give some insight into the range of harmonics required, and some idea of what scaling of these methods might be. In Kerr Schild coordinates, the hole, with boost in the z-direction, has the shape of a spheroid, $`{\displaystyle \frac{x^2}{a^2}}+{\displaystyle \frac{z^2}{b^2}}=1,`$ (41) where we have supressed the y-direction. Notice that the axes $`a`$,$`b`$ of the ellipsoid obey $`b^2/a^2=1v^2`$, which demonstrates that the eccentricity is directly proportional to the boost velocity, $`ϵ=v`$, for this case. Hence even ellipses with moderate ratio of axes, such as that for v=0.9, where the ratio is a little less than 0.5, have moderately large eccentricities. We will approximate the form Eq.(41) with an axisymmetric series of the form Eq.(40) (the general case would have nonaxisymmetric terms also). We find it more convenient to work with Legendre polynomials than with the spherical harmonics directly. Since we work with Legendre polynomials, we drop the y- coordinate in the spheroid expression, to obtain : $`R(\theta )`$ $`=`$ $`b/\sqrt{1ϵ^2\mathrm{sin}^2\theta }`$ (42) $`=`$ $`b/\sqrt{1ϵ^2(1q^2)},`$ (43) where $`q=\mathrm{cos}\theta `$. To obtain the expansion of expression Eq.(43) in terms of $`P_m`$, we first expand using the binomial theorem. $`R(\theta )/b`$ $`=`$ $`{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}\left(\begin{array}{c}1/2\\ s\end{array}\right)(1)^sϵ^{2s}(1q^2)^s.`$ (46) This converges for all $`v<1`$. Using the binomial theorem again for $`(1q^2)^s`$ we substitute $`(1q^2)^s`$ $`=`$ $`{\displaystyle \underset{r=0}{\overset{s}{}}}\left(\begin{array}{c}s\\ r\end{array}\right)(1)^rq^{2r}`$ (49) in (46) to obtain $`R(\theta )/b`$ $`=`$ $`{\displaystyle \underset{s=0}{\overset{\mathrm{}}{}}}\left(\begin{array}{c}1/2\\ s\end{array}\right)(1)^sϵ^{2s}{\displaystyle \underset{r=0}{\overset{s}{}}}\left(\begin{array}{c}s\\ r\end{array}\right)(1)^ra_{sr}`$ (54) where $`a_{sr}(q)`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{r}{}}}{\displaystyle \frac{2^{2n}(4n+1)(2r)!(r+n)!}{(2r+2n+1)!(rn)!}}P_{2n}(q)`$ (55) and we made the substitution $`q^{2r}`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{r}{}}}{\displaystyle \frac{2^{2n}(4n+1)(2r)!(r+n)!}{(2r+2n+1)!(rn)!}}P_{2n}(q).`$ (56) By exchanging the summations over $`r`$ and $`n`$ and then $`n`$ and $`s`$ it is possible to rewrite Eq.(54) as $`r(\theta )={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}C_{2n}P_{2n}(\theta ),`$ (57) where $`C_{2n}`$ $`=`$ $`2^{2n}(4n+1){\displaystyle \underset{s=n}{\overset{\mathrm{}}{}}}\left(\begin{array}{c}1/2\\ s\end{array}\right)(1)^sϵ^{2s}\times `$ (63) $`{\displaystyle \underset{r=n}{\overset{s}{}}}\left(\begin{array}{c}s\\ r\end{array}\right){\displaystyle \frac{(1)^r(2r)!(r+n)!}{(2r+2n+1)!(rn)!}}.`$ FIG.(18) gives the coefficients $`C_{2n}`$ for $`n=1,\mathrm{},10`$, and for several values of v. While FIG.(18) shows the exponential convergence of algorithm with $`n`$, it also shows that the coefficient of the convergence is small for $`v0.9`$. It can be seen that the number of required terms approaches 20 for $`v=0.9`$ if the error is required to be less than $`10^3`$. (The general sum would have polynomials of odd as well as even order, and for each l, a set of azimuthal quantum numbers spanning -2n to 2n). Hence in general, to compute the distorted apparent horizon would take a search over $`20^2`$ parameters in a minimization routine. This is equivalent to inverting a full matrix of this size, and would be expected to be slow. The boundary value problem is expected to be much faster. It is a fact that the boundary value problem as now implemented does not handle multiple black-hole spaces, so its speed is counteracted by the impossibility of using it in 2-hole cases. However, it may be possible to use a flow method which does recognize the existence of seperate black holes, run down to find the two holes each with some accuracy, and then to use this boundary soluion code to quickly get a highly accurate result. We are confident such a combined tool would be of great utility. ## VI Acknowledgements We wish to thank S.Klasky for providing and developing the 3d interpolation tools used in a part of this work as well as for discussions and work on applications of interpolation with excision. MFH wishes to thank J.Thornburg for stimulating discussions on apparent horizon location during the early part of this work. This work was supported by NSF PHY9800722 and PHY9800725 and PHY9800970. ## VII Appendix A: Evaluation of the residual The evaluation of Cartesian derivatives on $`\widehat{𝒮}`$ is carried out by constructing 3D finite difference stencils at each mesh point on $`\widehat{𝒮}`$. The finite difference stencil, denoted by $`𝒩`$, consists of 26 additional points around each mesh point. These 26 points are $`\pm \delta x`$, $`\pm \delta y`$ and $`\pm \delta z`$ away from the central mesh point as shown in FIG.(19). These points, as shown, are organized into three planes of constant $`z`$: $`z=z_0\delta z,z_0,z_0+\delta z`$. Each plane contains the nine nearest neighbors to the center point, including the center point itself in the case of $`z=z_0`$. We use a single discretization scale $`h`$ ($`\delta x=\delta y=\delta z=h`$) which is always proportional to the mesh spacing $`\delta \theta =\pi /(N_s1)`$. To define $`\phi (x,y,z)`$ at each stencil point $`𝐱𝒩`$ we use its split into radial and angular parts, $`\phi (x,y,z)=r\rho (\theta ,\varphi )`$. For each stencil point $`𝐱`$ we compute the corresponding spherical coordinates $`(r_x,\theta _x,\varphi _x)`$. This point can be thought of as a ray emanating from the origin of our spherical coordinate system (which coincides with the origin of our Cartesian coordinate system) along $`(\theta _x,\varphi _x)`$ of length $`r_x`$. FIG.(20) labels the point $`𝐱`$ as P. The dashed line from P to the origin is the ray from the origin. Its intersection with $`\widehat{𝒮}`$ is denoted by a filled square. The value of $`\phi `$ at $`𝐱`$ can be obtained by computing $`\rho (\theta _x,\varphi _x)`$ via biquartic interpolation where the truncation error has a leading order term which is fourth order in the grid spacing $`h`$. The interpolation is carried out with values of $`\widehat{\rho }`$ defined on mesh points of $`\widehat{𝒮}`$ using a 16 point stencil. FIG.(21) shows the choice of these stencil points in the interior of the mesh. At the poles a special choice is made of stencil points which takes into account the indentifications made at the poles. FIG.(22) shows a choice of stencil points for an interpolation point near the pole. This approach leads to a fourth order truncation error in $`\rho (\theta _x,\varphi _x)`$ at all points on $`\widehat{𝒮}`$. Then $`\phi `$ can be constructed for every $`𝐱𝒩`$ as $`\phi =r_x\rho (\theta _x,\varphi _x)`$. Using this approach $`\phi `$ is defined at any finite difference stencil point for every mesh point on $`\widehat{𝒮}`$. The finite difference expressions for $`\mathrm{\Delta }_i^h\phi `$, $`\mathrm{\Delta }_i^h\mathrm{\Delta }_j^h\phi `$ (corresponding to first and second derivatives) are computed at each of the mesh points. The residual is then evaluated on $`\widehat{𝒮}`$ using these finite difference approximations for the derivatives to $`O(h^2)`$, and metric data ($`\gamma _{ij}`$, $`_k\gamma _{ij}`$, $`K_{ij}`$) which are specified either analytically or interpolated from an enveloping 3D Cartesian grid. Because we use $`O(h^2)`$ finite difference approximations $`\mathrm{\Delta }_i^h\phi `$, $`\mathrm{\Delta }_i^h\mathrm{\Delta }_j^h\phi `$ to the derivatives, this approach leads to an $`O(h^2)`$ truncation error in evaluating $`\widehat{F}[\widehat{\rho }]`$ . Because of our special attention to points near the pole, $`\widehat{F}`$ is evaluated smoothly everywhere on $`\widehat{𝒮}`$. With a means for evaluating $`\widehat{F}`$ at any point in the domain of $`\widehat{S}`$ it is straightforward to generate $`\widehat{J}_{\overline{\mu }\overline{\nu }}`$ numerically using Eq.(11). The algorithm for this is summarized as follows: \| Specify metric data everywhere on $`\widehat{𝒮}`$ \| \| --- \| \| Evaluate $`\widehat{F}[\rho ]`$ everywhere on $`\widehat{𝒮}`$ \| \| For each point(labelled by $`\overline{\nu })`$ in $`\widehat{𝒮}`$ \| \| \| Perturb $`\rho _{\overline{\nu }}=\rho _{\overline{\nu }}+ϵ`$ \| \| \| Specify metric data on perturbed point \| \| \| Evaluate $`\widehat{F}[\rho _{\overline{\nu }}+ϵ]`$ at the $`\overline{\mu }`$-th point. \| \| \| Compute the $`\overline{\mu }\overline{\nu }`$ component of the Jacobian matrix \| \| \| using (11) \| \| End loop over points on $`\widehat{𝒮}`$. \| This gives the Jacobian matrix, $`\widehat{J}_{\overline{\mu }\overline{\nu }}`$, for $`\widehat{F}`$ evaluated. $`\widehat{J}_{\overline{\mu }\overline{\nu }}`$ is a $`(N_s^22N_s+2)\times (N_s^22N_s+2)`$ matrix which is used in Newton’s method as follows: \| Start with an initial guess surface $`\widehat{\rho }=\widehat{\rho }_0`$ \| \| --- \| \| while $`\widehat{F}>`$ stopping criterion \| \| \| Compute the Jacobian $`\widehat{J}_{\overline{\mu }\overline{\nu }}`$ for the current $`\widehat{\rho }`$ \| \| \| Evaluate $`\widehat{F}[\widehat{\rho }]`$ \| \| \| Solve $`\widehat{J}\delta \widehat{\rho }=\widehat{F}[\widehat{\rho }]`$ for $`\delta \widehat{\rho }`$ \| \| \| Update the surface $`\widehat{\rho }=\widehat{\rho }+\delta \widehat{\rho }`$ \|
warning/0002/hep-th0002026.html	ar5iv	text	# 1 Introduction ## 1 Introduction The so called triplectic quantization is a general Lagrangian gauge theory quantization procedure following the general lines of the field antifield or Batalin Vilkovisky (BV) method but with the requirement of extended BRST (BRST plus anti-BRST) invariance rather than just BRST. In the usual BV quantization the BRST invariance is translated into the so called master equation. At zero loop order this equation is well defined and its solution, together with the appropriate requirements corresponding to gauge fixing, leads to the construction of the complete structure of ghosts, antighosts, ghosts for ghosts, etc. At higher orders in $`\mathrm{}`$ one needs however to introduce some regularization procedure in order to give a well defined meaning to the mathematical objects involved in the formal master equation. Anomalies and Wess Zumino terms can this way be calculated at one loop order. In the triplectic quantization the extended BRST invariance is translated into a set of two master equations corresponding to the requirements of BRST and anti-BRST invariances respectively. As in the standard BV case, both equations have formally an expansion in loop order. One then expects that anomalies and Wess Zumino terms should show up at one loop order as long as one is able to introduce appropriate regularization schemes. These features are not present in the recently discussed case of Yang Mills theory. In that case only the zero loop order corrections are relevant, as there are no anomalies. The important features of calculation of anomalies and counterterms in the triplectic context have not yet been discussed in the literature. In this article we will discuss the W2 model where the one loop order corrections will nicely illustrate the behavior of the quantum master equations, compared with the standard BV case. We will also show how to fix the gauge by means of canonical transformations. ## 2 Triplectic quantization Considering some gauge theory, we enlarge the original field content $`\varphi ^i`$, adding all the usual gauge fixing structure: ghosts, antighosts and auxiliary fields associated with the original gauge symmetries. The resulting set will be denoted as $`\varphi ^A`$. Then we associate with each of these fields five new quantities, introducing the sets: $`\overline{\varphi }^A`$, $`\varphi _A^{\mathrm{\hspace{0.17em}1}}`$ , $`\varphi _A^{\mathrm{\hspace{0.17em}2}}`$,$`\pi _A^1`$ and $`\pi _A^2`$. The Grassmanian parities of these fields are: $`ϵ(\varphi ^A)=ϵ(\overline{\varphi }^A)ϵ_A`$, $`ϵ(\varphi _A^a)=ϵ(\pi _A^a)=ϵ_A+\mathrm{\hspace{0.17em}1}`$. In this way the ideas of extended BRST quantization in the antifield context previously discussed in are put in a completely anticanonical setting. The extended BRST invariance of the generating functional, defined on this 6n dimensional space, is equivalent to the fact that the quantum action $`W`$ is a solution of the two master equations: $$\frac{1}{2}\{W,W\}^a+V^aW=i\mathrm{}\mathrm{\Delta }^aW$$ (1) where the indices $`a=\mathrm{\hspace{0.17em}1},\mathrm{\hspace{0.17em}2}`$ correspond respectively to BRST and anti-BRST invariances and the extended form of the antibrackets, triangle and $`V`$ operators read $$\{F,G\}^a\frac{^rF}{\varphi ^A}\frac{^lG}{\varphi _A^a}+\frac{^rF}{\overline{\varphi }^A}\frac{^lG}{\pi _A^a}\frac{^rF}{\varphi _A^a}\frac{^lG}{\varphi ^A}\frac{^rF}{\pi _A^a}\frac{^lG}{\overline{\varphi }^A}$$ (2) $$\mathrm{\Delta }^a(1)^{ϵ_A}\frac{^l}{\varphi ^A}\frac{^l}{\varphi _A^a}+(1)^{ϵ_A}\frac{^l}{\overline{\varphi }^A}\frac{^l}{\pi _A^a}$$ (3) $$V^a=\frac{1}{2}ϵ^{ab}\left(\varphi _{Ab}^{}\frac{^r}{\overline{\varphi }^A}(1)^{ϵ_A}\pi _{Ab}\frac{^r}{\varphi ^A}\right).$$ (4) here and in the rest of the article, unless explicitly indicated, we are adopting the convention of summing over repeated indices. The Vacuum functional is normally defined including also an extra functional $`X`$ $$Z=[𝒟\varphi ][𝒟\varphi ^{}][𝒟\pi ][𝒟\overline{\varphi }][𝒟\lambda ]exp\{\frac{i}{\mathrm{}}\left(W+X\right)\}$$ (5) that represents gauge fixing and must satisfy the equations $$\frac{1}{2}\{X,X\}^aV^aX=i\mathrm{}\mathrm{\Delta }^aX$$ (6) An alternative way of gauge fixing, using canonical transformations rather than including the functional $`X`$ was proposed in . We will use this method in section 4 for gauge fixing W2 theory. Expanding the quantum action in powers of $`\mathrm{}`$: $`W=S+\mathrm{}M_1+\mathrm{}`$ we can look at the two first orders of the master equations $`{\displaystyle \frac{1}{2}}\{S,S\}^a`$ $`+`$ $`V^aS=\mathrm{\hspace{0.17em}0}`$ $`\{S,M_1\}^a`$ $`+`$ $`V^aM_1=i\mathrm{\Delta }^aS`$ (7) For a gauge theory with closed and irreducible algebra, corresponding to a classical action $`S_0[\varphi ^i]`$, a solution for the zero loop order action $`S`$ is: $$S=S_0+\varphi _A^a\delta _a\varphi ^A+\frac{1}{2}\overline{\varphi }_A\delta _2\delta _1\varphi ^A+\frac{1}{2}ϵ^{ab}\varphi _{Aa}^{}\pi _b^A$$ (8) where the $`\delta _a`$ represent gauge fixed BRST ($`a=1`$) and anti-BRST ($`a=2`$) transformations of the fields (in other words, for theories with closed algebra, the standard BRST extended algebra associated with the gauge theory). In this article we will not be dealing with the generalized BRST transformations of the triplectic formalism but just with standard transformations that do not involve the antifields. Let us consider now the one loop order equation. As it happens in the standard BV case, we need to introduce a regularization procedure in order to give a well defined meaning to the operator $`\mathrm{\Delta }^aS`$ as there are two functional derivatives acting on the same space time point. If we consider actions of the form (8) we see that the second term in the $`\mathrm{\Delta }^a`$ operator will not contribute and the important term in the action is just $`\varphi _A^a\delta _a\varphi ^A`$. That means we must regularize: $$\frac{^l}{\varphi ^A}\frac{^l}{\varphi _A^a}\left(\varphi _A^a\delta _a\varphi ^A\right).$$ (9) If the BRST algebra is such that the BRST and anti-BRST transformations are symmetrical, just changing ghosts by antighosts, then the same regularization can be used in both sectors. Moreover, we can use the same regularization used in the standard BV quantization. ## 3 Extended BRST invariance in W2 gravity The classical action corresponding to W2 gravity reads $$S_0=\frac{1}{2\pi }d^2x\left(\varphi \overline{}\varphi h(\varphi )^2\right)$$ (10) and the corresponding BRST anti BRST algebra, satisfying $`(\delta _1)^2=(\delta _2)^2=\delta _1\delta _2+\delta _2\delta _1=\mathrm{\hspace{0.17em}0}`$ is $`\delta _1\varphi `$ $`=`$ $`c__1\varphi `$ $`\delta _1h`$ $`=`$ $`\overline{}c__1hc__1+hc__1`$ $`\delta _1c__1`$ $`=`$ $`c__1c__1`$ $`\delta _1c__2`$ $`=`$ $`b`$ $`\delta _2\varphi `$ $`=`$ $`c__2\varphi `$ $`\delta _2h`$ $`=`$ $`\overline{}c__2hc__2+hc__2`$ $`\delta _2c__1`$ $`=`$ $`bc__2c__1c__1c__2`$ $`\delta _2c__2`$ $`=`$ $`c__2c__2`$ where we are representing BRST and anti-BRST transformations respectively as $`\delta _1`$ and $`\delta _2`$. The general form of the gauge fixed action, after functionally integrating over the auxiliary fields of the triplectic formalism: $`\overline{\varphi }^A`$, $`\varphi _A^a`$ and $`\pi _A^a`$ is $$S=S_0+\delta _1\delta _2B$$ (12) where $`B[\varphi ^A]`$ is a bosonic functional. Therefore the ultimate result of triplectic quantization would be to build up such an object. However it is not possible to find a bosonic functional $`B`$ that removes the degeneracy of the action $`S_0`$ using just the fields of algebra (3). We need more fields in order to obtain such a gauge fixing in W2 theory. Inspired in the extended algebra for the bosonic string from ref. we can introduce the bosonic fields $`L`$ and $`\lambda `$ and the fermionic fields $`\eta `$ and $`\overline{\eta }`$ and try transformations of the form $`\delta _1\eta `$ $`=`$ $`\eta c__1+\mathrm{\hspace{0.17em}2}\alpha c__1\eta `$ $`\delta _1L`$ $`=`$ $`a_1\eta +Lc__1\mathrm{\hspace{0.17em}2}\alpha c__1L`$ $`\delta _1\lambda `$ $`=`$ $`\lambda c__1\mathrm{\hspace{0.17em}2}\alpha c__1\lambda `$ $`\delta _1\overline{\eta }`$ $`=`$ $`a_2\lambda +\overline{\eta }c__1+\mathrm{\hspace{0.17em}2}c__1\overline{\eta }`$ $`\delta _2\eta `$ $`=`$ $`b_1\lambda +\eta c__2+\mathrm{\hspace{0.17em}2}\alpha c__2\eta `$ $`\delta _2L`$ $`=`$ $`b_2\overline{\eta }+Lc__2\mathrm{\hspace{0.17em}2}\alpha c__2L`$ $`\delta _2\lambda `$ $`=`$ $`\lambda c__2\mathrm{\hspace{0.17em}2}\alpha c__2\lambda `$ $`\delta _2\overline{\eta }`$ $`=`$ $`\overline{\eta }c__2+\mathrm{\hspace{0.17em}2}\alpha c__2\overline{\eta }`$ (13) Extended nilpotency is satisfied if $`a_1b_1+a_2b_2=\mathrm{\hspace{0.17em}0}`$ for any $`\alpha `$. We will choose $`a_1=a_2=b_1=b_2=1`$ and $`\alpha =\mathrm{\hspace{0.17em}1}`$. In this enlarged space we can choose the gauge fixing boson as $$B=L(h\stackrel{~}{h})$$ (15) where $`\stackrel{~}{h}`$ is a BRST anti-BRST invariant background field. If we redefine the fields as $`\lambda ^{}`$ $`=`$ $`\lambda +\overline{\eta }c__1+\mathrm{\hspace{0.17em}2}c__1\overline{\eta }+Lb\eta c__2^2Lc__1c__2Lc__1c__2`$ $`+`$ $`^2c__1Lc__2+\mathrm{\hspace{0.17em}2}c__1Lc__22c_2\eta \mathrm{\hspace{0.17em}2}c__2Lc__1+\mathrm{\hspace{0.17em}4}c__2c__1L\mathrm{\hspace{0.17em}2}bL`$ $`\eta ^{}`$ $`=`$ $`\eta +Lc__1\mathrm{\hspace{0.17em}2}c__1L`$ $`b^{}`$ $`=`$ $`\overline{\eta }+Lc__2`$ (16) the gauge fixing action gets $$\delta _1\delta _2\left(L(h\stackrel{~}{h})\right)=\lambda ^{}(h\stackrel{~}{h})+b^{}(\overline{}c__1hc__1+hc__1)+\eta ^{}(\overline{}h+h)c__2+L(\overline{}h+h)b.$$ (17) The two last terms cancel each other by a supersymmetric compensation in the path integral while the remaining two first terms correspond to the gauge fixing obtained in ref. in the standard BV scheme with just BRST invariance. Thus, the boson $`B`$ of eq. (15) would appropriately fix the gauge of W2 gravity. Now we will see in the next section how to arrive at this gauge fixing action of eq. (17) starting from the triplectic action. ## 4 Gauge fixing by canonical transformations One interesting way to get the gauge fixed action of the form $`S=S_0+\delta _1\delta _2B`$ from the triplectic action (8) is to perform canonical transformations in the triplectic space. These transformations have been studied in . For each of the antibrackets of eq. (2) with $`a=1,2`$ we introduce a generator $`F_a[\varphi ^A,\overline{\varphi }^A,\varphi _A^a,\pi _A^a]`$ and write out the set of transformations $`\varphi ^A`$ $`=`$ $`{\displaystyle \frac{F_a}{\varphi _A^a}}`$ $`\varphi _A^a`$ $`=`$ $`{\displaystyle \frac{F_a}{\varphi ^A}}`$ $`\overline{\varphi }^A`$ $`=`$ $`{\displaystyle \frac{F_a}{\pi _A^a}}`$ $`\pi _A^a`$ $`=`$ $`{\displaystyle \frac{F_a}{\overline{\varphi }^A}},`$ (18) where there is no sum over $`a`$. If the matrix $$T_a^{\alpha \beta }=\frac{^r^rF_a}{z_\alpha ^az_\beta }$$ (19) (where there is again no sum over $`a`$ and we are defining $`\{z^\alpha \}\{\varphi ^A,\overline{\varphi }^A\}`$ and $`\{z_\alpha ^a\}\{\varphi _A^a,\pi _A^a\}`$) is invertible, each of these transformations, for fixed $`a=1`$ or $`2`$ will not change the form of the corresponding antibracket. If additionally the generators $`F_1`$ , $`F_2`$ both with non singular matrices (19) satisfy also the constraints $`{\displaystyle \frac{F_1}{\varphi _A^{\mathrm{\hspace{0.17em}1}}}}`$ $`=`$ $`{\displaystyle \frac{F_2}{\varphi _A^{\mathrm{\hspace{0.17em}2}}}}`$ $`{\displaystyle \frac{F_1}{\pi _A^1}}`$ $`=`$ $`{\displaystyle \frac{F_2}{\pi _A^2}}.`$ (20) Then the complete set of transformations (4) including both $`a=1`$ and $`a=2`$ will leave the two antibrackets invariant, preserving the complete triplectic anticanonical structure. The constraints (4) restrict the possible dependence of the generators of these transformations on the variables $`\varphi _A^a`$ and $`\pi _A^a`$. Their general form can be written as $$F_a=\mathbf{\hspace{0.17em}1}_a+f_a$$ (21) with $`\mathrm{𝟏}_a`$ $`=`$ $`\varphi ^A\varphi _{Aa}^{}+\overline{\varphi }^A\pi _{Aa}^{}`$ $`f_1`$ $`=`$ $`g_1[\varphi ,\overline{\varphi }]+g_3^A[\varphi ,\overline{\varphi }]\pi _A^1+g_4^A[\varphi ,\overline{\varphi }]\varphi _A^{\mathrm{\hspace{0.17em}1}}`$ $`f_2`$ $`=`$ $`g_2[\varphi ,\overline{\varphi }]+g_3^A[\varphi ,\overline{\varphi }]\pi _A^2+g_4^A[\varphi ,\overline{\varphi }]\varphi _A^{\mathrm{\hspace{0.17em}2}},`$ (22) where we have explicitly separated an identity operator $`\mathrm{𝟏}_a`$ just for future convenience. Now going back to the equation (4) we see that general triplectic canonical transformations can be put in the form $`\varphi ^A`$ $`=`$ $`\varphi ^A+g^[\varphi ,\overline{\varphi }]`$ $`\varphi _A^a`$ $`=`$ $`\varphi _A^a+{\displaystyle \frac{^rg_a}{\varphi ^A}}[\varphi ,\overline{\varphi }]+{\displaystyle \frac{^rg_3^B}{\varphi ^A}}[\varphi ,\overline{\varphi }]\pi _B^a+{\displaystyle \frac{^rg_4^B}{\varphi ^A}}[\varphi ,\overline{\varphi }]\varphi _B^a`$ $`\overline{\varphi }^A`$ $`=`$ $`\overline{\varphi }^A+g_3^A[\varphi ,\overline{\varphi }]`$ $`\pi _A^a`$ $`=`$ $`\pi _A^a+{\displaystyle \frac{^rg_a}{\overline{\varphi }^A}}[\varphi ,\overline{\varphi }]+{\displaystyle \frac{^rg_3^B}{\overline{\varphi }^A}}[\varphi ,\overline{\varphi }]\pi _B^a+{\displaystyle \frac{^rg_4^B}{\overline{\varphi }^A}}[\varphi ,\overline{\varphi }]\varphi _B^a`$ The condition that a canonical transformation reproduces the gauge fixing corresponding to some boson B, after we express the result in terms of the transformed fields and impose the condition that $`\overline{\varphi }^A`$, $`\varphi _A^a`$ and $`\pi _A^a`$ are set to zero reads $$\frac{f_a^{}}{\varphi ^A}\delta _a\varphi ^A+\frac{1}{2}g_3^{}\delta _2\delta _1\varphi ^A\frac{1}{2}ϵ^{ab}\frac{f_a^{}}{\varphi ^A}\frac{f_b^{}}{\overline{\varphi }^A}=\delta _2\delta _1B[\varphi ^A],$$ (24) where we are defining the primed functions as the corresponding function, written in terms of $`\varphi ^A`$ and $`\overline{\varphi }^A`$, taken at $`\overline{\varphi }^A=\mathrm{\hspace{0.17em}0}`$ $$f_a^{}[\varphi ]=f_a[\varphi ,\overline{\varphi }(\varphi ,\overline{\varphi }^{})]\|_{_{\overline{\varphi }^{}=\mathrm{\hspace{0.17em}0}}}$$ (25) and a similar definition for $`g_i^{}`$. Considering our W2 case, two illustrative possibilities are to choose $`g_1`$ $`=`$ $`g_3=g_4=\mathrm{\hspace{0.17em}0}`$ $`g_2`$ $`=`$ $`\delta _1\left(L(h\stackrel{~}{h})\right)=L(\overline{}c__1hc__1+hc__1)+(\eta +Lc__1\mathrm{\hspace{0.17em}2}lc__1)(h\stackrel{~}{h}),`$ (26) or $`g_1`$ $`=`$ $`\delta _2\left(L(h\stackrel{~}{h})\right)=L(\overline{}c__2hc__2+hc__2)+(\overline{\eta }+Lc__2\mathrm{\hspace{0.17em}2}lc__2)(h\stackrel{~}{h})`$ $`g_2`$ $`=`$ $`g_3=g_4=\mathrm{\hspace{0.17em}0}`$ (27) In both cases we get the gauge fixing action (17) if we perform the corresponding transformation in the fields of action $`S`$ and then set all the primed antifields to zero. ## 5 One loop order The first point that must be investigated is the possible effect of the introduction of the fields $`L`$ , $`\lambda `$ , $`\eta `$ and $`\overline{\eta }`$ in the anomalies of the model. In other words, we must see if the cohomology of our extended formulation is the same as that from the original one. We can introduce a filtration $`𝒩`$ that counts the number of fields and expand the BRST anti-BRST operators according to this filtration $`\delta _1=\delta _1^{(0)}+\delta _1^{(1)}`$; $`\delta _2=\delta _2^{(0)}+\delta _2^{(1)}`$. The first order piece of the algebra reads $`\delta _1^{(0)}\eta `$ $`=`$ $`0`$ $`\delta _1^{(0)}L`$ $`=`$ $`\eta `$ $`\delta _1^{(0)}\lambda `$ $`=`$ $`0`$ $`\delta _1^{(0)}\overline{\eta }`$ $`=`$ $`\lambda `$ $`\delta _2^{(0)}\eta `$ $`=`$ $`\lambda `$ $`\delta _2^{(0)}L`$ $`=`$ $`\overline{\eta }`$ $`\delta _2^{(0)}\lambda `$ $`=`$ $`0`$ $`\delta _2^{(0)}\overline{\eta }`$ $`=`$ $`0.`$ (28) Looking at this algebra we realize that this fields form doublets with respect to both the BRST and anti-BRST transformations. By a doublet one means a pair of fields say $`u,v`$ whose transformations are of the form $`\delta u=v,\delta v=\mathrm{\hspace{0.17em}0}`$. Fields that show up just in doublets are absent from the cohomology of the BRST operator (or correspondingly homology of the anti BRST operator). Moreover, the cohomology ($`a=1`$) or homology ($`a=2`$) of the operator $`\delta _a`$ is contained in the cohomology (homology) of the corresponding $`\delta _a^{(0)}`$ , $`a=\mathrm{\hspace{0.17em}1}`$. Thus we conclude that the inclusion of the fields $`L,\lambda ,\eta `$ and $`\overline{\eta }`$ does not change the cohomology (homology)of the W2 theory. We can then consider the same quantum correction $`\mathrm{\Delta }S`$ to the first order master equations as for the standard formulation of W2. The calculation of $`\mathrm{\Delta }S`$ depends on the introduction of a regularization procedure and the result depends on the result. But all the possible results differ just by trivial terms (in the cohomological sense). The simplest way to right the results of adapting them to the extended symmetry is: $`(\mathrm{\Delta }^1S)_{Reg}`$ $`=`$ $`{\displaystyle \frac{1}{12\pi }}{\displaystyle d^2x\left(c__1^3h\right)}`$ $`(\mathrm{\Delta }^2S)_{Reg}`$ $`=`$ $`{\displaystyle \frac{1}{12\pi }}{\displaystyle d^2x\left(c__2^3h\right)}`$ (30) The presence of this term in the master equation at one loop order means a breaking in the BRST invariance. What one normally does in the BV quantization is then to introduce an extra Wess Zumino field $`\theta `$ representing the extra degree of freedom corresponding to the anomalous breaking of gauge invariance. The BRST extended version of this procedure would correspond to define this new field with the transformations: $`\delta _1\theta `$ $`=`$ $`c__1+c__1\theta `$ $`\delta _2\theta `$ $`=`$ $`c__2+c__2\theta .`$ (31) Then we verify that the counterterm $$M_1=\frac{1}{24\pi }d^2x\left(\theta \overline{}\theta h\theta \theta +h^2\theta \right)$$ (32) solves the master equations: $$\{S,M_1\}^a+V^aM_1=i(\mathrm{\Delta }^aS)_{Reg}$$ (33) for both $`a=\mathrm{\hspace{0.17em}1},2`$. As a remark we mention that a different approach could be taken to the addition of the field $`\theta `$ to the theory. We mean, one could include another ghost associated with the invariance of the classical action with respect to any transformations in $`\theta `$. As a result the master equation would never be solved. One would only be able to shift the anomaly to this new symmetry, leaving the original gauge symmetry unbroken, but not the BRST (and anti-BRST symmetries). We will keep here the point of view of that the field $`\theta `$ represents the new degree of freedom that shows up at quantum level as a consequence of the anomalous breaking of the original gauge symmetry. Thus we do not add any extra ghost. Then we see that at one loop order the triplectic quantization reproduces the so called Wess Zumino mechanism of restoring the gauge invariance of an anomalous gauge theory by means of the introduction of an extra degree of freedom. ## 6 Conclusion We have seen that, enlarging the space of fields, it is possible to formulate W2 gravity with extended BRST invariance. We have proven that these enlargement of the representation do not change the cohomology of the theory. We have calculated the one loop order corrections in the triplectic quantization for the model. We have seen that the anomalies and counterterms of the BRST and anti BRST sectors are essentially the same, up to changing ghosts by antighosts. The question that can then be raised then is: how general is this result? In the present case this happens because the BRST and anti BRST algebras are symmetric and thus one trivially concludes that the cohomology (homology) is the same for both symmetries (again, up to changing ghosts by antighosts). It seems an interesting future task to look, in some gauge theory, for a $`\delta _2`$ symmetry satisfying the extended algebra $`(\delta _2)^2=\delta _1\delta _2+\delta _2\delta _1=\mathrm{\hspace{0.17em}0}`$ but with a different homology. Acknowledgements: The authors are partially supported by CNPq., FINEP and FUJB (Brazilian Research Agencies).
warning/0002/cond-mat0002449.html	ar5iv	text	# Reactive Hall response \[ ## Abstract The zero temperature Hall constant $`R_H`$, described by reactive (nondissipative) conductivities, is analyzed within linear response theory. It is found that in a certain limit, $`R_H`$ is directly related to the density dependence of the Drude weight implying a simple picture for the change of sign of charge carriers in the vicinity of a Mott-Hubbard transition. This novel formulation is applied to the calculation of $`R_H`$ in quasi-one dimensional and ladder prototype interacting electron systems. \] It is now well known that in strongly correlated systems the, zero temperature (T=0), reactive part of the conductivity can be used as a criterion of a metallic or insulating ground state. In particular, following the work of Kohn, the imaginary part of the conductivity, $`\sigma ^{\prime \prime }(\omega 0)=2D/\omega `$, characterized by $`D`$ (now called the “Drude weight” or charge stifness), can be related to the ground state energy density $`ϵ^0`$ dependence on an applied fictitious flux $`\varphi `$ as $`D=(1/2)^2ϵ^0/\varphi ^2\|_{\varphi 0}`$. A similar question is posed by the doping of an insulating state, where it would be interesting to have a simple description of the charge carriers sign as probed in a Hall experiment. For instance, we would like to describe the doping of a Mott-Hubbard insulator; within a semiclassical approach it is expected that the Hall constant $`R_H+1/e\delta `$, hole-like (positive) near half-filling ($`\delta =1n`$, $`n`$=density), changing to $`R_H1/en`$, electron-like at low densities, the turning point depending on the interaction. Over the recent years, ingenious ways have been proposed for characterizing this sign change and strongly correlated electron systems, as the $`tJ`$ model have been studied. In particular, following the suggestion to focus at the T=0 Hall constant within linear response theory, the $`R_H`$ of a hole in the $`tJ`$ model was analyzed and a numerical method was proposed for calculating the Hall response in ladder systems. This activity is partly motivated by the physics of high temperature superconductors viewed as doped Mott-Hubbard insulators and related Hall measurements showing a change of the sign of carriers with doping. In this work, we will show that within a certain frequency-$`\omega `$, wavevector-$`q`$ limiting procedure, the T=0, $`\omega 0`$, thus “reactive” Hall constant, is simply related to the density dependence of the Drude weight. Following this point of view, we recover in a straightforward way: (i) the semiclassical expressions for $`R_H`$ at low density and near an insulating state, (ii) a physical picture of the sign change of carriers in the vicinity of a Mott-Hubbard transition and its dependence on interaction strength, (iii) a common expression used to describe the Hall constant in quasi-one dimensional conductors described by a band picture, (iv) good accord with $`R_H`$ for ladder systems calculated using the numerical method proposed in . The Hamiltonian In the following we will consider a generic Hamiltonian for fermions on a lattice, where for simplicity we describe the kinetic energy term by a one band tight binding model; it is straightforward to extend this formulation to a many-band or continuum system. The sites are labeled $`l(m)`$ along the $`x(y)`$-direction with periodic boundary conditions in both directions: $`H`$ $`=`$ $`(t){\displaystyle \underset{l,m}{}}e^{i\varphi ^x(t)}e^{iA_m}c_{l+1,m}^{}c_{l,m}+h.c.`$ (1) $`+`$ $`(t^{}){\displaystyle \underset{l,m}{}}e^{i\varphi _{m+1/2}^y(t)}c_{l,m+1}^{}c_{l,m}+h.c.`$ (2) $`+`$ $`\widehat{U},l=1,\mathrm{},L_x;m=1,\mathrm{},L_y.`$ (3) $`c_{l,m}(c_{l,m}^{})`$ is an annihilation (creation) operator at site $`(l,m)`$ and the spin is neglected as it enters in a trivial way in the formulation. The $`\widehat{U}`$ term can represent a many-particle interaction or a one particle potential. We take the lattice constant so as to consider a unit volume, electric charge $`e=1`$ and $`\mathrm{}=1`$. We add a magnetic field along the $`z`$-direction, modulated by a one component wavevector-$`q`$ along the $`y`$-direction, generated by the vector potential $`A_m`$; this allows to take the zero magnetic field limit smoothly: $`A_m`$ $`=`$ $`e^{iqm}{\displaystyle \frac{iB}{2\mathrm{sin}(q/2)}}e^{iqm}{\displaystyle \frac{iB}{q}}`$ (4) $`B_{m+1/2}`$ $`=`$ $`(A_{m+1}A_m)=Be^{iq(m+1/2)}`$ (5) (for convenience, we will present the long wavelength limit, substituting $`2\mathrm{sin}(q/2)q`$). Electric fields along the $`x,y`$ directions are generated by time dependent vector potentials: $`\varphi ^{x,y}(t)`$ $`=`$ $`{\displaystyle \frac{E^{x,y}(t)}{iz}},\varphi _{m+1/2}^y(t)=e^{iq(m+1/2)}\varphi ^y(t);`$ (6) $`E^x(t)`$ $`=`$ $`E^xe^{izt},E^y(t)=iE^ye^{izt};z=\omega +i\eta .`$ (7) Currents are defined through derivatives of the Hamiltonian expanded to second order in $`\varphi ^{x,y}`$: $`J^x={\displaystyle \frac{H}{\varphi ^x}},J_q^y={\displaystyle \frac{H}{\varphi ^y}},`$ (8) with the paramagnetic parts: $`j^x`$ $`=`$ $`t{\displaystyle \underset{l,m}{}}(ie^{iA_m}c_{l+1,m}^{}c_{l,m}+h.c.)`$ (9) $`j_q^y`$ $`=`$ $`t^{}{\displaystyle \underset{l,m}{}}e^{iq(m+1/2)}(ic_{l,m+1}^{}c_{l,m}+h.c.).`$ (10) The reactive Hall response From standard linear response theory we obtain: $`J^x`$ $`=`$ $`\sigma _{j^xj^x}E^x(t)+\sigma _{j^xj_q^y}E^y(t)`$ (11) $`J_q^y`$ $`=`$ $`\sigma _{j_q^yj^x}E^x(t)+\sigma _{j_q^yj_q^y}E^y(t).`$ (12) $`\mathrm{}`$ are ground state expectation values in the presence of the magnetic field, with the conductivities $`\sigma _{j^\alpha j^\beta }`$ $`=`$ $`{\displaystyle \frac{i}{z}}({\displaystyle \frac{^2H}{\varphi ^\alpha \varphi ^\beta }}\chi _{j^\alpha j^\beta }),`$ (13) $`\chi _{AB}`$ $`=`$ $`i{\displaystyle _0^{\mathrm{}}}𝑑te^{izt}[A(t),B].`$ (14) Now, in contrast to the usual derivation of the Hall constant expression, we will keep the $`q`$dependence explicit by converting the current-current to current-density correlations using the continuity equation: $`J^x`$ $`=`$ $`\sigma _{j^xj^x}E^x(t)+{\displaystyle \frac{1}{q}}\chi _{j^xn_q}E^y(t)`$ (15) $`J_q^y`$ $`=`$ $`{\displaystyle \frac{1}{q}}\chi _{n_qj^x}E^x(t)+({\displaystyle \frac{z}{q}})^2\chi _{n_qn_q}{\displaystyle \frac{i}{z}}E^y(t),`$ (16) with $`n_q=_{l,m}(ie^{iqm})c_{l,m}^{}c_{l,m}`$. At T=0, the response is non-dissipative so we will study the reactive (out-of phase) induced currents. Furthermore, at this point we will consider the “screening” (or slow) response in the $`y`$direction, by taking the $`(q,\omega )`$ limits in the order $`\omega 0`$ first and $`q0`$ last; in the usual “transport” (or fast) response the limits are in the opposite order. As we will discuss below, this approach leads to a simple physical picture for the Hall constant and it might be argued that at least for certain cases, for example for a system of finite size in the $`y`$direction, it is indeed the right one. The expressions (16) for the currents become: $`J^x_0`$ $`=`$ $`\sigma _{j^xj^x}^{\prime \prime }(\omega 0))(iE^x(t))`$ (17) $`+`$ $`{\displaystyle \frac{1}{q}}\chi _{j^xn_q}^{}(\omega =0)E^y(t)`$ (18) $`J_q^y_0`$ $`=`$ $`{\displaystyle \frac{1}{q}}\chi _{n_qj^x}^{}(\omega =0)E^x(t)`$ (19) $`+`$ $`({\displaystyle \frac{\omega }{q}})^2{\displaystyle \frac{1}{\omega }}\chi _{n_qn_q}^{}(\omega =0)(iE^y(t)),`$ (20) where the subscript zero denotes the leading order in $`\omega `$ response, $`\chi _{AB}^{}(\omega =0)={\displaystyle \underset{n>0}{}}{\displaystyle \frac{0\|A\|nn\|B\|0+h.c.}{E_nE_0}},`$ (21) and $`\|n(E_n)`$ are eigenstates (eigenvalues) of the Hamiltonian in the presence of the magnetic field. Now, following Kohn’s observation, we can identify the different terms as derivatives of the ground state energy density $`ϵ^0`$ of a fictitious Hamiltonian depending on static $`\varphi ^x,\mu _q`$ fields: $`H`$ $`=`$ $`(t){\displaystyle \underset{l,m}{}}(e^{i\varphi ^x}e^{iA_m}c_{l+1,m}^{}c_{l,m}+h.c.)`$ (22) $`+`$ $`(t^{}){\displaystyle \underset{l,m}{}}(c_{l,m+1}^{}c_{l,m}+h.c.)+\mu _qn_q+\widehat{U}.`$ (23) For $`H(\lambda ,\mu )`$, using the following identity, $`ϵ_{\mu \lambda }^0={\displaystyle \frac{^2ϵ^0}{\mu \lambda }}`$ $`=`$ $`0\|{\displaystyle \frac{^2H}{\mu \lambda }}\|0`$ (24) $``$ $`{\displaystyle \underset{m>0}{}}{\displaystyle \frac{0\|\frac{H}{\mu }\|mm\|\frac{H}{\lambda }\|0+h.c.}{E_mE_0}},`$ (25) we can rewrite the currents as: $`J^x_0`$ $`=`$ $`{\displaystyle \frac{ϵ_{\varphi ^x\varphi ^x}^0}{\omega }}(iE^x(t))+({\displaystyle \frac{1}{q}})ϵ_{\varphi ^x\mu _q}^0E^y(t)`$ (26) $`J_q^y_0`$ $`=`$ $`{\displaystyle \frac{1}{q}}ϵ_{\mu _q\varphi ^x}^0E^x(t){\displaystyle \frac{\omega }{q^2}}ϵ_{\mu _q\mu _q}^0(iE^y(t)).`$ (27) Finally, setting $`J_q^y_0=0`$ we determine the “reactive” Hall constant: $`R_H{\displaystyle \frac{1}{B}}{\displaystyle \frac{E^y}{J^x_0}}=({\displaystyle \frac{q}{B}}){\displaystyle \frac{ϵ_{\mu _q\varphi ^x}^0}{ϵ_{\varphi ^x\varphi ^x}^0ϵ_{\mu _q\mu _q}^0+ϵ_{\mu _q\varphi ^x}^0ϵ_{\varphi ^x\mu _q}^0}}.`$ (28) Neglecting the cross-terms $`ϵ_{\mu _q\varphi ^x}^0ϵ_{\varphi ^x\mu _q}^0`$ and Taylor expanding the numerator in $`B`$, we can rewrite $`R_H`$ as: $`R_H=q{\displaystyle \frac{\frac{^3ϵ^0}{B\mu _q\varphi ^x}}{ϵ_{\varphi ^x\varphi ^x}^0ϵ_{\mu _q\mu _q}^0}}=q{\displaystyle \frac{\frac{}{\mu _q}(\frac{^2ϵ^0}{B\varphi ^x})}{ϵ_{\varphi ^x\varphi ^x}^0ϵ_{\mu _q\mu _q}^0}}.`$ (29) Using (25) we find the final expression: $`R_H={\displaystyle \frac{\frac{D_q}{\mu _q}}{D\kappa _q}}`$ (30) where, $`D_q`$ $`=`$ $`{\displaystyle \frac{1}{2}}[0\|T_q^x\|0`$ (31) $``$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{0\|j^x\|mm\|j_q^x\|0+h.c.}{ϵ_mϵ_0}}],`$ (32) $`j_q^x`$ $`=`$ $`(t){\displaystyle \underset{l,m}{}}(ie^{iqm})(ic_{l+1,m}^{}c_{l,m}+h.c.),`$ (33) $`T_q^x`$ $`=`$ $`(t){\displaystyle \underset{l,m}{}}(ie^{iqm})(c_{l+1,m}^{}c_{l,m}+h.c.).`$ (34) $`D=\frac{1}{2}ϵ_{\varphi ^x\varphi ^x}^0`$, the Drude weight, is identical to $`D_q`$ by the replacement of $`j_q^x`$ ($`T_q^x`$) by $`j^x`$ ($`T^x`$). $`\kappa _q=ϵ_{\mu _q\mu _q}^0=n_q/\mu _q`$ is the compressibility corresponding to the density modulation $`n_q`$. Notice that the spatial dependence of $`j_q^x`$ and $`n_q`$ is the same as that of $`A_m`$. Taking the $`q0`$ limit, we obtain a particularly simple expression for $`R_H`$: $`R_H={\displaystyle \frac{1}{D}}{\displaystyle \frac{D}{n}}.`$ (35) A handwaving argument leading to expression (35) for $`t^{}0`$ is as follows: $`A_m`$ corresponds to a twist of boundary conditions on chain$`m`$, inducing an extra current on each chain proportional to $`D`$ (besides the uniform one induced by the flux $`\varphi ^x`$); minimization of the energy at fixed $`x`$current gives rise to an $`m`$dependent charge density. This induced charge density can then be canceled by the “Hall potential” $`\mu _q`$. Note that a similar idea, analyzing the Hall constant in terms of independent channels (edge states), exists in the literature of the Quantum Hall effect. This expression is appealing as it gives a direct, intuitive understanding for the change of sign of charge carriers in the vicinity of a metal-insulator transition. First, at low densities, $`Dn`$ giving $`R_H1/n`$; close to a Mott insulator $`D\delta =1n`$, implying $`R_H+1/\delta `$. Furthermore, we obtain a change of sign in the vicinity of a Mott transition at a density which depends on the interaction strength and is given by the position of the maximum of $`D`$. Second, for independent electrons, where $`D`$ is proportional to the kinetic energy, by taking the limit $`t^{}0`$ and calculating $`D`$ as a sum of $`D`$’s for individual $`x`$chains, we obtain from (35): $$D=\frac{2t}{\pi }\mathrm{sin}(\frac{\pi n}{2}),R_H=\frac{\pi }{2}\frac{1}{\mathrm{tan}(\frac{\pi n}{2})},$$ (36) an expression used for the Hall constant of quasi one-dimensional compounds. Considering that the $`t^{}0`$ limit might by subtle, it is of particular theoretical and experimental interest whether the Hall constant of quasi-one dimensional correlated systems is indeed given by the expression and thus related to the Drude weight of the individual chains. The same applies for the transverse Hall effect of weakly coupled planes. Examples In this section we present a generic picture for the behavior of the Hall constant for models of strongly correlated fermions showing a Mott-Hubbard metal-insulator transition. This picture emerges, on the one hand, by an exact calculation of $`R_H`$ for ladder systems using the numerical method of ref. and on the other hand, from the expression (35) assuming nearly decoupled chains ($`t^{}0`$) and calculating $`D(n)`$ for each chain analytically using the Bethe ansatz method. It is clear that this analytical approach refers to either ladder (with $`t^{}0`$) or quasi one-dimensional models. Three prototype models will be discussed: the Hubbard model, as the most experimentally relevant, the spinless fermions model (“t-V”) showing both a metallic and an insulating phase depending on interaction strength and the supersymmetric $`tJ`$ model. (i) The Hubbard model is given by the Hamiltonian: $`H`$ $`=(t){\displaystyle \underset{l,m}{}}(c_{l+1,m,\sigma }^{}c_{l,m,\sigma }+h.c.)`$ (37) $`+`$ $`(t^{}){\displaystyle \underset{l,m}{}}(c_{l,m+1,\sigma }^{}c_{l,m,\sigma }+h.c.)+U{\displaystyle \underset{l,m}{}}n_{l,m,}n_{l,m,}.`$ (38) $`c_{l,m,\sigma }(c_{l,m,\sigma }^{})`$ is an annihilation (creation) operator at site $`(l,m)`$ of a fermion with spin $`\sigma =,`$. $`R_H`$ extracted from a Bethe ansatz calculation of $`D(n)`$ for the one dimensional Hubbard model is shown in Fig. 1. This behavior is characteristic of correlated systems undergoing a metal-insulator transition at half-filling: at low densities $`R_H1/n`$, while near half-filling $`R_H+1/\delta `$, the position of change of sign of the carriers depending on the details of the interaction. (ii) The t-V model on a ladder is given by: $`H`$ $`=`$ $`(t){\displaystyle \underset{l,m}{}}(c_{l+1,m}^{}c_{l,m}+h.c.)`$ (39) $`+`$ $`(t^{}){\displaystyle \underset{l}{}}(c_{l,1}^{}c_{l,2}+h.c.)+V{\displaystyle \underset{l,m}{}}n_{l,m}n_{l+1,m}.`$ (40) Here and in the following $`l=1,\mathrm{},L_x,m=1,2`$. For a single chain, this model describes a metallic phase at all densities for $`V<2t`$, while for $`V>2t`$ it is an insulator at half-filling. In Fig. 2 we show $`R_H`$ calculated numerically on finite systems for two values of $`t^{}`$ and analytically from (35) in the $`t^{}0`$ limit. The numerical evaluation being especially sensitive to finite size effects for $`t^{}0`$, we study relatively large values of $`t^{}`$. Results for $`R_H`$ clearly show the difference between the metallic regime $`V=t`$, where at half-filling ($`n=0.5`$) we get $`R_H=0`$, while in the insulating regime $`V=4t`$, we are dealing with $`R_H(n0.5)\mathrm{}`$. (iii) The t-J model on a ladder is given by the Hamiltonian: $`H`$ $`=`$ $`(t){\displaystyle \underset{l,m}{}}(c_{l+1,m,\sigma }^{}c_{l,m,\sigma }+h.c.)`$ (41) $`+`$ $`(t^{}){\displaystyle \underset{l}{}}(c_{l,1,\sigma }^{}c_{l,2,\sigma }+h.c.)`$ (42) $`+`$ $`J{\displaystyle \underset{l,m}{}}(\stackrel{}{S}_{l,m}\stackrel{}{S}_{l+1,m}{\displaystyle \frac{1}{4}}n_{l,m}n_{l+1,m}).`$ (43) $`\stackrel{}{S}_{lm}`$ is the spin operator at site $`(l,m)`$ and the double occupancy on a site is forbidden. In Fig. 3 we show again $`R_H`$ calculated analytically for the “supersymmetric” model, $`J=2t`$, and by numerical evaluation for $`t^{}=0.5t`$ and different size systems. The above three examples show a remarkable agreement between the numerical evaluation of $`R_H`$ on finite size systems using the method of ref. (at finite $`t^{}`$) and the analytical calculation using (35) for $`t^{}0`$, indicating a relative insensitivity on the transverse coupling $`t^{}`$ for ladders. These results confirm the intuitive picture for the behavior of the Hall constant in the vicinity of a metal-insulator transition and present an intriguing link between the Hall constant and the Drude weight. It is possible that $`R_H`$ is dominated at low temperatures by correlations and not the relaxation mechanism so this formulation could have more general validity. In conclusion, the emerging simple physical picture raises the question of the relation of this novel formulation to the traditional semiclassical approach to the Hall constant, its range of validity, the role of relaxation in the description of the Hall effect and of the perspectives for an extension at finite temperatures. Part of this work was done during visits of (P.P.) and (M.L.) at IRRMA as academic guests of EPFL. X.Z. and F.N. acknowledge support by the Swiss National Foundation grant No. 20-49486.96, the EPFL, the Univ. of Fribourg and the Univ. of Neuchâtel.
warning/0002/gr-qc0002077.html	ar5iv	text	# Shell sources as a probe of relativistic effects in neutron star models ## I Introduction and overview For most astrophysical objects Newton’s classical theory of gravity gives a fully satisfactory description. Only when gravitational fields become strong need one consider the possibility that general relativistic effects may play a significant role. A standard index of field strength is $`GM/Rc^2`$, where $`M`$ is an object’s mass, and $`R`$ is its characteristic size. This index is of order unity for black holes and for the universe itself, and much smaller than unity for almost all stars, galaxies, and other astronomical entities. One exception is neutron stars; for compact neutron stars $`GM/Rc^2`$ is on the order of $`0.2`$, and exotic equations of state could lead to even larger values. Despite this, Newtonian gravity is used almost exclusively in studying neutron stars. The obvious reason is the significant increase in difficulty in giving a fully relativistic treatment of neutron star structure, and the enormous increase of difficulty in dealing with fullly relativistic dynamics, i.e. , with oscillations of neutron stars. Newtonian physics has been used even in studies (see for example Refs.) of neutron stars as sources of gravitational waves. Of course, Newtonian gravity per se has no gravitational waves, so the typical procedure is to compute gravitational wave generation as a postprocessing step. More specifically, Newtonian gravity is used to find the fluid motions inside a neutron star associated with some event (core collapse in a supernova, precession of a rotating neutron star, etc.). Those fluid motions are then used as sources in the “quadrupole formula” of general relativity, just as known charge motions would be used in the dipole formula of electromagnetism. If Newtonian theory predicts periodic oscillations, the “adiabatic approximation” can be used to give the damping of the oscillations due to gravitational wave emission: the energy of the oscillations is taken to decrease at the rate at which gravitational waves remove energy. The sufficiency of this approximate procedure caused little worry until relativistic modes of oscillation of neutron stars were discovered that had no counterpart in Newtonian theory. The so-called $`w`$ modes are qualitatively oscillations of the spacetime, like black hole quasinormal (QN) modes, rather than oscillations of the neutron star material like the $`f`$ (fluid) and $`p`$ (pressure) modes of the Newtonian description. Partly due to the existence of these $`w`$ modes, the general question of the sufficiency of Newtonian theory for dealing with gravitational wave processes was emphasized by Andersson and Kokkotas , and stirred considerable interest. To give a clear answer to this question requires a specific and astrophysically plausible event for which gravitational wave generation can be computed in both Newtonian theory and fully relativistically. One would want, for instance, initial data for the fluid and spacetime of a neutron star formed in a supernova core collapse, but the possibililty of giving such initial data is at least several years off. A more tractable model was needed. Three separate groups studied the problem of emission of gravitational waves by a relativistic neutron star due to the close passage of a perturbing particle. This model had the advantage of definitiveness; there was no freedom in choosing (and biasing) the initial spacetime perturbations to be particularly larger or smaller than the fluid perturbations. This model allowed an investigation of whether the excitation of $`w`$ modes was significant, but the model had three serious shortcomings. (i) It was too restrictive. The only parameters were the two constants specifying the particle orbit (say, energy at infinity and impact parameter). There could be significant excitation only if the particle passed close to the neutron star, and this constrained the perturbation to be neither very fast nor very slow. Not only was the timescale limited, but it was coupled to the choice of location of the perturbation. (ii) It was difficult to separate the radiation due to the neutron star (the radiation of interest) from the radiation coming (in some sense) from the orbiting particle. (iii) It was not clear how to compare the fully relativistic computation from the Newtonian computation; at least there was no attempt to do this. We introduce here a very different model for investigating the importance of relativistic effects in neutron stars, and possibly for answering other questions. We consider a perturbative spherical shell around a neutron star at some radius $`R_{\mathrm{shell}}`$. In that shell we dictate the time dependence of a multipole of surface mass-energy density of the shell. There is no equation of state of the shell material constraining our choice. The equations of motion of the shell fix the surface stress once we have specified the surface energy density, so picking that single function of time fully specifies the source. In this manner we can independently choose where the perturbation of the neutron star arises and what its timescale is. We can probe the response of the star to close and far perturbations with slow motions or fast motions. The shell probe has the nice feature that it is straightforward to do the calculation in Newtonian theory, so that a comparison can be made with the relativistic result. A further advantage is that both the Newtonian and relativistic computations allow a separation of the gravitational waves from the shell and those from the star. In the Newtonian computation, this is completely straightforward. The waves from the star are those found from the quadrupole formula applied to the motions only of the stellar fluid. In the relativistic calculation the separation is only approximate, and more care is needed. From the star+shell results it is necessary to subtract the radiation from the shell itself generated in the background spacetime of the star. (The details of this procedure will be given below.) Our main purpose in the present paper is to present the method of using a shell probe and to display its advantages. For that reason we limit the application of this method to the simplest model of a star, the homogeneous incompressible perfect fluid (HIF) model. (See for details about the QN modes of this model.) For this model the only mode associated with motions of the stellar material is a single $`f`$ mode. An important element of our results is that we will present answers not only about excitation of $`w`$ modes, but about the difference in the Newtonian and relativistic predictions of excitation of the $`f`$ mode. In the interest of brevity we limit the analysis to even parity perturbations. The excitation of $`w`$ modes for odd parity should not be remarkably different from the excitation in even parity, and odd parity motions do not couple to fluid motions. The paper is organized as follows: The shell source model is introduced in Sec. II, and the equations governing even parity perturbations due to the shell source are given in Sec. III. Some details of the computational implementation are given in Sec. IV along with a discussion of the method used for subtracting the shell contribution from the relativistic calculation of the gravitational wave due to the shell and star. Numerical results are presented and discussed in Sec. V, and conclusions are given in Sec. VI. Details of the Newtonian calculation are given in the Appendix. Throughout the paper we use geometric units $`G=c=1`$, the metric signature $`(+++)`$ and the conventions of Misner, Thorne and Wheeler. ## II Model: Perturbation of a static star by matter moving on a spherical shell We start with a static and spherically symmetric spacetime background metric $$ds^2=e^{\nu (r)}dt^2+e^{\lambda (r)}dr^2+r^2[d\theta ^2+\mathrm{sin}^2\theta d\phi ^2],$$ (1) describing both the interior and exterior of a star of a barotropic, ideal fluid of mass M and radius $`r=R`$. The stress energy of the fluid is $$T_{\alpha \beta }=(\rho +p)u_\alpha u_\beta +pg_{\alpha \beta },$$ (2) where $`\rho `$ is the mass-energy density and $`p=p(\rho )`$ is the pressure. The mass function $`m(r)`$ is defined by $`e^{\lambda (r)}=12m(r)/r`$, and the structure of the stellar interior is found by solving the hydrostatic equilibrium equations of general relativity. (See, e.g. , Eq. (3) of .) For simplicity, we limit considerations to homogeneous incompressible fluid (HIF) stellar models whose unperturbed interior metric is given by Eqs. (5) and (6) of . The exterior metric is simply the Schwarzschild metric, with $`m(r)`$ equal to a constant $`M`$, and $`\nu (r)=\lambda (r)`$. A spherical thin shell of coordinate radius $`r=R_{\mathrm{shell}}>R`$ surrounds the star. We treat the shell as a perturbation of the spacetime inside and outside the star and we analyze the perturbations only to first order in the parameter of the perturbation. In this order, the shell has spherical geometry described by the 3-metric, $$ds^2\|_{\mathrm{shell}}=\left(1\frac{2M}{R_{\mathrm{shell}}}\right)dt^2+R_{\mathrm{shell}}^2[d\theta ^2+\mathrm{sin}^2\theta d\phi ^2],$$ (3) induced by the Schwarzschild metric. The metric of the perturbed spacetime can be written as $$g_{\alpha \beta }=g_{\alpha \beta }^{(0)}+h_{\alpha \beta }$$ (4) where the “(0)” index denotes the background solution, that of Eq. (1). The Einstein equations to first order in perturbations are $$\delta G_\alpha ^\beta =8\pi [\delta T_\alpha ^\beta {}_{\mathrm{fluid}}{}^{}+\delta T_\alpha ^\beta {}_{\mathrm{shell}}{}^{}].$$ (5) The perturbed stress energy has two contributions. One, denoted by $`\delta T_\alpha ^\beta _{\mathrm{fluid}}`$, is that of the fluid star perturbations and is nonzero only inside the star. The other, denoted by $`\delta T_\alpha ^\beta _{\mathrm{shell}}`$, is the stress energy of the matter in the thin shell and is nonzero only outside the star. Its form, in the coordinates of Eq. (1), is $$\delta T_{\mathrm{shell}}^{\alpha \beta }=\sqrt{12M/r}S^{\alpha \beta }\delta (rR_{\mathrm{shell}}),$$ (6) where $$S^{\alpha \beta }=\underset{ϵ0}{lim}_{R_{\mathrm{shell}}ϵ}^{R_{\mathrm{shell}}+ϵ}T^{\alpha \beta }\frac{dr}{\sqrt{12M/r}}$$ (7) is the surface stress energy of the shell. (See e.g , .) ### A Even parity equations of motion of the matter in the shell From Eq. (5), we obtain the equations of motion of the shell, $`\delta T_{\mathrm{shell};\beta }^{\alpha \beta }=0`$, where ; denotes the covariant derivative in the Schwarzschild spacetime. Taking Eq. (6) into account leads to the restriction $`S^{r\alpha }=0`$ and to the partial differential equations, $$S_{\|b}^{ab}=0,S_{;\alpha }^{r\alpha }=0,a,b=t,\theta ,\phi $$ (8) for the components of the surface stress energy tensor. Here “<sub>\|</sub>” is the covariant derivative with respect to the shell 3-metric in Eq. (3). Due to the spherical symmetry of the shell, we can decompose $`S_{00},S_{0i},S_{ij};i,j=\theta ,\phi `$ in scalar, vector and tensor spherical harmonics respectively. Restricting attention to the even parity harmonics, we write these decompositions as: $`S_{00}={\displaystyle \underset{l}{}}S_{l0}^3(t)Y_{l0}(\theta )`$ (10) $`S_{0\theta }={\displaystyle \underset{l}{}}S_{l0}^4(t){\displaystyle \frac{}{\theta }}Y_{l0}(\theta )`$ (11) $`S_{ij}={\displaystyle \underset{l}{}}[S_{l0}^5(t)\mathrm{\Phi }_{l0ij}(\theta )+S_{l0}^6(t)\mathrm{\Psi }_{l0ij}(\theta )].`$ (12) Here $`Y_{l0}`$ is the scalar spherical harmonic and $`\mathrm{\Phi }_{l0\phi \phi }/\mathrm{sin}^2\theta =\mathrm{\Phi }_{l0\theta \theta }=Y_{l0},\mathrm{\Phi }_{l0\theta \phi }=0`$ and $`\mathrm{\Psi }_{l0\theta \theta }=^2/\theta ^2Y_{l0},\mathrm{\Psi }_{l0\phi \phi }=\mathrm{sin}\theta \mathrm{cos}\theta /\theta Y_{l0},\mathrm{\Psi }_{l0\theta \phi }=0`$ are the even parity Regge-Wheeler tensor harmonics. Azimuthal symmetry guarantees that the azimuthal index $`m`$ does not enter into any of the equations after multipole decomposition, so with no loss of generality we consider only $`m=0`$ axially symmetric motions. Upon substitution of Eq. (II A) into Eq. (8) we get, $`{\displaystyle \frac{R_{\mathrm{shell}}^2}{12M/R_{\mathrm{shell}}}}{\displaystyle \frac{dS_{l0}^3}{dt}}+l(l+1)S_{l0}^4=0`$ (14) $`{\displaystyle \frac{R_{\mathrm{shell}}^2}{12M/R_{\mathrm{shell}}}}{\displaystyle \frac{dS_{l0}^4}{dt}}S_{l0}^5+\left[l(l+1)1\right]S_{l0}^6=0`$ (15) $`{\displaystyle \frac{MR_{\mathrm{shell}}}{(12M/R_{\mathrm{shell}})^2}}S_{l0}^32S_{l0}^5+l(l+1)S_{l0}^6=0.`$ (16) As these equations show, we only have one degree of freedom. The choice of the surface mass-energy density $`S_{l0}^3(t)`$ uniquely determines all the other components of the shell’s stress energy through Eq. (II A). In this work, we make the choice $$S_{l0}^3(t)=ϵ\frac{e^{at^2}}{M},$$ (17) where $`ϵ`$ is the perturbation parameter. Since all perturbation equations will be proportional to $`ϵ`$ we will omit it henceforth. The use of a Gaussian time dependence for the surface density on the shell gives us a source that it localized in time and allows us a choice of timescale for the process that drives the stellar fluid motions. ## III Equations governing even parity perturbations A multipole decomposition of the even parity quantities in Eq. (5) leads to a set of coupled partial differential equations in the variables $`t,r`$ for the coefficients of the metric perturbation $`h_{\alpha \beta }`$ and velocity of the fluid star. We adopt the notation of Regge and Wheeler , Thorne and Campolattaro and of Moncrief . For simplicity, we make the Regge-Wheeler gauge choice, which for even parity means that the only nonvanishing metric perturbations, for a particular $`l`$, are $`h_{00}=e^{\nu (r)}H_0^{l0}(r,t)Y_{l0}(\theta )`$ (18) $`h_{0r}=H_1^{l0}(r,t)Y_{l0}(\theta )`$ (19) $`h_{rr}=e^{\lambda (r)}H_2^{l0}(r,t)Y_{l0}(\theta )`$ (20) $`h_{jk}=r^2K^{l0}(r,t)\mathrm{\Phi }_{l0jk},j,k=\theta ,\phi ,`$ (21) where $`\mathrm{\Phi }_{l0jk}`$ is one of the even parity tensor harmonics defined above, after Eq. (II A). It is useful to divide the perturbed Einstein equations (5) into those governing perturbations inside the star and those for perturbations outside the star. ### A The interior equations Inside the star the only nonzero stress energy is due to the perturbed fluid. Its independent components, are $`\delta T_0^0{}_{\mathrm{fluid}}{}^{}=\delta \rho `$ (23) $`\delta T_r^r{}_{\mathrm{fluid}}{}^{}=\delta T_\theta ^\theta {}_{\mathrm{fluid}}{}^{}=\delta T_\phi ^\phi {}_{\mathrm{fluid}}{}^{}=\delta p`$ (24) $`\delta T_0^a{}_{\mathrm{fluid}}{}^{}=(\rho +p)u_0\delta u^a,`$ (25) $`\delta T_a^0{}_{\mathrm{fluid}}{}^{}=(\rho +p)\delta u_au^0,a=r,\theta ,\phi ,`$ (26) where $`\delta \rho ,\delta p`$ are the Eulerian changes in density and pressure, $`u^0=e^{\nu (r)/2}`$ is the only nonzero component of the velocity of the unperturbed star and $`\delta u^\alpha `$ is the velocity of the perturbed fluid. It is convenient to introduce, at this point, the quantity $$\delta h\frac{\delta p}{\rho +p}.$$ (27) For barotropic fluids, $`\delta h`$ is the Eulerian perturbation of the relativistic enthalpy and $$\delta \rho =\frac{(p+\rho )^2}{p\gamma }\delta h,$$ (28) where $`\gamma `$ is the adiabatic index, $$\gamma \frac{p+\rho }{p}\frac{dp/dr}{d\rho /dr}.$$ (29) We decompose the stress energy components in Eqs. (III A) into spherical harmonics, and for a single multipole have $`\delta h=\delta h_l(r,t)Y_{l0}(\theta )`$ (30) $`\delta u^0={\displaystyle \frac{1}{2}}e^{\nu (r)/2}H_0^{l0}(r,t)Y_{l0}(\theta )`$ (31) $`\delta u^r={\displaystyle \frac{e^{[\nu (r)+\lambda (r)]/2}}{r^2}}{\displaystyle \frac{}{t}}W_{l0}(r,t)Y_{l0}(\theta )`$ (32) $`\delta u^\theta ={\displaystyle \frac{e^{\nu (r)/2}}{r^2}}{\displaystyle \frac{}{t}}V_{l0}(r,t){\displaystyle \frac{}{\theta }}Y_{l0}(\theta ).`$ (33) With a similar decomposition of the perturbed Einstein tensor into tensor harmonics, Eq. (5) leads to a set of coupled equations for $`H_0^{l0},H_1^{l0},K^{l0},H_2^{l0},W_{l0},V_{l0}`$ and $`\delta h_{l0}`$. One of the equations provides the important simplification, $$H_2^{l0}(r,t)=H_0^{l0}(r,t).$$ (34) Using this result, the other equations, can be reduced to a coupled system of two equations in which all fluid functions have been eliminated and the only dependent variables are the metric functions $`H_0^{l0},K^{l0}`$. For barotropic fluids, these equations are $`e^{[\lambda (r)3\nu (r)]/2}[e^{[3\nu (r)\lambda (r)]/2}K_{,r}^{l0}]_{,r}+2\left({\displaystyle \frac{1}{r}}+{\displaystyle \frac{1}{2}}\nu _{,r}\right)K_{,r}^{l0}e^{\lambda (r)\nu (r)}K_{,tt}^{l0}(l1)(l+2){\displaystyle \frac{e^{\lambda (r)}}{r^2}}K^{l0}H_{0,rr}^{l0}`$ (35) $`+({\displaystyle \frac{2}{r}}+{\displaystyle \frac{\lambda _{,r}}{2}}{\displaystyle \frac{5\nu _{,r}}{2}})H_{0,r}^{l0}+\{{\displaystyle \frac{l(l+1)}{r^2}}e^{\lambda (r)}{\displaystyle \frac{2}{r^2}}(1r\nu _{,r})e^{[\lambda (r)3\nu (r)]/2}[e^{[3\nu (r)\lambda (r)]/2}\nu _{,r}]_{,r}`$ (36) $`8\pi e^{\lambda (r)}(\rho +p)\}H_0^{l0}+e^{\lambda (r)\nu (r)}H_{0,tt}^{l0}=0,`$ (37) and $`K_{,rr}^{l0}+\left[{\displaystyle \frac{3}{r}}{\displaystyle \frac{\lambda _{,r}}{2}}+{\displaystyle \frac{p+\rho }{p\gamma }}\left({\displaystyle \frac{1}{r}}+{\displaystyle \frac{1}{2}}\nu _{,r}\right)\right]K_{,r}^{l0}{\displaystyle \frac{e^{\lambda (r)}}{2r^2}}(l1)(l+2)\left(1+{\displaystyle \frac{p+\rho }{p\gamma }}\right)K^{l0}e^{\lambda (r)\nu (r)}{\displaystyle \frac{p+\rho }{p\gamma }}K_{,tt}^{l0}`$ (38) $`\left(1{\displaystyle \frac{p+\rho }{p\gamma }}\right){\displaystyle \frac{H_{0,r}^{l0}}{r}}\left[{\displaystyle \frac{\lambda _{,r}}{r}}+{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{l(l+1)}{2r^2}}e^{\lambda (r)}+{\displaystyle \frac{p+\rho }{p\gamma }}\left({\displaystyle \frac{1}{r^2}}{\displaystyle \frac{l(l+1)}{2r^2}}e^{\lambda (r)}\right){\displaystyle \frac{\nu _{,r}}{r}}{\displaystyle \frac{p+\rho }{p\gamma }}\right]H_0^{l0}=0.`$ (39) For HIF models, the adiabatic index $`\gamma `$ is effectively infinite since $`\rho =`$constant, and thus all the terms proportional to $`(p+\rho )/p\gamma `$ in Eq. (39) are zero. ### B Reduction of the exterior equations to the Zerilli equation In the Schwarzschild spacetime outside the star, it can be shown that all the perturbation equations can be obtained from the two first order equations $`l(l+1)H_1^{l0}+2rH_{0,t}^{l0}2r^2K_{,tr}^{l0}+{\displaystyle \frac{6M2r}{12M/r}}K_{,t}^{l0}=D_{l0}(r,t)`$ (40) $`{\displaystyle \frac{2M}{r^2}}H_1^{l0}+\left(1{\displaystyle \frac{2M}{r}}\right)H_{1,r}^{l0}K_{,t}^{l0}H_{0,t}^{l0}=B_{l0}(r,t)`$ (41) together with the algebraic identity $`F\left[(l1)(l+2)+{\displaystyle \frac{6M}{r}}\right]H_{0,t}^{l0}\left[(l1)(l+2)+{\displaystyle \frac{2M(r3M)}{r(r2M)}}\right]K_{,t}^{l0}`$ (42) $`{\displaystyle \frac{2r^2}{12M/r}}K_{,ttt}^{l0}+2rH_{1,tt}^{l0}+M{\displaystyle \frac{l(l+1)}{r^2}}H_1^{l0}=0`$ (43) and the shell’s equations of motion, Eq. (II A). The source terms in Eqs. (40), (41) are $`D_{l0}=32\pi r{\displaystyle \frac{dS_{l0}^6}{dt}}\sqrt{1{\displaystyle \frac{2M}{r}}}\delta \left(rR_{\mathrm{shell}}\right)`$ (44) $`B_{l0}=16\pi \left(S_{l0}^4{\displaystyle \frac{dS_{l0}^6}{dt}}\right)\sqrt{1{\displaystyle \frac{2M}{r}}}\delta \left(rR_{\mathrm{shell}}\right).`$ (45) Equations (40), (41) and (43), in turn, can be combined into the single wave equation, $$\frac{^2Z_{l0}}{t^2}\frac{^2Z_{l0}}{r^2}+V_l(r)Z_{l0}=𝒮_{l0}(r,t)$$ (46) for the Zerilli function, $$Z_{l0}(r,t)\frac{r(r2M)}{(r\lambda +3M)(\lambda +1)}[H_0^{l0}rK_{,r}^{l0}]+\frac{r}{\lambda +1}K^{l0}.$$ (47) In Eq. (46), $`r^{}`$ is the usual tortoise coordinate, $$r^{}=r+2M\mathrm{log}[r/2M1]+\mathrm{constant},$$ (48) the constant $`\mathrm{\Lambda }`$ is $$\mathrm{\Lambda }\frac{(l1)(l+2)}{2},$$ (49) and the potential $`V_l`$ is $$V_l(r)=\frac{2\mathrm{\Lambda }^2(\mathrm{\Lambda }+1)r^3+6\mathrm{\Lambda }^2Mr^2+18\mathrm{\Lambda }M^2r+18M^3}{r^3(\mathrm{\Lambda }r+3M)^2}\left(1\frac{2M}{r}\right).$$ (50) The source term in Eq. (46) is $`𝒮_{l0}(r,t)={\displaystyle \frac{C_{l0}(r,t)f(r)}{h(r)}}\delta [rR_{\mathrm{shell}}]+{\displaystyle \frac{}{r^{}}}\left[{\displaystyle \frac{C_{l0}(r,t)\delta [rR_{\mathrm{shell}}]}{h(r)}}\right]`$ (51) $`+{\displaystyle \frac{16\pi }{R_{\mathrm{shell}}}}\left(1{\displaystyle \frac{2M}{R_{\mathrm{shell}}}}\right)^{3/2}S_{l0}^6(t)\delta \left(rR_{\mathrm{shell}}\right),`$ (52) where $`C_{l0}(r,t)={\displaystyle \frac{16\pi R_{\mathrm{shell}}^2}{l(l+1)\sqrt{12M/R_{\mathrm{shell}}}}}S_{l0}^3(t)`$ (53) $`f(r)={\displaystyle \frac{6M^2+3\mathrm{\Lambda }Mr+r^2\mathrm{\Lambda }(\mathrm{\Lambda }+1)}{r^2(3M+\mathrm{\Lambda }r)}}`$ (54) $`h(r)={\displaystyle \frac{3M+\mathrm{\Lambda }r}{r2M}}.`$ (55) Together with appropriate boundary conditions, the set of equations (37), (39) and (46) govern the even parity perturbations of the star due to the shell. ### C Boundary conditions #### 1 Regularity at the center of the star The system of equations (37-39) admits two linearly independent regular solutions. Near the center they admit the series expansion, $`K^{l0}(r,t)=k_0(t)r^l+k_2(t)r^{l+2}+O(r^{l+4})`$ (57) $`H_0^{l0}(r,t)=k_0(t)r^l+h_2(t)r^{l+2}+O(r^{l+4}).`$ (58) The $`r`$-independent “constants,” $`k_0`$, $`k_2`$, $`h_2`$ are a priori arbitrary. When Eq. (39) is expanded about $`r=0`$ and the above expansions are used, we find $$\left[\frac{3}{2}l+3+\frac{l^2}{2}\right]h_2+\left[\frac{l^2}{2}+\frac{11}{2}l+9\right]k_2\frac{8\pi }{3}\rho _0\left[2l+l^23\right]k_0=0$$ (59) where $`\rho _0,p_0,\gamma _0`$ are the values of the density and pressure at the center of the star and we used the fact that $`\gamma =\mathrm{}`$ for a HIF. Equation (59) allows us to eliminate $`k_0`$, so that the only remaining unknown coefficients are $`k_2`$ and $`h_2`$. These turn out to be fixed, up to an overall scaling, by conditions at the stellar surface. #### 2 Vanishing of the pressure at the stellar surface The pressure must vanish at the perturbed surface of the star by definition of that boundary. This amounts to the vanishing of the Lagrangian perturbation of the pressure at the stellar surface $`r=R`$, $$\delta p(R)\frac{e^{\lambda (R)/2}}{R^2}W(R,t)p_{,r}(R)=0.$$ (60) When the density $`\rho `$ vanishes at $`r=R`$, the second term in Eq. (60) is zero and this condition is the same as the vanishing of the Eulerian perturbation of the pressure. But this is not so when $`\rho (R)0`$ as is the case for a HIF. After differentiating Eq. (60) twice with respect to time and doing a decomposition in spherical harmonics, we can rewrite it purely in terms of the metric functions $`K^{l0},H_0^{l0}`$ and its derivatives. The resulting expression, when $`\rho (R)0`$, is $`\{{\displaystyle \frac{1}{r}}K_{,ttr}^{l0}{\displaystyle \frac{\nu _{,r}}{4}}{\displaystyle \frac{l(l+1)}{r^2}}e^{\nu (r)}K_{,r}^{l0}{\displaystyle \frac{1}{2}}{\displaystyle \frac{(l1)(l+2)}{r^2}}e^{\lambda (r)}K_{,tt}^{l0}e^{\lambda (r)\nu (r)}K_{,tttt}^{l0}{\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu _{,r}}{r}}(1r{\displaystyle \frac{\nu _{,r}}{2}})K_{,tt}^{l0}`$ (61) $`+{\displaystyle \frac{1}{r}}H_{0,rtt}^{l0}+{\displaystyle \frac{l(l+1)}{4r^2}}e^{\nu (r)}\nu _{,r}H_{0,r}^{l0}+[3{\displaystyle \frac{\nu _{,r}}{2r}}{\displaystyle \frac{1}{r^2}}+{\displaystyle \frac{l(l+1)}{2r^2}}e^{\lambda (r)}]H_{0,tt}^{l0}+(\nu _{,r})^2e^{\nu (r)}{\displaystyle \frac{l(l+1)}{4r^2}}H_0^{l0}\}_{r=R}=0.`$ (62) #### 3 The relation between the interior and the exterior metric functions If the density of the star is not zero at $`r=R`$, the first radial derivative of $`K^{l0},H_0^{l0}`$ is not continuous at $`r=R`$, although the functions themselves are. Denoting by a superscript “$`+`$” the exterior metric functions and by a superscript “$``$” the interior ones, we have, for a particular multipole $`l`$, that $`K^+(R,t)=K^{}(R,t)`$ (64) $`H_0^+(R,t)=H_0^{}(R,t)`$ (65) $`K_{,rtt}^+(R,t)=K_{,rtt}^{}(R,t){\displaystyle \frac{1}{2R^2}}\{l(l+1)(1{\displaystyle \frac{2M}{r}})K_{,r}^{}+2r[{\displaystyle \frac{r3M}{r2M}}K_{,tt}^{}+rK_{,rtt}^{}]`$ (66) $`l(l+1)(1{\displaystyle \frac{2M}{r}})H_{0,r}^{}2rH_{0,tt}^{}2M{\displaystyle \frac{l(l+1)}{r^2}}H_0\}_{r=R}.`$ (67) If $`\rho (R)=0`$ then both metric functions and their first radial derivatives would be continuous at $`r=R`$. ### D Fourier transform of the perturbation equations The simplest way to solve the above partial differential equations in $`r,t`$ is to write all time dependent quantities as Fourier integrals, reducing the problem to that of ordinary differential equations in $`r`$. Thus we write $`Z_{l0}(r,t)`$ as, $$Z_{l0}(r,t)=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}e^{i\omega t}\stackrel{~}{Z}_{l0}(r,\omega )𝑑\omega $$ (68) transforming Eq. (46) into a second order equation for $`\stackrel{~}{Z}_{l0}(r,\omega )`$, $$\frac{d^2\stackrel{~}{Z}_{l0}}{dr^2}+[\omega ^2V_l(r)]\stackrel{~}{Z}_{l0}=\stackrel{~}{𝒮}_{l0}(r,\omega )$$ (69) where $$\stackrel{~}{𝒮}_{l0}(r,\omega )=_{\mathrm{}}^{\mathrm{}}𝒮_{l0}(r,t)e^{i\omega t}𝑑t$$ (70) is the Fourier transform of the source term. We can proceed similarly with Eqs. (37) and (39). The resulting equations for $`\stackrel{~}{H}_0^{l0}(r,\omega ),\stackrel{~}{K}^{l0}(r,\omega )`$ can be obtained directly from Eqs. (37) and (39) by substituting $`i\omega `$ for $`_t`$, and by replacing functions of $`r,t`$ by their Fourier transforms. At infinity, we impose the boundary condition that the wave be purely outgoing $$Z_{l0}(r\mathrm{},utr^{})=\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}A_{l0}(\omega )e^{i\omega u}𝑑\omega .$$ (71) The even parity gravitational energy, radiated in a single multipole component, is then given by $$\frac{dE_l}{d\omega }=\frac{1}{64\pi ^2}\frac{(l+2)!}{(l2)!}\omega ^2\|A_{l0}(\omega )\|^2.$$ (72) From the wave function $`Z_{l0}(r,t)`$ we can construct the multipoles of the metric perturbations. The transverse traceless (TT) part of the metric perturbation is of particular interest, since it completely characterizes the radiation at infinity. Extracting the TT part is easiest if we realize that the Zerilli function $`Z_{l0}(r,t)`$ is numerically equal to Moncrief’s gauge invariant wavefunction. Since the Moncrief invariant can be evaluated in any gauge, we choose the gauge to be asymptotically flat so that all multipole perturbations fall off faster than $`1/r`$ except the TT components. From this procedure we find that the TT perturbations are related to the Zerilli function by $$h_{jk}^{TT}=\frac{1}{r}\underset{l=2}{}Z_{l0}(r,t)\sqrt{\frac{(l+2)!}{2(l2)!}}T_{jk}^{E2,l0}(\theta ),j,k=\theta ,\phi ,$$ (73) where $$T_{jk}^{E2,l0}(\theta )=\sqrt{2\frac{(l2)!}{(l+2)!}}\left[\mathrm{\Psi }_{l0jk}+\frac{l(l+1)}{2}\mathrm{\Phi }_{l0jk}\right]$$ (74) is the orthonormal even parity TT tensor harmonic and $`\mathrm{\Psi }_{l0},\mathrm{\Phi }_{l0}`$ are the tensor harmonics introduced in Eq. (II A). ## IV Computational implementation ### A Solution for $`A_{l0}(\omega )`$ If the background spacetime is due to a star, a solution of Eq. (69) must be found that corresponds to outgoing waves at infinity and that matches the regular solution of Eq. (37), (39) at the unperturbed surface of the star according to the junction conditions of Eq. (III C 3). The Green function solution is found in the usual way. (See e.g. , Ref. .) We define $`y_l^{out}(r,\omega )`$ as the homogeneous solution of Eq. (69) with the asymptotic form $$y_l^{out}(r,\omega )e^{i\omega r^{}}r^{}\mathrm{}.$$ (75) For our second independent solution of Eq. (69) we start by finding, in the stellar interior, a solution of Eqs. (37), (39) satisfying the condition Eq. (III C 1) at the stellar center. Our second solution $`y_l^{reg}(r,\omega )`$ is taken to be the homogeneous solution of Eq. (69) that joins to the interior solution through the matching conditons of Eq. (III C 3). We then define the Wronskian of these two homogeneous solutions, an $`r`$ independent quantity, to be $$W_l(\omega )=y_l^{reg}\frac{dy_l^{out}}{dr^{}}y_l^{out}\frac{dy_l^{reg}}{dr^{}}.$$ (76) With the above definitions, and from the Green function solution, we obtain (see, e.g. ,Ref. ) the Fourier amplitude $`A_{l0}`$ defined in Eq. (71) $$A_{l0}(\omega )=\frac{1}{W_l(\omega )}_{r<R_{\mathrm{shell}}}^{\mathrm{}}\stackrel{~}{𝒮}_{l0}(r,\omega )\frac{y_l^{reg}(r,\omega )}{12M/r}𝑑r.$$ (77) Combining Eq. (II A), (17), (52) and (70) we get explicitly, $`A_{l0}(\omega )={\displaystyle \frac{16\pi }{W_l(\omega )\sqrt{12M/R_{\mathrm{shell}}}}}\sqrt{{\displaystyle \frac{\pi }{a}}}e^{\omega ^2/(4a)}\{{\displaystyle \frac{R_{\mathrm{shell}}^2}{Mh(r)l(l+1)}}{\displaystyle \frac{dy^{reg}}{dr}}(R_{\mathrm{shell}},\omega )`$ (78) $`+{\displaystyle \frac{y_l^{reg}(R_{\mathrm{shell}},\omega )}{12M/R_{\mathrm{shell}}}}[{\displaystyle \frac{1}{(l1)(l+2)}}(1{\displaystyle \frac{2R_{\mathrm{shell}}^3\omega ^2}{Ml(l+1)}})+{\displaystyle \frac{R_{\mathrm{shell}}^2f(r)}{Mh(r)l(l+1)}}]\}.`$ (79) What is required for a solution, then, is a numerical determination of $`y_l^{reg}(R_{\mathrm{shell}},\omega )`$ and its radial derivative. ### B Numerical method to find $`y^{reg}`$ for a HIF The numerical problem of finding $`y_l^{reg}(R_{\mathrm{shell}},\omega )`$ and its derivative for a HIF, can be divided in two parts: Integration of Eqs. (37) and (39) from the center of the star to $`r=R`$ and integration of Eq. (69) from $`r=R`$ to $`r=R_{\mathrm{shell}}`$, so that we can evaluate $`y_l^{reg},dy_l^{reg}/dr`$ at this radius. To find $`K^{}(R,\omega ),H_0^{}(R,\omega )`$ and $`K_{,r}^{}(R,\omega )`$, for a HIF we first find an “A solution” by starting from the center with with $`k_0=1,k_2=0`$ and $`h_2`$ obtained from Eq. (59); we then integrate Eqs. (37), (39) out to $`r=R`$, to find the “A solution” there. The procedure for the “B solution” at $`r=R`$ is the same except we start with the central conditions $`k_0=0,k_2=1`$. The general solution for $`H_0^{}`$ and $`K^{}`$ can be written $$H_0^{}=\alpha H_0^A+\beta H_0^BK^{}=\alpha K^A+\beta K^B.$$ (80) where $`\alpha `$ and $`\beta `$ are arbitrary constants. The overall scale of $`y_l^{reg}(r,\omega )`$ is arbitrary. Since $`y_l^{reg}(r,\omega )`$ occurs both in the integrand in Eq. (77) and in the Wronskian in the denominator, the scale cancels out so only the ratio $`\alpha /\beta `$ is of importance. This ratio can be found by substituting the expression in Eq. (80) and its derivatives in the vanishing of the Lagrangian pressure condition Eq. (62). We can next compute $`K^+(R,\omega ),H^+(R,\omega )`$ and $`K_{,r}^+(R,\omega )`$ using Eq. (III C 3). The final step is to use Eq. (47) to find the starting values for integrating the Zerilli equation, $$y_l^{reg}(R,\omega )=\frac{R(R2M)}{(R\mathrm{\Lambda }+3M)(\mathrm{\Lambda }+1)}[H_0^{l0+}rK_{,r}^{l0+}]_{r=R}+\frac{r}{\mathrm{\Lambda }+1}K^{l0+}(R,\omega )$$ (81) $$\frac{dy_l^{reg}(R,\omega )}{dr}=\frac{3M}{(R\mathrm{\Lambda }+3M)(\mathrm{\Lambda }+1)}K^{+l0}(R,\omega )+\frac{R[6M^2+3M\mathrm{\Lambda }R+R^2\mathrm{\Lambda }(\mathrm{\Lambda }+1)]}{(R\mathrm{\Lambda }+3M)^2(\mathrm{\Lambda }+1)}\left[K_{,r}^{l0+}\frac{H_0^{l0+}}{R}\right]_{r=R_{\mathrm{shell}}},$$ (82) and to integrate the Zerilli equation Eq. (46) out to $`r=R_{\mathrm{shell}}`$ in order to find $`y_l^{reg}(R_{\mathrm{shell}},\omega )`$ and $`dy_l^{reg}/dr(R_{\mathrm{shell}},\omega )`$. ### C Radiation due only to the shell The waveform $`Z_{l0}`$, its Fourier transform $`A_{l0}(\omega )`$, and the energy computed in Eq. (72) refer, of course, to radiation from the star and the shell. As might be expected the radiation that can be attributed directly to the shell is much larger than the radiation from the perturbed fluid motions in the neutron star. In order to have the clearest comparison of Newtonian and relativistic predictions of radiation from the neutron star, it is useful to remove the contribution due to the shell. For the Newtonian computation presented in the appendix, this presents no problem. As is evident in Eq. (A22), the density perturbations due to the shell and due to the star are distinct, and it is evident in Eq. (A25) that their contribution to the quadrupole moment are distinct. This clear distinction does not exist in the relativistic calculation. The equations of Sec. III contain information about the shell entangled with information about the star. To get an approximate idea of what radiation can be ascribed to the star, and not to the shell, one can compute the waveform due (in a sense) to the shell itself, and subtract this waveform from the total star+shell waveform. This subtraction, however, is somewhat subtle. In particular, it is not useful to consider the shell in a flat spacetime background. If we consider a shell with the surface stress energy of Eqs. (II A), (17) radiating in a flat background, and the same (in some sense) shell, and surface stresses, radiating in another fixed background, there will be a large difference simply due to the different radial null geodesics, the spacetime lines on which perturbations propagate. If we are to have the same shell-only radiation as that due to the shell in the shell+star problem, we must have the two signals propagate on the same spacetime background. To accomplish this, we describe the waves at radius $`r>R`$ for the shell-only problem with the same Zerilli equation as we use in the shell+star problem. Equation (46) is then part of both problems for $`r>R`$. In the shell+star problem the perturbations in the interior are, of course, treated with the equations of Sec. III. For the shell-only problem, we need instead something like a Zerilli equation suitable to the fixed spacetime of the stellar interior. Such an equation is provided by the equation for the propagation of massless scalar perturbations in the neutron star background. These waves do not excite oscillations either of the spacetime or of the neutron star fluid. They are described only by an equation identical in form to Eq. (46), but with no source, and with the Zerilli potential of Eq. (50) replaced by $$V_l^{\mathrm{scalar}}=e^\nu \left[\frac{l(l+1)}{r^2}+\frac{\nu _{,r}\lambda _{,r}}{2r}e^\lambda \right].$$ (83) (The idea of freezing the fluid perturbations to study the $`w`$ modes described in is somewhat similar in spirit to our method but is different in practice.) With this mixed mathematical description we find the regular solution of the Zerilli equation and then compute, using Eq. (79), the Fourier transform $`A_{l0}^{\mathrm{shell}}`$ and, with Eq. (71), a wave function $`Z_{l0}^{\mathrm{shell}}`$. By subtracting these, respectively from the $`A_{l0}`$ and $`Z_{l0}`$ for the shell+star, we arrive at results $`A_{l0}^{\mathrm{star}}`$, $`Z_{l0}^{\mathrm{star}}`$ meant to describe radiation only due to the oscillations of the neutron star fluid. From the square of $`A_{l0}^{\mathrm{star}}`$, we can compute an energy spectrum and a total radiation energy attributed to the motion of the stellar fluid. Our method of finding the radiation from the shell alone is, of course, only an approximation. Einstein’s equations couple oscillations of the fluid and oscillations of the spacetime, so that there can be no completely meaningful way of separating the two. One sign of this is that the potential Eq. (83) is not unique. The potential does not influence the radial null geodesics, and any potential with the same general behavior at $`r0`$, is equally good. In situations for which the details of the potential are important, our method of computing shell-only radiation is not justified. But such details are not relevant for must of our models. This can be seen in the reasonable success of our method in removing appearances of shell radiation in the star-only results presented below. An additional important point to understand about our shell subtraction method is that in principle it should subtract the $`w`$ modes. Since the $`w`$ modes are due to the spacetime background, not to the fluid motion, the shell-only radiation should have $`w`$ modes, and the star-only waves that result from subtraction should have only waves due to fluid motion. This will be discussed further in connection with the numerical results presented below. ## V Numerical results We focus first on a very compact HIF stellar model of radius $`R=2.5M`$. By starting with an extreme, though astrophysically implausible, model we will be able to see relativistic features that will be absent in less compact, and more plausible, models. In Fig. 1 (a) we present the waveform $`Z_{20}`$, the Zerilli function for quadrupole ($`l=2`$) perturbations in the case that the Gaussian parameter in Eq. (17) is $`a=0.1M^2`$. The shell was placed at a large distance from the star, $`R_{\mathrm{shell}}=110M`$, to have the waveform clearly display the time profile of events. The first burst, at around $`u=tr^{}=160M`$ is radiation coming directly from the stress energy of the shell. Later, at around $`u=tr^{}100M`$, a burst arrives representing the ingoing radiation from the shell that has “reflected” off $`r=0`$. At around the same time, radiation from oscillations of the fluid and central region of the spacetime arrive. Two curves are shown in Fig. 1 (a). One is the waveform for the relativistic star+shell. The second curve is that for the shell itself, computed by the method described in Sec. IV C. The two curves are nearly identical up to around $`u=160M`$, confirming that the early radiation is that due to the shell. From $`u=160M`$ to around $`400M`$ or more, there are damped oscillations of the dominant (least damped) $`w`$ mode. The $`w`$ modes depend on the details of wave propagation in the innermost strong field regions, and are not the same for the star+shell problem and for the somewhat ad hoc spacetime we constructed for the shell-only problem. The $`w`$ mode frequency for the constant density $`R=2.5M`$ model is $`\omega =(0.42+i0.0217)M^1`$, while that for the spacetime of Sec. IV C, $`\omega =(0.438+i0.029)M^1`$, is more rapidly damped. We would, of course, find a different $`w`$ mode frequency if in the shell-only problem we used a potential other than that in Eq. (83). It is clear that there is no point in subtracting the shell-only radiation from the star+shell. The difference waveform would contain an artifact corresponding to the difference of the $`w`$ modes, and hence would be dominated by features completely irrelevant to fluid motions of the neutron star. Figure 1 (a) helps to demonstrate that subtraction is pointless when $`w`$ mode radiation is of importance in the waveform. It is only the fact that $`w`$ mode radiation is a minor feature for realistically compact stars that makes subtraction useful. In addition to the difference in frequency of the $`w`$ modes there is another, more important, difference between the two curves in Figure 1 (a). The star+shell curve shows an oscillation with imperceptible damping at a frequency less than half that of the $`w`$ mode. This is the $`f`$ mode of the fluid of the neutron star. This mode is missing, as it should be, from the shell-only computation. Both curves in Fig. 1(a) show relativistic results. A comparison between the Newtonian and relativistic star+shell results would not be of much use. The slow motion condition underlying the Newtonian approximation would be strongly violated since the radius of the shell $`R_{\mathrm{shell}}=110M`$ is much larger than the characteristic time scale of the shell stress-energy oscillations ($`\mathrm{\Delta }ta^{1/2}`$ several $`M`$) or of the modes of oscillation of the spacetime or of the stellar fluid. The Newtonian computation, based on the slow motion approximation (time scale $``$ light travel time across source) would be completely inappropriate. In order to construct a more justifiable Newtonian comparison for a $`R=2.5M`$ HIF model, we consider in Fig. 1(b), a shell with radius $`R_{\mathrm{shell}}=5M`$ and with Gaussian parameter $`a=0.001M^2`$ (and hence timescale $`30M`$). The waveform in this case is not of primary interest, since the radiation from the shell and from the stellar fluid will not be clearly distinguishable. The energy spectrum, however, gives the answer to the most important questions. It shows, for example, that there is negligible difference between the broad spectra of the Newtonian and the relativistic results. But the broad spectrum is due to the shell. Of more astrophysical interest is the $`f`$ mode excitation. The relativistic computation shows a lower frequency $`f`$ mode containing almost an order of magnitude more radiation energy than the $`f`$ mode peak of the Newtonian computation. This conclusion, of course, applies to a model that is too compact to be astrophysically relevant. The feature in the relativistic spectrum at $`\omega 0.3M^1`$ is due to the repetition of the shell radiation with a time delay of $`\mathrm{\Delta }t10M`$. This produces a modulation of the form $`\mathrm{cos}^2(\omega \mathrm{\Delta }t/2)`$, and hence a dramatic decrease at $`\omega \pi /\mathrm{\Delta }t0.3`$. The location of such features is dependent on $`R_{\mathrm{shell}}`$ and, of course, is unrelated to the physics of the neutron star. It is worthwhile noting that these features are absent in the Newtonian spectrum. Since that spectrum is based on the slow motion approximation the whole star+shell source is treated as if it radiates in phase, and there can be no repetition of the shell radiation. In Fig. 2 we show results for a HIF star of radius $`R=5M`$, a typical radius of a neutron star. The shell is located at $`R_{\mathrm{shell}}=10M`$ and the Gaussian parameter for the time profile of the shell stress energy is $`a=0.01M^2`$. In Fig. 2(a) the dotted curve gives the full Zerilli quadrupole waveform of the star+shell computation, and shows why subtraction of shell radiation is useful: the waveform is completely dominated by the shell radiation, and the much smaller $`f`$ mode is barely visible. The solid curve is the waveform from the shell-only computation. It is clear in Fig. 2(a) that the two waveforms are nearly identical at early times, and at late times are different in that the star+shell result has $`f`$ mode oscillations, while the shell-only waveform, of course, does not. This suggests that subtraction will be very effective in isolating the radiation due to the stellar fluid, and this turns out to be true. The star-only curve (the result of subtracting the shell-only from the star+shell) is given in the inset to Fig. 2(a) and shows the $`f`$ mode oscillation as a dominant feature. Figure 2(b) shows the energy spectrum for the subtracted (i.e. , star-only) case shown in the inset of Fig. 2(a). Both the relativistic and Newtonian spectra are given. The two spectra are roughly similar in general appearance, except for the strong feature in the relativistic spectrum at $`\omega 0.5M^1`$. Of particular interest in Fig. 2(b) is the excitation of the $`f`$ mode (shown in detail in the inset). Perhaps the most important feature of our model is that we can compare the excitation given by a relativistic and a Newtonian computation. In order to illustrate how the methods of this paper work for a very weakly relativistic model, Fig. 3 shows two spectra for a HIF star with $`R=20M`$. In both, the Newtonian star-only spectrum is compared with the relativistic star-only (i.e. , subtracted) spectrum, and for both $`R_{\mathrm{shell}}=22M`$. Figure 3(a) shows the case for a Gaussian parameter $`a=0.01M^2`$, while (b) shows the case $`a=0.00001M^2`$. The Newtonian and relativistic computations of $`f`$ mode excitation (shown in the inset) agree quite well in Fig. 3(a), and there is general agreement in the overall shape of the spectrum but, as in Fig. 2(b), the relativistic case has structure that is missing in the Newtonian case. This is due to the fact that the shell stress energy source has a timescale ($`a^{1/2}=10M`$) that is less than the light travel time across the star, and the star is not radiating in phase. By comparison, in Fig. 3(b) the shell timescale is $`a^{1/2}300M`$ and the slow approximation is justified. (Note that the $`f`$ mode oscillations, at $`\omega 0.01M^1`$ are also slow compared to the light travel time across the star.) Because of this, the Newtonian and relativistic spectra in Fig. 3(b) are in excellent overall agreement. ## VI Conclusions We have presented a method of probing the gravitational wave properties of neutron stars by using time varying stress energy in a spherical shell. In particular, we have shown that this method can give more useful answers about the neutron star physics than those given by studies of the scattering of gravitational waves or by the response of the star to a close particle orbit. The main motivation for considering such a probe is to compare relativistic and Newtonian computations of gravitational radiation. We have shown that the shell probe is well suited for this purpose, since both the Newtonian and relativistic computations can be carried out for the model. Two additional features of the shell probe have been shown to be important or useful. One is the possibility of approximately distinguishing the radiation that can be ascribed to the star from the radiation due to the shell. Since the radiation from the shell is stronger than that from the the stellar fluid, this separation is valuable in bringing out the physical radiation that is of primary interest. A fundamentally important feature of the shell probe model, is the ability to choose the timescale of the shell stress energy. This has allowed us to direct attention to the fact that the Newtonian approximation is not only a weak field approximation, but a “slow” approximation. That is, the quadrupole approximation used in a quasi-Newtonian gravitational wave calculation supposes that the light travel time across the source is much smaller that the period of the waves generated. We have shown that even for a weakly relativistic stellar model there is not good agreement in the details of the Newtonian and relativistic computations if the timescale for the shell stress energy is short. The results shown in the previous section make this clear. The structure and dynamics of the $`R=20M`$ model of Fig. 3(a) is well described by Newtonian physics, but there are large differences between the Newtonian and relativistic spectra when the star is excited by a short time scale perturbation. The difference between the Newtonian and relativistic results is even larger for the $`R=5M`$ model of Fig. 2(b). For smaller Gaussian parameters $`a`$ (i.e. , for driving perturbations with longer time scales) the “relativistic-only” structure seen in Figs. 2(b) and 3(a) decreases. What becomes clear from these results is that the question of whether Newtonian physics is adequate for neutron star dynamics is inseparable from the question of the timescale of the excitation of the neutron star. If the timescale is imposed from a distance many times the neutron stars radius, then the excitation will be slow and our results (based on a very limited exploration of models) suggest that Newtonian physics will suffice. On the other hand, rapid processes, due to impacts, collapse, etc., may be have a timescale only several times $`GM/c^3`$, the slow approximation may be violated, and Newtonian calculations may be significantly in error. As explained in Sec. I, the motivation for the shell probe for neutron star oscillations is a sequence of two questions. In Ref. Andersson and Kokkotas questioned whether Newtonian physics was adequate for neutron star physics. The second question is whether this can be adequately studied with the particle-scattering model and its inflexible timescales. The results presented here suggest that in cases in which neutron stars are excited on a very short time scale, those particle-scattering computations are not a sufficient basis for the conclusion that relativistic effects are unimportant in neutron star models of gravitational radiation sources. Our main purpose here has been to introduce the shell probe, and the motivation for it. We have applied it only to a single simple neutron star model. It is quite possible, of course, that for some equations of state the importance of relativistic effects might be quite different. If other equations of state are to be studied towards this end, we suggest that the method of time varying stress energy in a shell be considered as a good way of getting the clearest comparison of Newtonian and relativistc predictions. ## VII Acknowledgments We thank Nils Andersson for supplying the frequencies of the even parity modes of homogeneous stars. This work was partially supported by NSF grant PHY9734871. Z. A. was supported by PRAXIS XXI/BD/3305/94 grant from FCT (Portugal). ## A The newtonian limit of a HIF ### 1 The newtonian perturbation equations and their solution We consider a HIF, excited by a spherical shell whose mass is much smaller than the star’s mass. Assuming that the motion of the fluid consists of small perturbations around a spherically symmetric and static fluid ball (the equilibrium star), we can decompose all perturbative scalar quantities, namely the gravitational potential, the pressure, and the density perturbations, in scalar spherical harmonics, and can decompose the fluid’s velocity in even parity vector spherical harmonics. We write these decompositions as $`U=U_{eq}(r)+{\displaystyle \underset{l}{}}\delta U_l(r,t)Y_{l0}(\theta )`$ (A1) $`p=p_{eq}(r)+\rho _{eq}{\displaystyle \underset{l}{}}\delta h_l(r,t)Y_{l0}(\theta )`$ (A2) $`\rho =\rho _{eq}+{\displaystyle \underset{l}{}}\left\{{\displaystyle \frac{e^{at^2}}{M}}\delta [rR_{\mathrm{shell}}]{\displaystyle \frac{W_l(r,t)}{r^2}}\rho _{eq}\delta [rR]\right\}Y_{l0}(\theta )`$ (A3) $`v^r={\displaystyle \frac{1}{r^2}}{\displaystyle \underset{l}{}}{\displaystyle \frac{}{t}}W_{l0}(r,t)Y_{l0}(\theta )`$ (A4) $`v^\theta ={\displaystyle \frac{1}{r^2}}{\displaystyle \underset{l}{}}{\displaystyle \frac{}{t}}V_{l0}(r,t){\displaystyle \frac{}{\theta }}Y_{l0}(\theta ),`$ (A5) where $`U_{eq},p_{eq},\rho _{eq}=\mathrm{const}`$ are the gravitational potential, pressure and constant density of the spherical (equilibrium) star. The perturbation of the density has two contributions: one from the matter in the shell and the other from the star itself. In writing the first one, we supposed that $`\delta T_{00,\mathrm{shell}}\delta \rho _{\mathrm{shell}}`$ is dominant over all the other components of the stress energy tensor of the shell (weak field, slow motion approximations) and that $`\delta T_{00,\mathrm{shell}}`$ is given by Eqs. (6) and (17). The density due to the fluid perturbation is zero everywhere inside the star (since the star is incompressible), but not at the unperturbed surface of the star. In writing the perturbation of the pressure, we used the definition Eq. (27) and the fact that $`p_{eq}\rho _{eq}`$ in the weak field limit. We then substitute these expressions in the fluid equations for a perfect barotropic fluid, which are Poisson’s equation for the gravitational potential, $$^2U=4\pi \rho ,$$ (A6) the continuity equation, $$\rho _{,t}=\stackrel{}{}.(\rho \stackrel{}{v}),$$ (A7) and Euler’s equation, $$\stackrel{}{v}_{,t}+(\stackrel{}{v}.\stackrel{}{})\stackrel{}{v}=\stackrel{}{}U\frac{1}{\rho }\stackrel{}{}p,$$ (A8) and we keep terms only to first order in the perturbations. The resulting linearized fluid equations can then be reduced to two equations: one for $`\delta U_l`$ and another for $`\delta h_l`$. These equations are the Newtonian limit of the relativistic even parity perturbation equations derived by Lindblom et al.. In the special case of a HIF, they reduce to two decoupled second order equations, $$\delta U_{l,rr}+\frac{2}{r}\delta U_{l,r}\frac{l(l+1)}{r^2}\delta U_l=4\pi \left\{\frac{e^{at^2}}{M}\delta [rR_{\mathrm{shell}}]\frac{W_{l0}(r,t)}{r^2}\rho _{eq}\delta [rR]\right\}$$ (A9) and $$\delta h_{l,rr}+\frac{2}{r}\delta h_{l,r}\frac{l(l+1)}{r^2}\delta h_l=0.$$ (A10) Since the left hand sides of Eqs. (A9) and (A10) are equivalent to the multipole decomposition of the Laplacian, it is straightforward to write down the solutions that are well behaved at the center of the star and that vanish at infinity: $$\delta U_l(r,t)=\{\begin{array}{cc}\alpha _l(t)r^l+\frac{4\pi }{2l+1}\frac{r^l}{MR_{\mathrm{shell}}^{l1}}e^{at^2},\hfill & rR\hfill \\ & \\ \frac{\beta _l}{r^{l+1}}+\frac{4\pi }{2l+1}\frac{r^l}{MR_{\mathrm{shell}}^{l1}}e^{at^2},\hfill & RrR_{\mathrm{shell}}\hfill \\ & \\ \frac{\beta _l}{r^{l+1}}+\frac{4\pi }{2l+1}\frac{R_{\mathrm{shell}}^{l+2}}{Mr^{l+1}}e^{at^2},\hfill & rR_{\mathrm{shell}}\hfill \end{array}$$ (A11) and $$\begin{array}{cc}\delta h_l(r,t)=\mu _l(t)r^l,\hfill & r<R.\hfill \end{array}$$ (A12) The three constants can be easily determined, by requiring the vanishing of the Lagrangian pressure at $`r=R`$, as in Eq. (60), by requiring the continuity of $`\delta U_l`$ at the surface of the star, and by integrating Eq. (A9) about $`r=R`$ and by using the linearized Euler equation for the radial component of the fluid’s velocity. The result is $`\beta _l=R^{2l+1}\alpha _l`$ (A14) $`\mu _l=(2l+1){\displaystyle \frac{\alpha _l}{3}}`$ (A15) $`{\displaystyle \frac{d^2\alpha _l}{dt^2}}+\omega _l^2\alpha _l={\displaystyle \frac{16\pi ^2}{(2l+1)^2MR_{\mathrm{shell}}^{l1}}}l\rho _{eq}e^{at^2}`$ (A16) $`\omega _l^2={\displaystyle \frac{8\pi \rho _{eq}}{3}}{\displaystyle \frac{l(l1)}{2l+1}}.`$ (A17) The frequency $`\omega _l`$ is the $`f`$ mode of vibration of the star. ### 2 Damping of the $`f`$ mode oscillation Once set into vibration at its $`f`$ mode frequency by the shell perturbation, the star would oscillate forever, since there is no damping mechanism for the perfect fluid in Newtonian theory. This is clear from Eq. (A16), which is the equation of an undamped harmonic oscillator. But in practice the star will radiate away its vibrational energy through gravitational waves, over a long period of time. The time, $`\tau _l`$, for gravitational wave damping of the $`f`$ mode oscillation can be computed, in the weak field, slow motion, approximation using energy conservation to be, $$\tau _l=\frac{4l(l1)^2(2l+1)[(2l1)!!]^2}{3(l+1)(l+2)\omega _l^{2l+2}R^{2l+1}}.$$ (A18) To introduce this damping we replace Eq. (A16) by $$\frac{d^2\alpha _l}{dt^2}+\omega _l^2\alpha _l+\frac{2}{\tau _l}\frac{d\alpha _l}{dt}=\frac{16\pi ^2}{(2l+1)^2MR_{\mathrm{shell}}^{l1}}l\rho _{eq}e^{at^2}.$$ (A19) (For a similar procedure see and ). The retarded solution of this equation is $$\alpha _l(t)=\frac{16\pi ^2l\rho _{eq}}{(2l+1)^2MR_{\mathrm{shell}}^{l1}\omega _l}_0^{\mathrm{}}𝑑ve^{v/\tau _l}e^{a[tv]^2}\mathrm{sin}[\omega _lv].$$ (A20) From Eq. (A14) and the linearized Euler equation we obtain $$\frac{W_l(r,t)}{R^2}\rho _{eq}=\frac{2l+1}{4\pi }R^{l1}\alpha _l.$$ (A21) Combined with Eq. (A20) and Eq. (A3) for $`l=2`$, this leads to the following expression for the quadrupole perturbation of the density $$\delta \rho _2(r,t)=\frac{e^{at^2}}{M}\delta [rR_{\mathrm{shell}}]+\frac{8\pi R\rho _{eq}}{5\omega _2MR_{\mathrm{shell}}}_0^{\mathrm{}}𝑑ve^{v/\tau _2}e^{a[tv]^2}\mathrm{sin}[\omega _2v]\delta [rR].$$ (A22) ### 3 The energy radiated by a vibrating, semi-Newtonian HIF The energy radiated in gravitational waves by an oscillating Newtonian star can be computed by regarding the star’s gravitational field as a small perturbation of Minkowski spacetime. For details see Ref. . The energy radiated in the quadrupole is $$\frac{dE_2}{d\omega }=\frac{1}{32\pi ^2}\omega ^6\|A_{20}(\omega )\|^2$$ (A23) where $`A_{20}`$ is the Fourier amplitude of the mass quadrupole, $$A_{20}(\omega )=_{\mathrm{}}^{\mathrm{}}𝑑ue^{i\omega u}I_{20}(u).$$ (A24) In the slow motion approximation, $`I_{20}(u)`$ is simply $$I_{20}(u)=\frac{16\pi }{5\sqrt{3}}_0^{\mathrm{}}\delta \rho _2(u,r)r^4𝑑r.$$ (A25) In this approximation, Eqs. (A22) and (A24) can be combined to give the Fourier amplitude $$A_{20}(\omega )=\frac{16\pi }{5\sqrt{3}M}\sqrt{\frac{\pi }{a}}e^{\omega ^2/(4a)}\left\{R_{\mathrm{shell}}^4+\frac{8\pi \rho _{eq}}{5}\frac{R^5}{R_{\mathrm{shell}}}\frac{[\omega _2^2\omega ^2+2i\omega /\tau _2]}{(\omega _2^2\omega ^2)^2+4\omega ^2/\tau _2^2}\right\},$$ (A26) which can then be used in Eq. (A23) to compute the Newtonian quadrupole energy spectrum. The first term in curly brackets in the Fourier amplitude Eq. (A26) is the contribution only of the shell to the total energy radiated. The second term is the contribution only of the star’s oscillations to the energy. This should be contrasted with the relativistic procedure in which no such exact identification of the shell and star contributions can be made.
warning/0002/physics0002024.html	ar5iv	text	# Soliton electro-optic effects in paraelectrics ## Abstract The combination of charge separation induced by the formation of a single photorefractive screening soliton and an applied external bias field in a paraelectric is shown to lead to a family of useful electro-optic guiding patterns and properties. Apart from their inherent interest as peculiar products of nonlinearity, spatial solitons hold the promise of allowing viable optical steering in bulk environments . Photorefractive screening solitons differ from other known manifestations of spatial self-trapping for their peculiar ease of observation and versatility , and recent experiments in photorefractive strontium-barium-niobate (SBN) and potassium-niobate (KNbO<sub>3</sub>) have demonstrated two conceptual applications of their guiding properties. In the first case, a tunable directional coupler was realized making use of two independent slab-solitons ; in the second, self-induced phase-matching was observed to enhance second-harmonic-generation . Although results suggest a means of obtaining all-optical functionality, actual implementation is hampered by the generally slow nonlinear response , that can be ”accelerated” only at the expense of stringent intensity requirements. In contrast, non-dynamic guiding structures have been observed by fixing a screening soliton , or in relation to the observation of spontaneous self-trapping during a structural crystal phase-transition . One possible method of obtaining acceptable dynamics is to make directly use of the electro-optic properties of the ferroelectrics involved, in combination with the internal photorefractive space charge field deposited by the soliton. Since photorefractive charge-activation is wavelength dependent, one can induce charge separation in soliton-like structures at one active wavelength (typically visible), and then read the electro-optic index modulation at a different, nonphotorefractive, wavelength (typically infrared) . For noncentrosymmetric samples (such as the above mentioned crystals) that typically host screening soliton formation, the electro-optic index of refraction modulation is proportional to the static crystal polarization P, and thus to the electric field (linear electro-optic effect). For these, no electro-optic modulation effects are possible: for whatever value of external constant electric field E<sub>ext</sub>, the original soliton supporting guiding pattern remains unchanged. In centrosymmetrics, such as photorefractive potassium-lithium-tantalate-niobate (KLTN), solitons are supported by the quadratic electro-optic effect . In this case, the ”nonlinear” combination of the internal photorefractive field with an external electric field can give rise to new and useful soliton-based electro-optic phenomena, which we here study for the first time. The basic mechanism leading to screening soliton formation is the following: a highly diffracting optical beam ionizes impurities hosted in the lattice of an electro-optic crystal. An externally applied electric field makes these mobile charges drift to less illuminated regions, forming a double layer that renders the resultant electric field in the illuminated region lower. For an appropriate electro-optic sample, this leads to a self-lensing and soliton propagation, when beam diffraction is exactly compensated. For slab solitons, i.e. those self-trapped beams that originate from a beam that linearly diffracts only in one transverse dimension (x), for a given soliton intensity full-width-half-maximum (FWHM) $`\mathrm{\Delta }`$x, a given ratio between the soliton peak intensity and the (generally artificial) background illumination, $`u_0^2=I_{peak}/I_b`$ (intensity ratio), solitons form for a particular value of applied external biasing field $`\overline{E}`$. The soliton-supporting electric field E is expressed by E=(V/L)(1+I(x)/I<sub>b</sub>)<sup>-1</sup>, where V is the external applied voltage, L is the distance between the crystal electrodes (thus $`\overline{E}`$=V/L), and I(x) is the soliton optical intensity confined in the x transverse dimension . This electric field, a result of a complex nonlinear light-matter interaction, is present even when the generating optical field is blocked, and the sample is illuminated with a nonphotorefractively active light. Charge separation is smeared out only by slow recombination, associated with dark conductivity, characterized by considerably long decay times. The nonphotorefractively active illumination, although not leading to any further evolution in the internal charge field, will feel the index inhomogeneity due to the quadratic electro-optic response described by the relation $`\mathrm{\Delta }`$n =-(1/2) n<sup>3</sup> g$`{}_{eff}{}^{}ϵ_{0}^{2}`$($`ϵ_r`$-1)<sup>2</sup>E<sup>2</sup>, where n is the crystal index of refraction, $`g_{eff}`$ is the effective electro-optic coefficient for a given scalar configuration, $`ϵ_0`$ is the vacuum dielectric constant, and $`ϵ_r`$ is the relative dielectric constant. The actual electric field in the crystal is now E=(V/L)(1+I(x)/I<sub>b</sub>)<sup>-1</sup>-(V/L)+E<sub>ext</sub>, where E<sub>ext</sub> (in general $`\overline{E}`$) is the externally applied electric field after the nonlinear processes have occurred (the ”read-out” field). The index pattern induced is $$\mathrm{\Delta }n=\mathrm{\Delta }n_0\left(\frac{1}{1+I(x)/I_b}1+\frac{E_{ext}}{V/L}\right)^2,$$ (1) where $`\mathrm{\Delta }`$n<sub>0</sub> =(1/2) n<sup>3</sup>g$`{}_{eff}{}^{}ϵ_{0}^{2}`$($`ϵ_r`$-1)<sup>2</sup>(V/L)<sup>2</sup>. In Fig.(1) we show two families of induced index patterns associated with two solitons at different saturation levels. In Fig.(1a) a 7$`\mu `$m FWHM soliton at $`\lambda `$=514 nm wavelength ($`\mathrm{\Delta }`$n$`{}_{0}{}^{}`$5.4$`\times 10^4`$, for n=2.45) with an intensity ratio $`u_0^2`$=4, leads to three characteristic pattern regimes: for $`\eta =`$E$`{}_{ext}{}^{}/(V/L)`$1, the soliton supporting potential is formed. For $`\eta `$0, an antiguiding hump appears, whereas for intermediate values of $`\eta `$, a twin-waveguide potential forms. Analogous results can be predicted for a strongly saturated regime shown in Fig.(1b), where a 11$`\mu `$m soliton is formed for $`u_0^2`$22. Experiments are carried out with an apparatus that is well documented in literature . An enlarged TEM<sub>00</sub> Gaussian beam from a CW Argon-ion laser operating at $`\lambda `$ =514nm, is focused be means of an f=150mm cylindrical lens onto the input facet of an $`3.7^{(x)}\times 4.6^{(y)}\times 2.4^{(z)}`$ mm sample of zero-cut paraelectric KLTN, at T=20 C (with a critical temperature T<sub>c</sub>=11 C), giving rise to an approximately one-dimensional x-polarized Gaussian beam of $`\mathrm{\Delta }`$x$``$11 $`\mu `$m (”soliton” beam), and the entire crystal is illuminated with a second, homogeneous beam (”background” beam) from the same laser, polarized along the y axis. Both the focused and the plane-wave beams copropagate along the z-direction. The constant voltage V is applied along the crystal x direction, the crystal itself being doped with Vanadium and Copper impurities, and photorefractively active at the laser wavelength. Guiding patterns can be investigated either by illuminating the crystal with an infrared beam (as mentioned above), or simply by using the same soliton-forming wavelength, but at a lower intensity, since photorefractive temporal dynamics are proportional to beam intensity. Here we use this read-out method, and in what follows all read/write experiments are at $`\lambda `$=514nm, with I<sub>read</sub>/I$`{}_{write}{}^{}`$20. By changing the value of the applied readout voltage, V<sub>ext</sub>, we can explore the optical potential described by Eq.(1), through the variable $`\eta `$. Beam distribution is investigated by imaging the facets of the sample onto a CCD camera by means of a second lens placed after the sample (along the z direction). In Fig.(2) the observation of a single photorefractive screening soliton is shown. The 11$`\mu `$m soliton is observed with an intensity ratio u$`{}_{0}{}^{2}`$ 22 at V<sub>exp</sub>=1.33 kV, annulling linear diffraction to 24 $`\mu `$m . Soliton formation takes approximately 3 min, for an I$`{}_{peak}{}^{}`$1.8 kW/m<sup>2</sup> (I$`{}_{b}{}^{}`$80 W/m<sup>2</sup>), measured directly before the sample, thus meaning that erasure during readout would take, at the very least, about 1 hr (i.e. longer than the duration of any one of our experiments). Had we used an IR read-out beam, decay would be halted indefinetly. Given the sample g<sub>eff</sub>=0.12m<sup>4</sup>C<sup>-2</sup>, $`ϵ_r`$9000, $`\mathrm{\Delta }n_0`$ 6.9$`\times 10^4`$, the expected value for soliton formation would be V$`{}_{th}{}^{}`$1.27 kV. In Fig.(3) we show the same region of the crystal invested by the less intense (but otherwise identical to the soliton generating) ”read” beam at various values of $`\eta `$. For $`\eta `$=1 the output beam is identical to the soliton (apart from the actual intensity). For low values of $`\eta `$ ($`\eta <`$0.4) the index pattern given by Eq.(1) is antiguiding, and the output beam is scattered and split into two diffracting beams (beam ”bursting”, see Fig.(1b)). As $`\eta `$ is increased, the defocusing is weakened and for $`\eta 0.45`$ the sample gives rise to a beam-splitting on the twin-waveguide structure formed by the two-hump potential, as shown in Fig.(1). The distance between the two beams is $`20\mu `$m. As opposed to previous defocusing, in this case light is exciting a guided mode. Next we shift the crystal with respect to the optical beam in the x direction, so as to launch it directly into one of the twin-guides for intermediate values of $`\eta `$. For an $`\eta `$= 0.45, shifting the crystal by 10$`\mu `$m, the beam is guided by the side hump, as shown in Fig.(4b). In this forward guiding condition, we change $`\eta `$ to $`\eta `$=0.8. The potential commutes from a double-hump twin-waveguide to a single guiding pattern (see Fig.1). The optical beam is redirected as shown in Fig.(4c). It is therefore possible to realize, by means of the formation of a single photorefractive centrosymmetric screening soliton, three qualitatively different optical circuits: a single waveguide, a double waveguide beam-splitter, and an antiguiding beam-stopper. If the crystal is shifted so as to launch the guided beam into one of the twin-guides, it is possible to deviate the beam, maintaining its strong confinement, realizing an electro-optic switch. Had we used a longer sample, launching the beam in a twin-waveguide leads to a tunable directional coupler, as shown in Fig.(5). The observed phenomena represent an important step in the achievement of viable soliton based components in two major aspects. The first is that the observed phenomenology occurs with the formation of a single soliton, that is only used to deposit a pattern of charge displacement (a peculiar volume hologram), whereas switching from one regime to the other occurs only through the change of the applied electric field. Thus switching dynamics are only limited by capacity charging times, as all other electro-optic devices. Secondly, whereas screening soliton formation requires a constant applied external field, during read-out, the use of independent electrodes can allow the formation of composite circuitry in cascade, all from a single soliton. The work of E.D. and M.T. was carried out in the framework of an agreement between Fondazione Ugo Bordoni and the Italian Communications Administration. Research carried out by A.J.A. is supported by a grant from the Ministry of Science of the State of Israel.
warning/0002/gr-qc0002062.html	ar5iv	text	# Initial data and spherical dust collapse ## I Introduction It is now well known that, subject to a number of physically reasonable assumptions, final state singularities arise in solutions of Einstein’s equations in a range of settings (see e.g. ). What is not known in general - and this is one of the key outstanding questions in classical general relativity at the moment - is the nature of the resulting singularities and in particular whether and under what conditions they may be naked (NS) or black holes (BH) . A great deal of effort has gone into the study of this question over the recent years. Given the complexity of the full Einstein’s equations, these studies have mainly concentrated on the spherically symmetric collapse and fall into a number of categories. One of the mathematically most developed of these, due to Christodoulou , concentrates on the spherical gravitational collapse of a scalar field and shows that naked singularities do indeed occur , but that in a certain sense they are unstable . Another group of works has concentrated on showing the occurrence of NS solutions in a variety of spherical symmetric space–times with several examples of field–sources, including dust , perfect fluids , imperfect fluids and radiation . In particular it has been shown that in spherical dust collapse, given any initial density profile for the collapsing cloud, the corresponding velocity profile may be chosen such that the collapse may eventually result either in a BH or a NS . In this way, both groups of works demonstrate that the end result of spherical collapse depends upon the nature of the initial data. Our aim here is to make a more thorough study of the outcomes of the inhomogeneous spherical dust collapse and in particular to study the nature of subsets of initial data $``$ corresponding to NS and BH. The structure of the paper is as follows. In section 2 we give a brief description of the spherical dust collapse. In Section 3 we summarise physical and other constraints that need to be satisfied by functions that represent the initial data in these models. Section 4 contains our main results in the form of a number of Lemmas and Propositions and finally Section 5 contains our conclusions. ## II Spherically symmetric dust collapse The inhomogeneous spherically symmetric dust collapse can be represented by the Lemaitre-Tolman-Bondi (LTB) line element which is given by $$ds^2=dt^2+\frac{R_{}^{^{}}{}_{}{}^{2}}{1+E}dr^2+R^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$ (1) where $`r,\theta ,\varphi `$ are the comoving coordinates. The dot and prime denote differentiation with respect to $`t`$ and $`r`$ respectively and $`R=R(r,t)`$ and $`E=E(r)`$ are $`C^2`$ real functions such that $`R(r,t)0`$ and $`E(r)>1`$. The matter–density is given by $$\rho (t,r)=\frac{M^{^{}}}{R^2R^{^{}}}$$ (2) where $`M=M(r)`$ is another $`C^2`$ real function such that $`M(r)>0`$. The evolution equation for the case of $`\dot{R}<0`$ (corresponding to gravitational collapse) takes the form $$\dot{R}=\sqrt{\frac{M}{R}+E}$$ (3) which can be solved to give $$tt_c(r)=\frac{R^{3/2}G(ER/M)}{\sqrt{M}},$$ (4) where $`t_c=t_c(r)`$ is another real $`C^2`$ function that corresponds to the time of arrival of each shell $`r`$ to the central singularity and $`G`$ is a positive function given by $`G(x)=`$ $`{\displaystyle \frac{arcsin\sqrt{x}}{x^{3/2}}}{\displaystyle \frac{\sqrt{1x}}{x}},for1x>0`$ (5) $`G(x)=`$ $`2/3,forx=0`$ (6) $`G(x)=`$ $`{\displaystyle \frac{arcsinh\sqrt{x}}{(x)^{3/2}}}{\displaystyle \frac{\sqrt{1x}}{x}},for0>x>\mathrm{}.`$ (7) Using the coordinate freedom to rescale $$R(0,r)=r,$$ (8) equation (4) gives $$t_c(r)=\frac{r^{3/2}G(p)}{\sqrt{M}},$$ (9) where $`p=r\frac{E}{M}`$. The collapse is then simultaneous for all shells if $`t_c^{}(r)=0`$, which is the case for the homogeneous dust collapse. We note that the metric (1) can be matched at the boundary, say $`r=r_d=const.`$, to the Schwarzschild metric in the exterior region . Thus the scenario here is that of a collapsing compact matter region matched in the exterior to the Schwarzschild geometry. We shall refer to a singularity as naked if there is a family of future directed non–spacelike geodesics which terminate at the singularity in the past. Here we consider the outgoing radial null geodesics which, as can be seen from (1), correspond to the solution of the differential equation $$\frac{dt}{dr}=\frac{R^{}}{\sqrt{1+E}}.$$ (10) One can now rewrite this equation as $$\frac{dR}{du}=\frac{1}{u^{}}\left(R^{}+\dot{R}\frac{dt}{dr}\right)=\left(1\sqrt{\frac{E+\mathrm{\Lambda }/X}{1+E}}\right)H(X,u),$$ (11) where $$H(X,u)=(\eta _u\beta _u)X+(\mathrm{\Theta }_u(\eta _u\frac{3}{2}\beta _u)X^{3/2}G(PX))\sqrt{P+1/X}$$ (12) and $`X={\displaystyle \frac{R}{u}},\eta =r{\displaystyle \frac{M^{}}{M}},\eta _u=\eta {\displaystyle \frac{u}{ru^{}}},\beta =r{\displaystyle \frac{E^{}}{E}},\beta _u=\beta {\displaystyle \frac{u}{ru^{}}},`$ (13) $`P={\displaystyle \frac{uE}{M}},p=r{\displaystyle \frac{E}{M}},\mathrm{\Theta }_u=\mathrm{\Theta }{\displaystyle \frac{\sqrt{r}}{\sqrt{u}u^{}}},\mathrm{\Lambda }={\displaystyle \frac{M}{u}},`$ (14) $`\mathrm{\Theta }={\displaystyle \frac{\sqrt{M}}{\sqrt{r}}}t_c^{}(r)={\displaystyle \frac{1+\beta \eta }{\sqrt{1p}}}+(\eta {\displaystyle \frac{3}{2}}\beta )G(p),`$ (15) with the positive real function $`u=u(r)`$ being monotonically increasing and such that $`u(0)=0`$. Later we will specify $`u(r)=r^\alpha ,\alpha >0`$. In the cases where $`E(r)=0`$ we will take $`\beta (r)=p(r)=P(r)=0`$ by convention. A subscript $`\mathrm{`}0\mathrm{`}`$ will denote the limit of the associated functions as $`r0`$ (respectively $`u0`$). We should emphasise that the existence of such limits cannot be assumed a priori and need to be ensured for the given set of initial data under consideration, as we shall do in the following. It can be shown from that a sufficient condition for the existence of a naked singularity in spherical symmetric dust collapse is that the following algebraic equation in $`X_0`$ $$\left(1\sqrt{\frac{E_0+\mathrm{\Lambda }_0/X_0}{1+E_0}}\right)H(X_0,0)X_0=0,$$ (16) possesses a real positive root. These roots give the possible values of tangents for the outgoing geodesics at the singularity such that the associated integral curves terminate at the singularity in the past. We note that given the limiting nature of (16), the forms of $`E`$ and $`M`$ in a neighbourhood of $`r=0`$ will play an important role in determining the possible solutions to this equation. An interesting outcome of all the studies of the spherical dust collapse in the literature is that for the initial functions $`E`$ and $`M`$ chosen so far, the most general form of the equation (16) becomes a polynomial of degree not greater than four . As we shall see in the next section, this feature is very important in constraining the way the subsets of initial data corresponding to NS and BH are distributed in $``$. As a result it is important to ask whether a quartic is the most general form equation (16) can take. Another important outcome for these studies is that the occurrence of BH or NS as final outcomes of collapse depend on the choice of initial data. The question arises as to the nature of the subsets of the initial data that lead to each of these outcomes and how robust are these outcomes with respect to perturbations in the initial data. Before discussing these issues in Section 4, we briefly consider, in the next Section, some constraints that are to be satisfied by the functions $`E`$, $`M`$ (and $`\mathrm{\Theta }`$) on physical grounds. ## III Physical constraints In the case of spherical dust collapse the initial data are given in terms of two functions; namely the mass function for the dust cloud, $`M=M(r)`$, and the energy function $`E=E(r)`$, which is related to the initial radial velocity $`V_I(r)=\dot{R}(0,r)`$ of the shells by $$E(r)=V_I^2(r)\frac{M(r)}{r}.$$ (17) Here we briefly summarise the constraints that these functions, as well as $`\mathrm{\Theta }`$, need to satisfy in order to be physically acceptable. To begin with, to ensure that the curvature is initially well behaved at the regular centre of the matter distribution, we demand the condition that the quadratic curvature scalar (the so called Kretschmann scalar) given by $$K=R^{ijkl}R_{ijkl}=\frac{12M^2}{R^4R^2}\frac{32MM^{}}{R^5R^{}}+\frac{48M^2}{R^6}$$ (18) remains bounded. We should note here that in this case the Ricci scalar is the only term that remains relevant near origin, if the origin is regular. We have employed the Kretschmann scalar because in addition to encoding all the information we need from the Riemann invariants, it enables us to ensure that the centre is regular. To ensure this, as well as the finiteness of the density distribution at the initial epoch, $`M`$ in the neighbourhood of the origin needs, in general, to be of the form $$M(r)=O(r^a),a3.$$ (19) If we require $`\rho (0,0)0`$ then we need in general $$M(r)=r^3g(r)+O(r^m),m>3,$$ (20) where $`g`$ is a $`C^2`$ function of $`r`$ that remains finite as $`r0`$. The condition for the absence of shell-crossing is given by $$\mathrm{\Theta }(r)0.$$ (21) We also note that the apparent horizon is given by $`R(r,t)=M(r)`$ and in order to ensure that the initial hypersurface does not contain any trapped surfaces we shall require the condition $`M(r)/R(r,0)<1`$. By a regular initial data we will mean initial data that satisfies the physical conditions given above. For simplicity in the following we shall, unless otherwise stated, assume that, in a neighbourhood of $`r=0`$, the functions $`E=E(r)`$ and $`M=M(r)`$ can be expressed as $$M(r)=\underset{i=3}{\overset{\mathrm{}}{}}M_ir^i,E(r)=\underset{j=0}{\overset{\mathrm{}}{}}E_jr^j,M_i,E_j\mathrm{}.$$ (22) We shall refer to the first and second non-vanishing powers of $`M`$ and $`E`$ by $`i_0`$, $`j_0`$ and $`i_1`$, $`j_1`$ respectively. From here on we shall take $`u(r)=r^\alpha ,\alpha >0`$ . Letting $`R=X_0r^\alpha `$ in the neighbourhood of the singularity then $`Kr^{2(i_03\alpha )}`$ which implies that $`K`$ diverges as $`r0`$ only along geodesics with $`\alpha >i_0/3`$ which for the case $`i_03`$ gives $`\alpha >1`$. We shall use this condition in the next section in the proofs of the Lemmas. Finally we recall that for homogeneous initial data with $`E(r)0`$ we have $`M(r)^2=kE(r)^3`$, $`k`$ constant, which for initial data of the form (22) gives $`t_c(r)=const`$, which in turn implies a simultaneous collapse . The next definition will be useful in what follows: Definition: A LTB initial data set is said to be centrally homogeneous if, in a neighbourhood of the origin, given $`E(r)0`$, $`M(r)^2E(r)^3`$ or given $`E(r)=0`$, $`M(r)R(0,r)^3`$. In particular, a LTB initial data set with $`E(r)0`$ is centrally homogeneous if near the origin $$M(r)=\underset{i=3}{\overset{\mathrm{}}{}}M_ir^i,E(r)=\underset{j=2}{\overset{\mathrm{}}{}}E_jr^j;M_3>0,E_20,$$ (23) In the following we shall refer to perturbations which break the central homogeneity of the initial data as centrally inhomogeneous perturbations. ## IV Initial data and spherical dust collapse In this section we study the final outcomes of the spherical dust collapse as a function of the choice of initial data, by employing a family of outgoing null geodesics from the origin. In particular, we study the subsets of regular initial data which give rise to NS solutions in (16). As was shown in , finiteness of $`\mathrm{\Lambda }_0`$ is a necessary condition for the existence of NS solutions. To ensure this we require that $`i_0\alpha `$. We now note that allowing $`\mathrm{\Theta }_u,P`$ or $`p`$ to diverge as $`r0`$ would make equation (16) singular at $`r=0`$, for any regular initial data with $`i_0\alpha `$ and $`\alpha >1`$. Therefore in the following we shall assume that $`\mathrm{\Theta }_{u_0},P_0`$ and $`p_0`$ are finite. We begin with a simple (well known) result which demonstrates that in the homogeneous setting the set of initial data which lead to BH is full. Lemma 1: Consider the spherically symmetric dust collapse with initial data given by the functions $`E`$ and $`M`$ in the form (22) and assumed to be homogeneous. Then equation (16) has no NS solutions. Proof: The homogeneity condition for the functions (22) implies $`2i_0=3j_0`$, in the case $`E(r)0`$. On the other hand, the requirement of finiteness of $`p_0`$ implies $`1+j_0i_00`$ which results in $`i_03`$. From condition (18) we obtain $`i_0=3`$ and the homogeneity implies $`j_0=2`$. This gives $`\mathrm{\Theta }(r)=0`$ and the equation (16) has then the solutions $$X_0=0X_0=\frac{\sqrt{\mathrm{\Lambda }_0}}{1\alpha },$$ (24) which for $`\alpha >1`$ give $`X_00`$. The same conclusion holds for the homogeneous marginally bound case where $`E(r)=0`$ and $`i_0=3`$ $`\mathrm{}`$. We note that the existence of real positive roots to equation (8) basically characterizes the formation and time development of the apparent horizon as the collapse develops . To understand further the structure of the initial data in the inhomogeneous dust collapse, it is important, as a first step, to establish how general are the conditions for which (16) is a polynomial of degree $`4`$. The following Proposition makes precise the forms that equation (16) may take in order for the spherical symmetric dust collapse to result in a NS solution. Proposition 1: Consider the spherically symmetric dust collapse with initial data given by the functions $`E`$ and $`M`$ in the form (22) and assumed to be regular. Then the only non-divergent forms of the equation (16) are polynomials of degree not greater than four. Proof: For equation (16) to remain finite in the limit $`r0`$, the limiting quantities $`P_0,E_0,\mathrm{\Lambda }_0`$ and $`\mathrm{\Theta }_{u_0}`$ need to remain finite. Let $`x_0=\sqrt{X_0}`$ and recall that $`u=r^\alpha `$ and that $`j_0`$ and $`i_0`$ are the orders of the first non-vanishing coefficients of $`E`$ and $`M`$ respectively. We proceed by considering all different combinations of these limiting quantities in turn. (i) Let $`P_0=E_0=0`$, $`\mathrm{\Theta }_{u_0}0`$ and $`\mathrm{\Lambda }_00`$. Then equation (16) becomes $$\left(\frac{1}{3}\eta _{u_0}1\right)x_0^4\frac{1}{3}\sqrt{\mathrm{\Lambda }_0}\eta _{u_0}x_0^3+\mathrm{\Theta }_{u_0}(x_0\sqrt{\mathrm{\Lambda }_0})=0.$$ (25) The exponent $`\alpha `$ can then be chosen such that $`\mathrm{\Theta }_u,\eta _u`$ and $`\mathrm{\Lambda }`$ remain finite at $`r=0`$. (ii) For $`P_0=E_0=0`$, $`\mathrm{\Theta }_{u_0}0`$ and $`\mathrm{\Lambda }_0=0`$, equation (25) gives $$x_0=0\left(x_0^3=\frac{\mathrm{\Theta }_{u_0}}{1\eta _{u_0}/3}\eta _{u_0}3\right).$$ (26) (iii) For $`P_0=E_0=0`$, $`\mathrm{\Theta }_{u_0}=0`$ and $`\mathrm{\Lambda }_0=0`$, equation (25) reduces to $$\left(\frac{1}{3}\eta _{u_0}1\right)x_0^4=0,$$ (27) which has the solutions $`x_0=0`$ or $`\frac{1}{3}\eta _{u_0}1=0`$. (iv) For $`P_0=E_0=0`$, $`\mathrm{\Theta }_{u_0}=0`$ and $`\mathrm{\Lambda }_00`$, equation (25) gives $$x_0=0\left(x_0=\frac{\sqrt{\mathrm{\Lambda }_0}\eta _{u_0}}{\eta _{u_0}3}\eta _{u_0}3\right).$$ (28) (v) Let $`P_00`$, $`E_0=0`$, $`\mathrm{\Theta }_{u_0}0`$ for any finite $`\mathrm{\Lambda }_0`$. We first consider $`p_00`$ (i.e., $`1+j_0i_0=0`$) which, since $`P_00`$ (i.e., $`\alpha +j_0i_0=0`$), implies $`\alpha =1`$ and $`1+\beta _0\eta _0=0`$. Requiring $`i_03`$ gives $`\eta _0\frac{3}{2}\beta _00`$ which either contradicts the assumption $`P_00`$ or violates the shell-crossing condition. Suppose now $`p_0=0`$ which implies $`\alpha <1`$ and in turn $`\mathrm{\Theta }_{u_0}=0`$, contradicting the assumption $`\mathrm{\Theta }_{u_0}0`$ and resulting in equation (16) to become singular. (vi) Let $`E_00`$ for any finite values of $`P_0,\mathrm{\Theta }_{u_0}`$ and $`\mathrm{\Lambda }_0`$. Then the only way to make $`E_00`$ (i.e., $`j_0=0`$) and $`p_0`$ finite (i.e., $`1+j_0i_00`$) is to assume $`i_01`$, which makes (18) diverge and thus results in equation (16) to become singular. (vii) Let $`E_0=\mathrm{\Theta }_{u_0}=\mathrm{\Lambda }_0=0`$ and $`P_00`$. Then from (v) above, with $`p_00`$, we obtain $`1+\beta _0\eta _0=0`$ and $`\alpha =1`$. Now since $`\mathrm{\Theta }_{u_0}=0`$, we necessarily have $`\eta _0\frac{3}{2}\beta _0=0`$ which gives $`i_0=3`$ and $`j_0=2`$ and equation (16) becomes identically satisfied for all $`x_0`$. The other possibility is $`p_0=0`$ which implies $`\alpha <1`$ which in general results in $$x_0=0\left[(\eta _{u_0}\beta _{u_0}1)x_0(\eta _{u_0}\frac{3}{2}\beta _{u_0})G(P_0x_0^2)\sqrt{P_0+\frac{1}{x_0^2}}=0\right].$$ (29) Now since $`P_00`$ implies $`\alpha +j_0i_0=0`$, the first term in the square bracket in the above equation vanishes. On the other hand, the second term would vanish if $`x_0=0`$ or $`\eta _0\frac{3}{2}\beta _0=0`$ (i.e. $`i_0\frac{3}{2}j_0=0`$) which, since $`p_0=0`$, implies $`i_0<3`$, hence making $`K`$ diverge initially. The same term can also vanish if $$x_0^2=\frac{1}{P_0},$$ (30) for any $`P_0<0`$. $`\mathrm{}`$ This Proposition shows that in the case of spherical dust collapse, equation (16) can only take a restrict number of forms with the most general being a polynomial of degree 4 (corresponding to $`E_0=0`$, $`\mathrm{\Theta }_00`$, $`P_0=0`$ and $`\mathrm{\Lambda }_00`$). Now an important question is the way the outcome of the spherically symmetric dust collapse relates to the central homogeneity or otherwise of the initial data. The following Proposition makes this precise. Proposition 2: Consider the spherically symmetric dust collapse with initial data given by the functions $`E`$ and $`M`$ in the form (22) and assumed to be regular. Let $`p_00`$. For the occurrence of NS solutions in (16) the initial data set must be centrally homogeneous. Proof: We shall show that if $`E`$ and $`M`$ do not satisfy (23)<sup></sup><sup></sup>*In fact all that is required is that in the case of $`p_00`$ the power expansion of $`E`$ and $`M`$ have lower powers of $`r`$ given by 2 and 3 respectively. for $`p_00`$, then equations (25)–(29) have no real positive roots. We shall proceed by considering the following cases: (i) Let $`p_0=0`$, $`\alpha <1`$. In this case $`\mathrm{\Lambda }_0=0`$, $`\mathrm{\Theta }_{u_0}=0`$ and a non–zero solution of (27) requires $`\eta _0<3`$ (ie, $`i_0<3`$) which makes $`K`$ divergent. Case (vi) of Proposition 1 gives no NS solutions either. In the case of (vii), however, we may find NS solutions for $`P_0<0`$ and $`3>j_0>2`$ and $`i_0=3`$. (ii) Let $`p_0=0`$, $`\alpha =1`$. This implies $`P_0=0`$ and $`\mathrm{\Lambda }_0=0`$ which, by requiring $`\mathrm{\Theta }_0>0`$ in (26), means that we would need to have $`\eta _{u_0}=\eta _0<3`$, which does not satisfy the condition on $`K`$. However, NS may arise in the case (iii) of Proposition 1 with $`i_0=3`$ and $`j_0>2`$. (iii) Let $`p_0=0`$, $`\alpha >1`$. In this case $$\mathrm{\Theta }_u=(1\eta /3)r^{\frac{3}{2}(1\alpha )},$$ (31) which implies that in order to make $`\mathrm{\Theta }_{u_0}`$ finite we require $`i_0=3`$. But since $`p_0=0`$ (ie, $`1+j_0i_0>0`$) then in addition to $`P_0=0`$ we also require $`j_0>2`$, which demonstrates that in this case the presence of NS does not require centrally homogeneous initial data. (iv) Let $`p_00`$, $`\alpha <1`$. From $`p_00`$ we have $`1+j_0i_0=0`$. In the case of $`P_0=0`$ and $`\alpha <1`$ we must have $`1+j_0i_0>0`$ which is a contradiction. On the other hand, if $`P_00`$ and $`p_00`$, then $`\alpha =1`$ which is contrary to our assumptions and hence there are no NS solutions in this case. (v) Let $`p_00`$, $`\alpha =1`$. These conditions imply $`P_00`$ and therefore from the case (vii) of Proposition 1 we may have NS solutions only if $`i_0=3`$ and $`j_0=2`$. (vi) Let $`p_00`$, $`\alpha >1`$. As in case (iii) above, $`\mathrm{\Theta }`$ must vanish as $`r0`$. If we assume $`K0`$ initially we need $`i_0=3`$, but since $`p_00`$ (i.e., $`1+j_0i_0=0`$), then $`j_0=2`$. On the other hand, if we take $`K=0`$ initially then from $`1+j_0i_0=0`$ we obtain $`1+\beta _0\eta _0=0`$ which in order to make $`\mathrm{\Theta }_{u_0}`$ finite necessitates $`\eta _0\frac{3}{2}\beta _0=0`$, which implies $`i_0\frac{3}{2}j_0=0`$ and hence $`j_0=2`$ and $`i_0=3`$. The above considerations show that for $`E(r)0`$ and $`p_00`$ we only obtain a NS solution if the lowest order of the powers of $`r`$ in $`M`$ and $`E`$ are $`i_0=3`$ and $`j_0=2`$ respectively. In the case of $`E(r)=0`$ we have $`\beta (r)=p(r)=P(r)=0`$ and $`\mathrm{\Theta }(r)=1\frac{1}{3}\eta `$. If $`\eta _0>3`$, then $`\mathrm{\Theta }<0`$ which violates the shell-crossing condition. If $`\eta _0<3`$, then $`i_0<3`$ which makes $`K`$ divergent. So, the only possibility is to have $`\eta _0=3`$ (ie, $`i_0=3`$) $`\mathrm{}`$. We note that this result essentially follows from the non-shell-crossing conditions and the initial regularity of $`K`$. The above Proposition shows that unless $`p_0=0`$, inhomogeneous neighbourhoods of the centrally homogeneous initial data do not in general result in NS. Now $`p_0=0`$ implies that as we approach the origin the solution describing the matter distribution tends to a marginally bound solution (given by (4) and (5)), which is rather special. This then indicates that in general in order to have NS only small departures from homogeneity can be allowed near the origin. In this sense, NS in spherical dust collapse may be said to be mathematically unstable. It is, however, important to note that even though NS may be unstable with respect to general perturbations (in this case centrally inhomogeneous perturbations), they may nevertheless stabilise if restricted classes of perturbations are allowed. In particular this seems to be the case if only physically motivated perturbations are allowed. To see this more precisely, recall that given the forms (23) for the functions $`M`$ and $`E`$, then the only way to break the central homogeneity of (23) is by letting $`M_3=0`$ or $`E_2=0`$. It turns out, however, that $`M_3=0`$ would result in a density profile given by $`\rho (r)=M^{}(r)/r^2`$ that tends to zero at the center $`r=0`$; a result contrary to the physical expectation of the density being higher at the centreIn fact the condition for $`\rho ^{}(0,r)>0`$ is $`rM^{\prime \prime }(r)<2M^{}(r)`$.. Similarly, letting $`E_2=0`$, implies $`M_3=0`$ for $`p_00`$ which would again result in $`\rho `$ to become zero at the centre. This gives an important demonstration of the fact that instability deduced with respect to general perturbations can become stabilised once the class of perturbations are restricted (in this case to physically motivated ones). Having demonstrated the relation between NS and central homogeneity, it is of importance to be able to determine the subset of the initial data that result in NS solutions. The following Lemma makes this precise. Lemma 2: Consider the spherically symmetric dust collapse with centrally homogeneous initial data given by the functions $`E`$ and $`M`$ in the form (23) and assumed to be regular. Then the set of initial data corresponding to NS as final states possess open intervals in $`M_i`$ and $`E_j`$. Proof: Considering the case of $`\mathrm{\Lambda }_0=0`$ with $`\alpha =5/3`$ gives $$\mathrm{\Theta }_{u_0}=3\frac{E_3}{E_2}\left(\frac{1}{\sqrt{1p_0}}\frac{3}{2}G(p_0)\right)+4\frac{M_4}{M_3}\left(G(p_0)\frac{1}{\sqrt{1p_0}}\right).$$ (32) From the case (ii) of Proposition 1 we find that the occurrence of a NS solution necessitates $`\mathrm{\Theta }_{u_0}>0`$. We also know that $`G(p)\frac{1}{\sqrt{1p}}<0`$ for $`1p>\mathrm{}`$ (see (5)) and that $`\frac{1}{\sqrt{1p}}\frac{3}{2}G(p)>0`$ for $`1p>0`$ and $`\frac{1}{\sqrt{1p}}\frac{3}{2}G(p)<0`$ for $`0>p>\mathrm{}`$. Therefore the condition $`\mathrm{\Theta }_{u_0}>0`$ is always satisfied for the intervals $$M_3]1,+\mathrm{}[,M_4]\mathrm{},0[,E_2]0,1[,E_3]\mathrm{},0[,$$ (33) or $$M_3]0,+\mathrm{}[,M_4]\mathrm{},0[,E_2]\mathrm{},0[,E_3]\mathrm{},0[,$$ (34) with all other coefficients $`M_i`$ and $`E_j`$ (with $`j2`$) in (23) being real and arbitrary. In this way we have found, for both $`E_2<0`$ and $`E_2>0`$, open intervals in all coefficients $`M_i,i3`$ and $`E_j,j2`$ such that the final state of collapse is a NS. $`\mathrm{}`$ As an example of this result, Figure 1 depicts the outcomes of the dust collapse as a function of the two parameters $`M_3`$ and $`E_5`$, (while for the sake clarity the other coefficients $`M_i`$ and $`E_j`$ in the initial data are kept fixed). As can be seen, there exists open neighbourhoods of the initial data (in the $`M_3E_5`$ plane) resulting in each outcome. The definition of central homogeneity given in Section 3 constraints the functions $`E`$ and $`M`$ by imposing lower bounds on the degrees of the polynomials in the expansions of these functions around the origin. It turns out that for the presence of NS, central homogeneity also places upper bounds on the second non-vanishing coefficients of $`E`$ and $`M`$. This is due to the fact that the function $`\mathrm{\Theta }`$ vanishes at the lowest order for $`r=0`$, so in order to ensure the finiteness of $`\mathrm{\Theta }_{u_0}`$ we will have to analyse the second non-vanishing terms in $`E`$ and $`M`$ and fine tune them in line with the choice of $`\alpha `$ . To begin with, we note that for a centrally homogeneous initial data set we have $$\mathrm{\Theta }_{u_0}=\underset{r0}{lim}\left[j_1\frac{E_{j_1}}{E_2}r^{j_12}\left(\frac{1}{\sqrt{1p}}\frac{3}{2}G(p)\right)+i_1\frac{M_{i_1}}{M_3}r^{i_13}\left(G(p)\frac{1}{\sqrt{1p}}\right)\right]r^{\frac{3}{2}(1\alpha )}.$$ (35) The necessary and sufficient condition for $`\mathrm{\Theta }_{u_0}`$ to be positive and finite is given by $`\left(j_1={\displaystyle \frac{3}{2}}(\alpha 1)+2i_1>{\displaystyle \frac{3}{2}}(\alpha 1)+3sgn(E_{j_1})=sgn(E_2)\right)`$ (36) $`\left(j_1>{\displaystyle \frac{3}{2}}(\alpha 1)+2i_1={\displaystyle \frac{3}{2}}(\alpha 1)+3M_{i_1}<0\right)`$ (37) $`\left(j_1={\displaystyle \frac{3}{2}}(\alpha 1)+2i_1={\displaystyle \frac{3}{2}}(\alpha 1)+3j_1{\displaystyle \frac{E_{j_1}}{E_2}}\left({\displaystyle \frac{1}{\sqrt{1p_0}}}{\displaystyle \frac{3}{2}}G(p_0)\right)+i_1{\displaystyle \frac{M_{i_1}}{M_3}}\left(G(p_0){\displaystyle \frac{1}{\sqrt{1p_0}}}\right)>0\right),`$ (38) in the case $`E(r)0`$ and $$i_1=\frac{3}{2}(\alpha 1)+3M_{i_1}<0$$ (39) in the case $`E(r)=0`$, where $`sgn`$ denotes the sign function. Now recall that in the case $`i_0=3`$ in order to ensure that the geodesics originating at the singular point $`r=0`$ can come out and, at the same time, to have $`K`$ divergent as $`r0`$, we need to consider outgoing radial null geodesics such that $`\alpha >1`$ (see Section 3). We note also that the cases in Proposition 1 where real positive solutions with $`\alpha >1`$ may occur in (16) correspond to $`\mathrm{\Theta }_{u_0}>0`$, i.e. cases (i) and (ii). Cases (iii), (iv) and (vii) require $`\alpha 1`$ and will not be considered in what follows. Given these definitions and notations on $`i_1`$ and $`j_1`$, necessary conditions for the existence of a NS solution in the case of centrally homogeneous initial data are made precise through the following Lemma. Lemma 3: Consider the spherically symmetric dust collapse with centrally homogeneous initial data given by the functions $`E`$ and $`M`$ in the form (23) and assumed to be regular. Then in order for (16) to have NS solutions it is necessary to have $`\left(j_1(2,5]i_1>3sgn(E_{j_1})=sgn(E_2)\right)\left(j_1>2i_1(3,6]M_{i_1}<0\right)`$ (40) $`\left(j_1(2,5]i_1(3,6]j_1{\displaystyle \frac{E_{j_1}}{E_2}}\left({\displaystyle \frac{1}{\sqrt{1p_0}}}{\displaystyle \frac{3}{2}}G(p_0)\right)+i_1{\displaystyle \frac{M_{i_1}}{M_3}}\left(G(p_0){\displaystyle \frac{1}{\sqrt{1p_0}}}\right)>0\right),`$ (41) if $`E(r)0`$ and $$i_1(3,6]M_{i_1}<0,$$ (42) if $`E(r)=0`$. Proof: In order to have a finite value of $`\mathrm{\Lambda }_0`$ for centrally homogeneous initial data we need $`\alpha 3`$. Since for $`\alpha >1`$ we need $`\mathrm{\Theta }_{u_0}>0`$ to have NS solutions to (16) then the necessary conditions follow directly from (36) and (39). $`\mathrm{}`$ Finally, for the case of the spherical dust collapse we give necessary and sufficient conditions on the functions $`M`$ and $`E`$, given by (22) near the origin, for the existence of NS solution to (16) in the case $`p_00`$. The results that follow generalise previous results given in where sufficient conditions for the occurrence of NS were obtained considering the functional forms (23). Lemma 4: Consider the spherically symmetric dust collapse with initial data given by the functions $`E`$ and $`M`$ in the form (22) and assumed to be regular. Let $`p_00`$ and $`\mathrm{\Lambda }_0=0`$. Then there are NS solutions to (16) if and only if $`\mathrm{\Theta }_{u_0}>0`$ and the initial data is centrally homogeneous. Proof: The necessary condition of central homogeneity was already established in Proposition 2. Now, the case (ii) of Proposition 1 demonstrates that for centrally homogeneous initial data and $`\alpha >1`$ equation (26) has positive real solutions if and only if $`\mathrm{\Theta }_{u_0}>0`$. $`\mathrm{}`$ Lemma 5: Consider the spherically symmetric dust collapse with initial data given by the functions $`E`$ and $`M`$ in the form (22) and assumed to be regular. Let $`p_00`$ and $`\mathrm{\Lambda }_00`$. Then there are NS solutions to (16) if and only if $$\mathrm{\Theta }_{u_0}]0,M_3^{\frac{3}{2}}\left(\frac{13}{3}\frac{5\sqrt{3}}{2}\right)[]M_3^{\frac{3}{2}}\left(\frac{13}{3}+\frac{5\sqrt{3}}{2}\right),+\mathrm{}[$$ (43) and the initial data is centrally homogeneous. Proof: From Proposition 2 we know that the initial data must be centrally homogeneous. Now, the case (i) of the Proposition 1 demonstrates that the existence of a NS solution depends on the solutions of a quartic equation. For centrally homogeneous initial data these solutions exist and are positive and real if and only if $`\mathrm{\Theta }_{u_0}]0,\mathrm{\Lambda }_0^{3/2}(13/35\sqrt{3}/2)[]\mathrm{\Lambda }_0^{3/2}(13/3+5\sqrt{3}/2),+\mathrm{}[`$, with $`\mathrm{\Lambda }_0=M_3`$. $`\mathrm{}`$ ## V Conclusions We have studied the final outcomes of the inhomogeneous spherical dust collapse as a function of initial data, given in the form (22) near the origin, with the help of families of radial null geodesics. We have found all the possible cases where NS solutions can arise in such collapse and have demonstrated that assuming regular initial data then the most general form of the equation (12) is a quartic. We have defined the notion of central homogeneity and have proved that for the occurrence of naked singularities the initial data must in general be centrally homogeneous. Mathematically this result indicates that in general one would expect NS to be unstable to centrally inhomogeneous perturbations. It turns out, however, that such perturbations are not physically reasonable, in the sense that they would require the density to tend to zero as $`r0`$, contrary to the physical expectation that the density increases as we approach the centre. In this way our results show that NS in the setting considered here remain robust when physically motivated perturbations are allowed. They also provide us with an example of the fact that stability and instability of a particular phenomenon will crucially depend on the class of perturbations allowed. In this case, the occurrence of NS, though mathematically unstable to general centrally inhomogeneous perturbations, can become stabilised once only physically motivated perturbations (in the sense made precise above) are allowed. This is a potentially important point to bear in mind in the general debates regarding the generic presence of naked singularities in gravitational collapse. It may also be noted that as far as the existence of naked singularities and black holes as end product of collapse is concerned, rather general results are available, with $`M`$ and $`E`$ being just $`C^2`$, without any further restrictions (see e.g. ). However, these results deal with only the existence part, and give no insight on possible distribution of these outcomes. We believe this latter has been achieved here for the first time in an explicit manner for a wide class of physically reasonable initial data. Finally we should add that most of our results are readily generalisable to initial data characterised by the functional form (19), see , and that generalizations to more generic settings than spherical dust are currently being investigated and will be the subject of a future publication. Acknowledgments We thank Malcolm MacCallum and Brien Nolan for reading the manuscript. FCM thanks Centro de Matemática, Universidade do Minho, for support and FCT (Portugal) for grant PRAXIS XXI BD/16012/98. RT benefited from PPARC UK Grant No. L39094. PSJ thanks the Astronomy Unit, QMW, for hospitality and CERN for grant CERN/S/FAE/1172/97.
warning/0002/hep-ph0002071.html	ar5iv	text	# 1 Introduction ## 1 Introduction The Compton amplitude for the scattering of a virtual photon off a hadron $$\gamma _1^{}+p_1\gamma _2^{()^{}}+p_2$$ (1.1) provides one of the basic tools to understand the short–distance behavior of the nucleon and to test Quantum Chromodynamics (QCD) at large space–like virtualities. In this kinematic regime the Compton amplitude is dominated by the singularities of the light–cone and can be described by the light–cone expansion in terms of the contributions of different twist . Many investigations were devoted to the process of deep–inelastic forward scattering in the past, for which the leading, next–to–leading order, and partly the 3–loop order corrections were calculated . There the absorptive part of the Compton amplitude, the hadronic tensor $`W_{\mu \nu }`$, is described by the four independent structure functions $`F_1(x,Q^2),F_2(x,Q^2),g_1(x,Q^2)`$ and $`g_2(x,Q^2)`$ .<sup>1</sup><sup>1</sup>1In the case of general electroweak currents eight structure functions contribute in the massless limit, cf. . A more general view is obtained studying the non–forward process (1.1) without restricting it to the absorptive contributions only. In this more general kinematics a systematic description of a whole class of physical processes in the deep–inelastic region can be given and common properties of these processes can be derived. Moreover, relations obtained for forward scattering become more transparent under this generalized point of view. It is the aim of this paper to calculate the Compton amplitude for non–forward scattering in this kinematic range in lowest order in QCD. General spin states are considered which allows us to study besides the unpolarized contributions also those due to target polarization. The present analysis is devoted to the study of the twist–2 terms. The twist–decomposition of the Compton amplitude in the generalized Bjorken region is necessary since the scaling violations of the respective amplitude functions differ due to the set of the contributing operators and their anomalous dimensions . The different contributions to the Compton amplitude are described by non–perturbative amplitude functions, which can be measured experimentally in the various non–forward processes. Analogously to the case of forward scattering some of the amplitude functions are experimentally easier accessible than others<sup>2</sup><sup>2</sup>2 Both the measurement of the unpolarized longitudinal structure function $`F_L(x,Q^2)=F_2(x,Q^2)2xF_1(x,Q^2)`$ and the polarized structure function $`g_2(x,Q^2)`$ are experimentally much more demanding than the measurement of the structure functions $`F_2(x,Q^2)`$ and $`g_1(x,Q^2)`$.. Therefore it is important to know the relations between the different amplitude functions due to which predictions can be made for experiment. One of the main objectives of this investigation is to derive these relations for the twist–2 parts of the amplitude functions. The paper is organized as follows. In Section 2 the twist–2 contributions to the Compton amplitude are calculated. We apply the non–local light–cone expansion to express the operator products. The non–forward expectation values are calculated for general spin states and the corresponding Lorentz–structure is derived. A helicity basis of both (virtual) photons is constructed (Section 3) which is subsequently used to determine the twist–2 contributions to the operator matrix elements. The light–cone expansion does not a priori lead to an explicit gauge–invariant representation of the different contributions to the Compton amplitude. Since the Compton amplitude itself is gauge invariant, the expansion has to be formulated in such a way that the respective contributions under discussion form gauge–invariant sub–sets. In Section 4 we show that this is the case for the twist–2 contributions for non–forward scattering. The helicity projections of the Compton amplitude are discussed in detail in Section 5. In Section 6 we derive the relations between the twist–2 contributions to the different amplitude functions, which are generalizations of the relations known in the case of forward scattering . Section 7 contains the conclusions. ## 2 The Compton Amplitude The Compton amplitude for the general case of non–forward scattering is given by $$T_{\mu \nu }(p_+,p_{},q)=id^4xe^{iqx}p_2,S_2\|T(J_\mu (x/2)J_\nu (x/2))\|p_1,S_1.$$ (2.1) Here, $`p_+`$ $`=`$ $`p_2+p_1,p_{}=p_2p_1=q_1q_2,`$ (2.2) $`q`$ $`=`$ $`\frac{1}{2}\left(q_1+q_2\right),p_1+q_1=p_2+q_2,`$ (2.3) where $`q_1(q_2)`$ and $`p_1(p_2)`$ denote the four–momenta of the incoming (outgoing) photon and hadron, respectively, and $`S_1,S_2`$ are the spins of the initial– and final–state hadron. Representations of the Lorentz–structure of the Compton amplitude were given in Refs. . In the following we will consider the Compton amplitude in the generalized Bjorken region only which is defined by the conditions $`\nu =qp_+\mathrm{},q^2\mathrm{},`$ (2.4) keeping the variables $`\xi ={\displaystyle \frac{q^2}{qp_+}},\eta ={\displaystyle \frac{qp_{}}{qp_+}}={\displaystyle \frac{q_1^2q_2^2}{2\nu }}`$ (2.5) fixed. The subsequent analysis will be performed demanding furthermore that the vector $`q_1`$ is spacelike and the vector $`q_2`$ can be space–, light–, or time–like, where $`q_1`$ $`=`$ $`q+\frac{1}{2}p_{}`$ (2.6) $`q_2`$ $`=`$ $`q\frac{1}{2}p_{}.`$ (2.7) In the generalized Bjorken region the Compton amplitude is dominated by the light–cone singularities. It is therefore possible to apply the light–cone expansion for its representation. In the following we use the non–local light-cone expansion , which is a summed–up form of the local light–cone expansion with respect to the spin indices. The respective expressions in the local light–cone expansion can be obtained from the former one by a Taylor expansion. In this way more compact representations can be obtained, cf. section 6. In this paper we study the twist–2 contributions to the non–forward Compton amplitude. They are obtained from the expectation values of the non-local twist–2 light cone operators, cf. . Here the notion of twist is used in its original form as canonical dimension - spin for the local operators which are summed up to the non–local operators. In calculating the non–forward operator expectation values it turns out that the twist–decomposition performed for the operators is not necessarily complete in the case of the expectation values unlike the case for forward scattering. The emergence of new hadronic mass scales, such as $`p_+.p_{}`$ or also of off–shell terms $`p_+^2,p_{}^2`$ open the possibility that expectation values of operators of lower twist mix with expectation values of operators of higher twist by virtue of these terms, cf. e.g. . The light cone expansion in the non–forward case bears therefore additional complications to be dealt with. The impact of the hadronic mass scales $`p_+^2,p_{}^2`$ and $`p_+.p_{}`$, which is of relevance for the contributions beyond leading twist, can in general not be dealt with as target mass– or final–state mass corrections only cf. for a discussion. Furthermore, the scales $`p_+^2`$ and $`p_{}^2`$ dependend on the factorization of non–perturbative quantities as the distribution amplitudes and are arbitrary in this sense. In size they are comparable to the mass scales of the expectaion values of higher twist operators. ### 2.1 Operator Structure The product of the two currents $`\widehat{T}_{\mu \nu }(x)=iRT\left[J_\mu \left({\displaystyle \frac{x}{2}}\right)J_\nu \left({\displaystyle \frac{x}{2}}\right)S\right]`$ (2.8) is then given by <sup>3</sup><sup>3</sup>3Here and in the following, the normal product symbols for the operator products will be omitted. $`\widehat{T}^{\mu \nu }(x)=e^2{\displaystyle \frac{\stackrel{~}{x}^\lambda }{2\pi ^2(x^2iϵ)^2}}RT\left[\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma ^\mu \gamma ^\lambda \gamma ^\nu \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma ^\mu \gamma ^\lambda \gamma ^\nu \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\right]S.`$ (2.9) $`\stackrel{~}{x}`$ denotes a light–like vector corresponding to $`x`$, $`\stackrel{~}{x}=x+{\displaystyle \frac{\zeta }{\zeta ^2}}[\sqrt{x.\zeta ^2x^2\zeta ^2}x.\zeta ],`$ (2.10) and $`\zeta `$ is a subsidiary vector. The leading order expressions turn out to be independent of $`\zeta `$<sup>4</sup><sup>4</sup>4The respective terms are suppressed $`\sqrt{\|\zeta ^2/q^2\|}`$ and do therefore not contribute at the twist–2 level but might be of relevance for the higher twist terms beginning with twist–3. $`e`$ denotes the charge of the fermion field $`\psi `$, which is either a quark– or an antiquark field. $`\widehat{T}_{\mu \nu }(x)`$, Eq. (2.9), refers therefore to the contribution of one of these fields for a single flavor. In the subsequent treatment there is no essential structural difference considering quark or antiquark fields since we work in the massless approximation. The expectation values of the bilocal quark and antiquark operators between nucleon states introduced below are of course different and do also depend on the quark flavor. The complete Compton amplitude is obtained summing Eq. (2.9) over all quark and antiquark flavors contributing in the kinematic domain considered. The operators $`\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma ^\mu \gamma ^\lambda \gamma ^\nu \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)`$ are bilocal operators on the light ray $`\stackrel{~}{x}`$. These are renormalized and time ordered operators . One may rewrite the operator $`\widehat{T}_{\mu \nu }(x)`$ in terms of a symmetric and an asymmetric contribution by $`\widehat{T}_{\mu \nu }(x)=e^2{\displaystyle \frac{\stackrel{~}{x}^\lambda }{i\pi ^2(x^2iϵ)^2}}\left[S_{\alpha \mu \lambda \nu }O^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})+i\epsilon _{\mu \lambda \nu \sigma }O_5^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})\right],`$ (2.11) where $`S_{\alpha \mu \lambda \nu }=g_{\alpha \mu }g_{\lambda \nu }+g_{\lambda \mu }g_{\alpha \nu }g_{\mu \nu }g_{\lambda \alpha }.`$ (2.12) The essential objects are the bilocal light–ray operators $`O^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma ^\alpha \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma ^\alpha \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\right],`$ (2.13) $`O_5^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma _5\gamma ^\alpha \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)+\overline{\psi }\left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\gamma _5\gamma ^\alpha \psi \left({\displaystyle \frac{\stackrel{~}{x}}{2}}\right)\right].`$ (2.14) These operators do still contain higher twist contributions. In the present paper we are going to discuss the twist–2 contributions only and have therefore to perform a twist decomposition. The scalar twist–2 operators on the light–cone are given by $`O({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})`$ $`=`$ $`\stackrel{~}{x}_\alpha O^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})`$ (2.15) $`O_5({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})`$ $`=`$ $`\stackrel{~}{x}_\alpha O_5^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}}).`$ (2.16) Their general definition off the light–cone can be found in Refs. . The operators satisfy the condition $`\mathrm{}O_{(5)}^{q,\mathrm{traceless}}(\kappa x,\kappa x)=0.`$ (2.17) Here $`\kappa `$ parametrizes the position on the $`x`$–ray. The twist–2 vector operators can now be constructed referring to the scalar operator by $`O_\sigma ^{q,\mathrm{twist2}}(\kappa \stackrel{~}{x},\kappa \stackrel{~}{x})`$ $`=`$ $`{\displaystyle _0^1}𝑑\tau _\sigma O_{\mathrm{traceless}}^q(\kappa \tau x,\kappa \tau x)\|_{x\stackrel{~}{x}}`$ (2.18) $`=`$ $`{\displaystyle _0^1}𝑑\tau \left[_\sigma +\frac{1}{2}(\mathrm{ln}\tau )x_\sigma \mathrm{}\right]O^q(\kappa \tau x,\kappa \tau x)\|_{x=\stackrel{~}{x}}.`$ These operators satisfy the relations $`^\sigma O_\sigma ^{q,\mathrm{traceless}}(\kappa x,\kappa x)=0,\mathrm{}O_\sigma ^{q,\mathrm{traceless}}(\kappa x,\kappa x)=0.`$ (2.19) ### 2.2 Operator Matrix Elements Up to now we have considered an expression for the product of the two electromagnetic currents. As the next step we form matrix elements of the respective operators. Let us first consider the matrix element of the scalar operator. We use the kinematic decomposition into a Dirac– and a Pauli–type contribution, the latter of which vanishes in the case of forward scattering, $`p_{},\eta 0`$. $`e^2p_2,S_2\left\|O({\displaystyle \frac{x}{2}},{\displaystyle \frac{x}{2}})\right\|p_1,S_1`$ $`=`$ $`i\overline{u}(p_2,S_2)\gamma xu(p_1,S_1){\displaystyle Dze^{ixp_z/2}f(z_1,z_2,p_ip_jx^2,p_ip_j,\mu _R^2)}`$ $`+`$ $`i\overline{u}(p_2,S_2)x\sigma p_{}u(p_1,S_1){\displaystyle Dze^{ixp_z/2}g(z_1,z_2,p_ip_jx^2,p_ip_j,\mu _R^2)}`$ Here $`\mu _R`$ denotes the renormalization scale and $`Dz`$ $`=`$ $`{\displaystyle \frac{1}{2}}dz_1dz_2\theta (1z_1)\theta (1+z_1)\theta (1z_2)\theta (1+z_2)`$ $`=`$ $`dz_+dz_{}\theta (1+z_++z_{})\theta (1+z_+z_{})\theta (1z_++z_{})\theta (1z_+z_{})`$ $`p_z`$ $`=`$ $`p_+z_++p_{}z_{}`$ $`z_\pm `$ $`=`$ $`{\displaystyle \frac{1}{2}}(z_2\pm z_1)`$ $`\sigma _{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{i}{2}}\left[\gamma _\alpha \gamma _\beta \gamma _\beta \gamma _\alpha \right].`$ (2.22) We consider all hadronic mass scales $`p_i.p_j0`$ as small compared to the large invariants $`p_\pm .q`$ and $`q^2`$. Therefore the on–shell relations $`\gamma _\mu p_1^\mu u(p_1,S_1)`$ $`=`$ $`0`$ $`\overline{u}(p_2,S_2)\gamma _\mu p_2^\mu `$ $`=`$ $`0`$ (2.23) hold. Under these assumptions we show that the scalar matrix element satisfies the condition (2.17): $`\mathrm{}e^2p_2,S_2\left\|O({\displaystyle \frac{x}{2}},{\displaystyle \frac{x}{2}})\right\|p_1,S_1\|_{x\stackrel{~}{x}}`$ $`=`$ $`i\mathrm{}[\overline{u}(p_2,S_2)\gamma xu(p_1,S_1){\displaystyle }Dze^{ixp_z/2}f(z_1,z_2)`$ $`+\overline{u}(p_2,S_2)x\sigma p_u(p_1,S_1){\displaystyle }Dze^{ixp_z/2}g(z_1,z_2)]\|_{x\stackrel{~}{x}}`$ $`=`$ $`i{\displaystyle Dze^{i\stackrel{~}{x}p(z)/2}f(z_1,z_2)\left[\frac{i}{2}p_{z\mu }+2\stackrel{~}{x}_\mu \left(\frac{i}{2}p_z\right)^2\right]\overline{u}(p_2,S_2)\gamma ^\mu u(p_1,S_1)}`$ $`+`$ $`i{\displaystyle Dze^{i\stackrel{~}{x}p(z)/2}f(z_1,z_2)\left[\frac{i}{2}p_{z\mu }+2\stackrel{~}{x}_\mu \left(\frac{i}{2}p_z\right)^2\right]\overline{u}(p_2,S_2)\sigma ^{\mu \nu }p_\nu u(p_1,S_1)}0.`$ Analogously to the construction of the vector operator out of the scalar operator (2.18) one obtains the matrix element $`e^2p_2,S_2\left\|O^\mu ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})\right\|p_1,S_1`$ $`=`$ $`i{\displaystyle _0^1}d\lambda _x^\mu \{\overline{u}(p_2,S_2)\gamma xu(p_1,S_1){\displaystyle }Dze^{i\lambda xp_z/2}f(z_1,z_2)`$ $`+\overline{u}(p_2,S_2)x\sigma p_{}u(p_1,S_1){\displaystyle }Dze^{i\lambda xp_z/2}g(z_1,z_2)\}\|_{x\stackrel{~}{x}}`$ $`=`$ $`i{\displaystyle }Dz_x^\mu \{\overline{u}(p_2,S_2)\gamma xu(p_1,S_1)e^{ixp_z/2}{\displaystyle _0^1}{\displaystyle \frac{d\lambda }{\lambda ^2}}f({\displaystyle \frac{z_1}{\lambda }},{\displaystyle \frac{z_2}{\lambda }})`$ $`+\overline{u}(p_2,S_2)x\sigma p_{}u(p_1,S_1)e^{ixp_z/2}{\displaystyle _0^1}{\displaystyle \frac{d\lambda }{\lambda ^2}}g({\displaystyle \frac{z_1}{\lambda }},{\displaystyle \frac{z_2}{\lambda }})\}\|_{x\stackrel{~}{x}}.`$ The expectation values of the operators $`O^\mu `$ and $`Q_5^\mu `$ are finally expressed by $`e^2p_2,S_2\left\|O^\mu ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})\right\|p_1,S_1`$ $`=`$ $`i{\displaystyle Dze^{i\stackrel{~}{x}p_z/2}F(z_1,z_2)\left[\overline{u}(p_2,S_2)\gamma ^\mu u(p_1,S_1)\frac{i}{2}p_z^\mu \overline{u}(p_2,S_2)\gamma \stackrel{~}{x}u(p_1,S_1)\right]}`$ $`+`$ $`i{\displaystyle Dze^{i\stackrel{~}{x}p_z/2}G(z_1,z_2)\left[\overline{u}(p_2,S_2)\sigma ^{\mu \nu }p_{}^{}{}_{\nu }{}^{}u(p_1,S_1)\frac{i}{2}p_z^\mu \overline{u}(p_2,S_2)\sigma ^{\alpha \beta }\stackrel{~}{x}_\alpha p_{}^{}{}_{\beta }{}^{}u(p_1,S_1)\right]}.`$ and $`e^2p_2,S_2\left\|O_5^\mu ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})\right\|p_1,S_1`$ $`=`$ $`i{\displaystyle Dze^{i\stackrel{~}{x}p_z/2}F_5(z_1,z_2)\left[\overline{u}(p_2,S_2)\gamma _5\gamma ^\mu u(p_1,S_1)\frac{i}{2}p_z^\mu \overline{u}(p_2,S_2)\gamma _5\gamma \stackrel{~}{x}u(p_1,S_1)\right]}`$ $`+`$ $`i{\displaystyle Dze^{i\stackrel{~}{x}p_z/2}G_5(z_1,z_2)\left[\overline{u}(p_2,S_2)\gamma _5\sigma ^{\mu \nu }p_{}^{}{}_{\nu }{}^{}u(p_1,S_1)\frac{i}{2}p_z^\mu \overline{u}(p_2,S_2)\gamma _5\sigma ^{\alpha \beta }\stackrel{~}{x}_\alpha p_{}^{}{}_{\beta }{}^{}u(p_1,S_1)\right]},`$ respectively. The new functions $`F,GH`$ and $`F_5,G_5H_5`$ are defined by $`H_{(5)}(z_1,z_2)={\displaystyle _0^1}{\displaystyle \frac{d\lambda }{\lambda ^2}}h_{(5)}({\displaystyle \frac{z_1}{\lambda }},{\displaystyle \frac{z_2}{\lambda }}),`$ (2.28) with $`h_{(5)}f,g,f_5,g_5`$. In the same manner as for the scalar operator it can be shown that the conditions (2.19) for the vector operator are satisfied by the matrix element. The functions $`H`$ and $`H_5`$ obey the relations $`H(z_1,z_2)`$ $`=`$ $`H(z_1,z_2)`$ (2.29) $`H_5(z_1,z_2)`$ $`=`$ $`+H_5(z_1,z_2),`$ (2.30) which follow from Eqs. (2.13,2.14) interchanging $`x`$ and $`x`$, in the limit in which contributions due to $`\zeta `$, Eq. (2.10), can be disregarded. The contributions due to the above operators to the Compton amplitude read $`T_{\mu \nu }(p_+,p_{},q)`$ $`=`$ $`i{\displaystyle d^4xe^{iqx}p_2,S_2\|T(J_\mu (x/2)J_\nu (x/2))\|p_1,S_1}`$ (2.31) $`=`$ $`{\displaystyle }d^4xe^{iqx}\{{\displaystyle \frac{\stackrel{~}{x}^\lambda }{i\pi ^2(x^2i\epsilon )^2}}[S_{\alpha \mu \lambda \nu }p_2\left\|O^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})\right\|p_1`$ $`+i\epsilon _{\mu \lambda \nu \sigma }p_2\left\|O_5^\alpha ({\displaystyle \frac{\stackrel{~}{x}}{2}},{\displaystyle \frac{\stackrel{~}{x}}{2}})\right\|p_1]\}.`$ We rewrite Eq. (2.31) referring to the distribution functions $`F_{(5)},G_{(5)}`$, $`T_{\mu \nu }(p_+,p_{},q)={\displaystyle d^4xe^{i(qp_z/2)x}\frac{\stackrel{~}{x}^\lambda }{\pi ^2(x^2i\epsilon )^2}}`$ $`\times \{S_{\alpha \mu \lambda \nu }\{{\displaystyle }Dz[\overline{u}(p_2,S_2)\gamma ^\alpha u(p_1,S_1){\displaystyle \frac{i}{2}}p_z^\mu \overline{u}(p_2,S_2)\gamma \stackrel{~}{x}u(p_1,S_1)]F(z_+,z_{})`$ $`+{\displaystyle }Dz[\overline{u}(p_2,S_2)\sigma ^{\alpha \beta }p_{}^{}{}_{\beta }{}^{}u(p_1,S_1){\displaystyle \frac{i}{2}}p_z^\alpha \overline{u}(p_2,S_2)\sigma ^{\beta \gamma }\stackrel{~}{x}_\beta p_{}^{}{}_{\gamma }{}^{}u(p_1,S_1)]G(z_+,z_{})\}`$ $`+i\epsilon _{\mu \lambda \nu \sigma }\{{\displaystyle }Dz[\overline{u}(p_2,S_2)\gamma _5\gamma ^\sigma u(p_1,S_1){\displaystyle \frac{i}{2}}p_z^\sigma \overline{u}(p_2,S_2)\gamma _5\gamma \stackrel{~}{x}u(p_1,S_1)]F_5(z_+,z_{})`$ $`+{\displaystyle }Dz[\overline{u}(p_2,S_2)\gamma _5\sigma ^{\sigma \alpha }p_{}^{}{}_{\alpha }{}^{}u(p_1,S_1){\displaystyle \frac{i}{2}}p_z^\sigma \overline{u}(p_2,S_2)\gamma _5\sigma ^{\alpha \beta }\stackrel{~}{x}_\alpha p_{}^{}{}_{\beta }{}^{}u(p_1,S_1)]G_5(z_+,z_{})\}\}.`$ Finally the Fourier transform to momentum space is performed, $`T_{\mu \nu }(q,p_+,p_{})`$ $`=`$ $`2{\displaystyle }Dz{\displaystyle \frac{1}{Q^2+i\epsilon }}\{\overline{u}(p_2,S_2)\mathrm{\Gamma }_{\mu \nu }^F(q,p_+,p_{})u(p_1,S_1)F(z_+,z_{})`$ (2.33) $`+\overline{u}(p_2,S_2)\mathrm{\Gamma }_{\mu \nu }^{F5}(q,p_+,p_{})u(p_1,S_1)F_5(z_+,z_{})`$ $`+\overline{u}(p_2,S_2)\mathrm{\Gamma }_{\mu \nu }^G(q,p_+,p_{})u(p_1,S_1)G(z_+,z_{})`$ $`+\overline{u}(p_2,S_2)\mathrm{\Gamma }_{\mu \nu }^{G5}(q,p_+,p_{})u(p_1,S_1)G_5(z_+,z_{})\}.`$ The matrices $`\mathrm{\Gamma }_{\mu \nu }^O`$ are given by $`\mathrm{\Gamma }_{\mu \nu }^F(q,p_z)`$ $`=`$ $`\left[Q_\mu \gamma _\nu +Q_\nu \gamma _\mu g_{\mu \nu }\gamma _\alpha Q^\alpha \right]`$ (2.34) $`{\displaystyle \frac{1}{2}}\left[p_{z\mu }\gamma _\nu +p_{z\nu }\gamma _\mu g_{\mu \nu }\gamma _\alpha p_z^\alpha \right]`$ $`+{\displaystyle \frac{1}{Q^2+i\epsilon }}\gamma _\alpha Q^\alpha [Q_\nu p_{z\mu }+Q_\mu p_{z\nu }g_{\mu \nu }Q.p_z]`$ $``$ $`\left[q_\mu \gamma _\nu +q_\nu \gamma _\mu g_{\mu \nu }\gamma _\alpha q^\alpha \right]\left[p_{z\mu }\gamma _\nu +p_{z\nu }\gamma _\mu \right]`$ $`+{\displaystyle \frac{1}{Q^2+i\epsilon }}\gamma _\alpha q^\alpha [p_{z\nu }p_{z\mu }+q_\nu p_{z\mu }+q_\mu p_{z\nu }g_{\mu \nu }q.p_z]`$ $`\mathrm{\Gamma }_{\mu \nu }^{F5}(q,p_z)`$ $`=`$ $`i\gamma _5\epsilon _{\mu \nu \lambda \sigma }\left[Q^\lambda \gamma ^\sigma {\displaystyle \frac{1}{2}}p_z^\sigma \gamma ^\lambda +{\displaystyle \frac{1}{Q^2+i\epsilon }}Q^\lambda p_z^\sigma \gamma _\alpha Q^\alpha \right]`$ (2.35) $``$ $`i\gamma _5\epsilon _{\mu \nu \lambda \sigma }\left[q^\lambda \gamma ^\sigma +{\displaystyle \frac{1}{Q^2+i\epsilon }}q^\lambda p_z^\sigma \gamma _\alpha q^\alpha \right]`$ $`\mathrm{\Gamma }_{\mu \nu }^G(q,p_z)`$ $`=`$ $`\left[Q_\mu \sigma _{\nu \alpha }p_{}^\alpha +Q_\nu \sigma _{\mu \alpha }p_{}^\alpha g_{\mu \nu }\sigma _{\alpha \beta }p_{}^\beta Q^\alpha \right]`$ $`{\displaystyle \frac{1}{2}}\left[p_{z\mu }\sigma _{\nu \alpha }p_{}^\alpha +p_{z\nu }\sigma _{\mu \alpha }p_{}^\alpha g_{\mu \nu }\sigma _{\beta \alpha }p_{}^\alpha p_z^\beta \right]`$ $`+{\displaystyle \frac{1}{Q^2+i\epsilon }}\sigma _{\beta \alpha }p_{}^\alpha Q^\beta [Q_\nu p_{z\mu }+Q_\mu p_{z\nu }g_{\mu \nu }Q.p_z]`$ $``$ $`\left[q_\mu \sigma _{\nu \alpha }p_{}^\alpha +q_\nu \sigma _{\mu \alpha }p_{}^\alpha g_{\mu \nu }\sigma _{\beta \alpha }p_{}^\alpha q^\beta \right]\left[p_{z\mu }\sigma _{\nu \alpha }p_{}^\alpha +p_{z\nu }\sigma _{\mu \alpha }p_{}^\alpha \right]`$ $`+{\displaystyle \frac{1}{Q^2+i\epsilon }}\sigma _{\beta \alpha }p_{}^\alpha q^\beta [p_{z\mu }p_{z\nu }+q_\nu p_{z\mu }+q_\mu p_{z\nu }g_{\mu \nu }q.p_z]`$ $`\mathrm{\Gamma }_{\mu \nu }^{G5}(q,p_z)`$ $`=`$ $`i\gamma _5\epsilon _{\mu \nu \lambda \sigma }\left[Q^\lambda \sigma ^{\sigma \alpha }p_\alpha {\displaystyle \frac{1}{2}}p_z^\sigma \sigma ^{\lambda \alpha }p_\alpha +{\displaystyle \frac{1}{Q^2+i\epsilon }}Q^\lambda p_z^\sigma \sigma ^{\alpha \beta }Q_\alpha p_\beta \right]`$ (2.37) $``$ $`i\gamma _5\epsilon _{\mu \nu \lambda \sigma }\left[q^\lambda \sigma ^{\sigma \alpha }p_\alpha +{\displaystyle \frac{1}{Q^2+i\epsilon }}q^\lambda p_z^\sigma \sigma ^{\alpha \beta }q_\alpha p_\beta \right],`$ where $`Q=q{\displaystyle \frac{p_z}{2}}.`$ (2.38) The equivalence–sign denotes that the other contributions vanish between the bi–spinor states according to the assumption that $`p_i.p_j0`$. ### 2.3 Lorentz Structure Let us now discuss the Lorentz structure of the Compton amplitude in more detail. The matrices $`\mathrm{\Gamma }_{\mu \nu }^O/(Q^2+i\epsilon )`$ consist out of two parts which are $`1/Q^2`$ and $`1/Q^4`$, respectively. The numerators of the former terms depend on $`q`$ only, whereas the second terms contain the vectors $`q`$ and $`p_z=z_+p_++z_{}p_{}`$. The $`(z_+,z_{})`$–dependence may be factored out of the quantities $`\overline{u}(p_2,S_2)\gamma _{\mu _1}\mathrm{}\gamma _{\mu _k}(\gamma _5)u(p_1,S_1)`$ (2.39) and the form factors related to the contributions $`1/Q^2`$ and $`1/Q^4`$ are of different structure, graded by the $`z_\pm `$ dependence in the numerators. Due to the different numerator structure individual variations in $`q`$, $`p_+`$ and $`p_{}`$ may allow to disentangle the different contributions experimentally. The Compton amplitude is given by $`T_{\mu \nu }(q,p_+,p_{})`$ $`=`$ $`2\overline{u}(p_2,S_2)[\mathrm{\Gamma }_{\mu \nu }^F(q,p_+,p_{})+\mathrm{\Gamma }_{\mu \nu }^{F5}(q,p_+,p_{})`$ (2.40) $`+\mathrm{\Gamma }_{\mu \nu }^G(q,p_+,p_{})+\mathrm{\Gamma }_{\mu \nu }^{G5}(q,p_+,p_{})]u(p_1,S_1),`$ with $`\mathrm{\Gamma }_{\mu \nu }^F(q,p_+,p_{})`$ $`=`$ $`\left[q_\mu \gamma _\nu +q_\nu \gamma _\mu g_{\mu \nu }\gamma _\alpha q^\alpha \right]F_1(\xi ,\eta )`$ (2.41) $`\gamma _\mu F_{1,\nu }(\xi ,\eta )\gamma _\nu F_{1,\mu }(\xi ,\eta )+\gamma _\alpha q^\alpha F_{2,\mu \nu }(\xi ,\eta )`$ $`\mathrm{\Gamma }_{\mu \nu }^{F5}(q,p_+,p_{})`$ $`=`$ $`i\gamma _5\epsilon _{\mu \nu \lambda \sigma }\left[q^\lambda \gamma ^\sigma F_1^5(\xi ,\eta )+q^\lambda \gamma _\alpha q^\alpha F_2^{\sigma ,5}(\xi ,\eta )\right]`$ (2.42) $`\mathrm{\Gamma }_{\mu \nu }^G(q,p_+,p_{})`$ $`=`$ $`\left[q_\mu \sigma _{\nu \alpha }p_{}^\alpha +q_\nu \sigma _{\mu \alpha }p_{}^\alpha g_{\mu \nu }\sigma _{\beta \alpha }p_{}^\alpha q^\beta \right]G_1(\xi ,\eta )`$ (2.43) $`\sigma _{\mu \alpha }p_{}^\alpha G_{1,\nu }(\xi ,\eta )\sigma _{\nu \alpha }p_{}^\alpha G_{1,\mu }(\xi ,\eta )p_{}^\alpha +\sigma _{\beta \alpha }p_{}^\alpha q^\beta G_{2,\mu \nu }(\xi ,\eta )`$ $`\mathrm{\Gamma }_{\mu \nu }^{G5}(q,p_+,p_{})`$ $`=`$ $`i\gamma _5\epsilon _{\mu \nu \lambda \sigma }\left[q^\lambda \sigma ^{\sigma \alpha }p_\alpha G_1^5(\xi ,\eta )+q^\lambda \sigma ^{\alpha \beta }q_\alpha p_\beta G_2^{\sigma ,5}(\xi ,\eta )\right].`$ (2.44) The different amplitude functions are given by $`H_1(\xi ,\eta )`$ $`=`$ $`{\displaystyle Dz\frac{1}{Q^2+i\epsilon }H(z_+,z_{})}`$ (2.45) $`H_k^\sigma (\xi ,\eta )`$ $`=`$ $`{\displaystyle Dz\frac{p_+^\sigma z_++p_{}^\sigma z_{}}{(Q^2+i\epsilon )^k}H(z_+,z_{})}={\displaystyle Dz\frac{p_+^\sigma t+\pi _\sigma z_{}}{(Q^2+i\epsilon )^k}H(z_+,z_{})}`$ (2.46) $`H_{2,\mu \nu }(\xi ,\eta )`$ $`=`$ $`{\displaystyle }Dz{\displaystyle \frac{1}{(Q^2+i\epsilon )^2}}[p_{z\mu }p_{z\nu }+q_\nu p_{z\mu }+q_\mu p_{z\nu }g_{\mu \nu }q.p_z]H(z_+,z_{})`$ $`=`$ $`{\displaystyle }Dz{\displaystyle \frac{1}{(Q^2+i\epsilon )^2}}[p_{+\mu }p_{+\nu }t^2+(q_\nu p_{+\mu }+q_\mu p_{+\nu })tg_{\mu \nu }q.p_z`$ $`\pi _\mu \pi _\nu z_{}^2+(q_\nu \pi _\mu +q_\mu \pi _\nu )z_{}+(p_{+\nu }\pi _\mu +p_{+\mu }\pi _\nu )tz_{}]`$ $`\times H(z_+,z_{}),`$ where $`t`$ $`=`$ $`z_++\eta z_{}`$ (2.48) $`\pi _\sigma `$ $`=`$ $`p_\sigma \eta p_{+\sigma }`$ (2.49) and $`q.p_z=q.p_+t`$. Here $`H`$ denotes either $`F`$ or $`G`$. The amplitude functions in Eqs. (2.412.44) depend on the scaling variables $`\xi `$ and $`\eta `$. In the limit of forward scattering ($`\eta 0,\xi x_B`$) the structure functions $`F_{1,2}(x_B)`$ and $`g_{1,2}(x_B)`$ occur as the limit of the amplitude functions of respective helicity projections from the Dirac–type terms. The relations between these amplitude functions are derived in Section 6. ## 3 Kinematic Relations For convenience we construct the helicity basis for the photons $`\gamma _1^{}`$ and $`\gamma _2^{()}`$ in the Breit–frame $`p_+`$ $`=`$ $`p_1+p_2=(2E_p;\stackrel{}{0})`$ $`p_{}`$ $`=`$ $`p_1p_2=(0;2\stackrel{}{p})=(0;0,0,2p_3)`$ $`q`$ $`=`$ $`\frac{1}{2}\left(q_1+q_2\right)=(q_0;q_1,0,q_3).`$ (3.1) The matrix elements $$𝖳_{kl}=\epsilon _{2,k}^\mu T_{\mu \nu }\epsilon _{1,l}^\nu ,k,lϵ\{0,1,2,3\}$$ (3.2) are independent of the reference frame. In the generalized Bjorken region we classify the various contributions to $`𝖳_{kl}`$ by their $`\nu `$–dependence, where the leading terms are kept respectively. There the kinematic invariants are given by $`q_1.q_1`$ $`=`$ $`\nu (\xi \eta )`$ (3.3) $`q_2.q_2`$ $`=`$ $`\nu (\xi +\eta )`$ (3.4) $`q.p_+`$ $`=`$ $`\nu `$ (3.5) $`q.p_{}`$ $`=`$ $`\eta \nu `$ (3.6) $`q.q`$ $`=`$ $`\xi \nu `$ (3.7) $`q.p_z`$ $`=`$ $`q^2Q^2=(z_++z_{}\eta )\nu t\nu `$ (3.8) $`p_+^2`$ $``$ $`p_{}^2p_+p_{}0.`$ (3.9) To define the helicity basis we introduce the two reference vectors $`n_0`$ $`=`$ $`(1;0,0,0)`$ (3.10) $`n_2`$ $`=`$ $`(0;0,1,0).`$ (3.11) The polarization vectors of the photons $`\gamma _1^{}`$ and $`\gamma _2^{}`$ are then given by $`\epsilon _{0\mu }^{(1)}`$ $`=`$ $`{\displaystyle \frac{q_{1\mu }}{\sqrt{\|q_1^2\|}}}={\displaystyle \frac{q_{1\mu }}{\nu ^{1/2}}}{\displaystyle \frac{1}{\sqrt{\|\xi \eta \|}}}`$ $`\epsilon _{0\mu }^{(2)}`$ $`=`$ $`{\displaystyle \frac{q_{2\mu }}{\sqrt{\|q_2^2\|}}}={\displaystyle \frac{q_{2\mu }}{\nu ^{1/2}}}{\displaystyle \frac{1}{\sqrt{\|\xi +\eta \|}}}`$ (3.12) $`\epsilon _{1\mu }^{(i)}`$ $`=`$ $`n_{2\mu }`$ (3.13) $`\epsilon _{2\mu }^{(i)}`$ $`=`$ $`{\displaystyle \frac{1}{N_{2i}}}\epsilon _{\mu \alpha \beta \gamma }n_0^\alpha n_2^\beta q_i^\gamma `$ (3.14) $`\epsilon _{3\mu }^{(i)}`$ $`=`$ $`{\displaystyle \frac{1}{N_{3i}}}[q_{i\mu }q_i.n_0n_{0\mu }q_i.q_i],`$ (3.15) with $`i=1,2`$ and $`N_{21}`$ $`=`$ $`{\displaystyle \frac{\nu }{\mu }}\sqrt{\left\|1+{\displaystyle \frac{\mu ^2}{\nu }}(\xi \eta )\right\|}`$ $`N_{22}`$ $`=`$ $`{\displaystyle \frac{\nu }{\mu }}\sqrt{\left\|1+{\displaystyle \frac{\mu ^2}{\nu }}(\xi +\eta )\right\|}`$ $`N_{31}`$ $`=`$ $`{\displaystyle \frac{\nu ^{3/2}}{\mu }}\sqrt{\|\xi \eta \|}\sqrt{\left\|1+{\displaystyle \frac{\mu ^2}{\nu }}(\xi \eta )\right\|}`$ $`N_{32}`$ $`=`$ $`{\displaystyle \frac{\nu ^{3/2}}{\mu }}\sqrt{\|\xi +\eta \|}\sqrt{\left\|1+{\displaystyle \frac{\mu ^2}{\nu }}(\xi +\eta )\right\|},`$ (3.16) with $`\mu ^2=p_+^2`$. With respect to the vector $`q_2`$, Eq. (3.4), the above notation applies for space- or time-like vectors, for which $`\|\xi \pm \eta \|0`$ and all the above terms are regular. The polarization vectors obey $`\epsilon _{k\mu }^{(i)}.\epsilon _l^{(i)\mu }=s_k\delta _{kl},`$ (3.17) with $`s_k=1`$ for $`k=0,1,2`$ and $`s_k=+1`$ for $`k=3`$ if both vectors are space-like and analogous relations if $`q_2`$ is time-like. $`\epsilon _{1,2\mu }^{(i)}`$ are the transversal and $`\epsilon _{3\mu }^{(i)}`$ is the longitudinal polarization vector of the photon $`\gamma _i^{}`$. As will be shown below the hadronic mass scale $`\mu ^2=p_+^2`$ occurring as a normalization parameter cancels in the helicity projections of the Compton amplitude in leading order in $`1/\nu `$, cf. Section 5. For the limit to the forward case and the understanding of the results in the limit $`\mu ^2\nu `$ it is useful to rewrite the polarization vectors (33.15) as $`\epsilon _{0\rho }^{(1(2))}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\|\xi \|}}}\left[{\displaystyle \frac{q_\rho }{\nu ^{1/2}}}\pm {\displaystyle \frac{p_\rho }{2\nu ^{1/2}}}\right]{\displaystyle \frac{1}{\sqrt{\|1\eta /\xi \|}}}`$ (3.18) $`\epsilon _{1\rho }^{(1(2))}`$ $`=`$ $`n_{2\rho }`$ (3.19) $`\epsilon _{2\rho }^{(1(2))}`$ $`=`$ $`{\displaystyle \frac{\mu }{\nu }}\left\|1{\displaystyle \frac{\mu ^2}{2\nu }}(\xi \eta )\right\|\epsilon _{\rho \alpha \beta \gamma }n_0^\alpha n_2^\beta \left(q^\gamma \pm {\displaystyle \frac{1}{2}}p_{}^\gamma \right)`$ (3.20) $`\epsilon _{3\rho }^{(1(2))}`$ $`=`$ $`{\displaystyle \frac{1}{\nu ^{1/2}}}{\displaystyle \frac{1}{\sqrt{\|\xi \|}}}\left\|1{\displaystyle \frac{\mu ^2}{2\nu }}(\xi \eta )\right\|\left[q_\rho \pm {\displaystyle \frac{1}{2}}p_\rho +\mu n_{0\rho }(\xi \eta )\right]{\displaystyle \frac{1}{\sqrt{\|1\eta /\xi \|}}}.`$ (3.21) In the limit $`p_\rho ,\eta 0`$ the polarization vectors of the first and the second photon become identical. On the other hand, in the limit that terms of $`O(\mu ^2)`$ can be neglected against terms of $`O(\nu )`$ the polarization vectors $`\epsilon _{3\rho }^{(i)}`$ and $`\epsilon _{0\rho }^{(i)}`$ become identical for $`i=1,2`$. The latter aspect induces that current conservation for the twist–2 contributions, cf. section 4, enforces that helicity projections in the direction of the longitudinal polarization vectors vanish in the generalized Bjorken region in lowest order QCD. If $`q_2`$ is light–like the respective set of polarization vectors reads $`\epsilon _{0\mu }^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}q_0^{(2)}}}q_\mu ={\displaystyle \frac{1}{\sqrt{2}q_0^{(2)}}}(q_0,\stackrel{}{q}_2)`$ (3.22) $`\epsilon _{1\mu }^{(2)}`$ $`=`$ $`n_{2\mu }`$ (3.23) $`\epsilon _{2\mu }^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{q_0^{(2)}}}\epsilon _{\mu \alpha \beta \gamma }n_0^\alpha n_2^\beta q_2^\gamma ,`$ (3.24) with $`q_0^{(2)}={\displaystyle \frac{\nu }{\mu }}.`$ (3.25) A fourth linearly independent vector associated to the above set is $`\stackrel{~}{\epsilon }_{0\mu }^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}q_0^{(2)}}}(q_0,\stackrel{}{q}_2)`$ (3.26) with $`\stackrel{~}{\epsilon }_{0\mu }^{(2)}+\epsilon _{0\mu }^{(2)}=\sqrt{2}n_{0\mu }.`$ (3.27) The vectors $`\stackrel{~}{\epsilon }_{0\mu }^{(2)},\epsilon _{0\mu }^{(2)},\epsilon _{1\mu }^{(2)}`$ and $`\epsilon _{2\mu }^{(2)}`$ span the Minkowski space. For later use we also note that the vectors $`p_+^\mu ,p_{}^\mu ,q^\mu \mathrm{and}n_2^\mu `$ span the Minkowski space for the general case of non–forward scattering, excluding special cases as forward scattering $`p_{}=0`$ and vacuum–meson transition $`p_+=p_{}`$. ## 4 Current Conservation The conservation of the electromagnetic current $`_\mu ^xJ^\mu (x)=0`$ (4.1) implies for the Compton amplitude $`T_{\mu \nu }(p_+,p_{},q)`$ $`=`$ $`i{\displaystyle }d^4e^{iq_2x}p_2,S_2\|RT(J_\mu (0)J_\nu (x)\|p_1,S_1`$ (4.2) $`=`$ $`i{\displaystyle }d^4e^{iq_1x}p_2,S_2\|RT(J_\mu (x)J_\nu (0)\|p_1,S_1`$ the relations $`q_2^\mu T_{\mu \nu }=T_{\mu \nu }q_1^\nu =0.`$ (4.3) Expanding the Compton amplitude according to the (non–local) operator product expansion for deep inelastic non–forward scattering in the generalized Bjorken region the terms obtained forming the matrix elements (2.2) are not necessarily yet the twist–2 contributions only. The explicit calculation shows, that still terms of the order $`(\mu ^2/\nu )^{1/2+k},k0`$ are contained. These terms are of higher twist and vanish for $`\nu \mathrm{}`$ if compared to the leading twist–2 contributions. To prove the current conservation for the twist–2 contributions to the Compton amplitude these terms have to be dealt with in common with the respective higher twist contributions resulting from the operator matrix elements of the higher twist operators. We are firstly considering the helicity projections of the Compton amplitude onto the states $`\epsilon _{0\mu }^{(2)}`$ and $`\epsilon _{0\nu }^{(1)}`$, respectively, if $`q_2`$ is either space- or time-like. The remaining index is contracted by all the four helicity vectors. As each of the sets of four helicity vectors corresponding to the virtual photons $`\gamma _{1,2}^{}`$ spans the complete Minkowski–space, vanishing of the projections in all components proves the current conservation for the twist–2 contributions. We list the individual projections for the contributions to the different distribution functions $`F,F_5,G`$ and $`G_5`$ separately. $`𝖳_{00}^F`$ $`=`$ $`{\displaystyle \frac{2}{\nu }}\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1){\displaystyle \frac{1}{\sqrt{\|\xi ^2\eta ^2\|}}}{\displaystyle DzF(z_+,z_{})}`$ (4.4) $`𝖳_{01}^F,𝖳_{10}^F,𝖳_{02}^F,𝖳_{20}^F`$ $`=`$ $`O\left({\displaystyle \frac{1}{\sqrt{\nu }}}\right)`$ (4.5) $`𝖳_{03}^F`$ $`=`$ $`{\displaystyle \frac{2}{\nu }}\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1){\displaystyle \frac{1}{\|\xi \eta \|\sqrt{\|\xi ^2\eta ^2\|}}}{\displaystyle DzF(z_+,z_{})}`$ (4.6) $`𝖳_{30}^F`$ $`=`$ $`{\displaystyle \frac{2}{\nu }}\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1){\displaystyle \frac{1}{\|\xi +\eta \|\sqrt{\|\xi ^2\eta ^2\|}}}{\displaystyle DzF(z_+,z_{})}`$ (4.7) $`𝖳_{00}^{F5}`$ $`=`$ $`0`$ (4.8) $`𝖳_{01}^{F5},𝖳_{10}^{F5},𝖳_{02}^{F5},𝖳_{20}^{F5}`$ $`=`$ $`O\left({\displaystyle \frac{1}{\sqrt{\nu }}}\right)`$ (4.9) $`𝖳_{03}^{F5},𝖳_{30}^{F5}`$ $`=`$ $`O\left({\displaystyle \frac{1}{\nu }}\right)`$ (4.10) $`𝖳_{00}^G`$ $`=`$ $`{\displaystyle \frac{2}{\nu }}\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1){\displaystyle \frac{1}{\sqrt{\|\xi ^2\eta ^2\|}}}{\displaystyle DzG(z_+,z_{})}`$ (4.11) $`𝖳_{01}^G,𝖳_{10}^G,𝖳_{02}^G,𝖳_{20}^G`$ $`=`$ $`O\left({\displaystyle \frac{1}{\sqrt{\nu }}}\right)`$ (4.12) $`𝖳_{03}^G`$ $`=`$ $`{\displaystyle \frac{2}{\nu }}\overline{u}(p_2,S_2)\sigma _{\beta \alpha }p_{}^\alpha q^\beta u(p_1,S_1){\displaystyle \frac{1}{\|\xi \eta \|\sqrt{\|\xi ^2\eta ^2\|}}}{\displaystyle DzG(z_+,z_{})}`$ $`𝖳_{30}^G`$ $`=`$ $`{\displaystyle \frac{2}{\nu }}\overline{u}(p_2,S_2)\sigma _{\beta \alpha }p_{}^\alpha q^\beta u(p_1,S_1){\displaystyle \frac{1}{\|\xi +\eta \|\sqrt{\|\xi ^2\eta ^2\|}}}{\displaystyle DzG(z_+,z_{})}`$ $`𝖳_{00}^{G5}`$ $`=`$ $`0`$ (4.15) $`𝖳_{01}^{G5},𝖳_{10}^{G5},𝖳_{02}^{G5},𝖳_{20}^{G5}`$ $`=`$ $`O\left({\displaystyle \frac{1}{\sqrt{\nu }}}\right)`$ (4.16) $`𝖳_{03}^{G5},𝖳_{30}^{G5}`$ $`=`$ $`O\left({\displaystyle \frac{1}{\nu }}\right).`$ (4.17) Note that $`\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1)`$ $``$ $`\nu `$ $`\epsilon _{\alpha ,\beta ,\gamma ,\delta }p_\pm ^\gamma q^\delta `$ $``$ $`\nu ,`$ (4.18) etc. The contributions to $`𝖳_{kl}^{H(5)}`$ of $`O(1/\sqrt{\nu })`$ or $`O(1/\nu )`$ are of higher twist since the expectation values of the twist–2 operators are multiplied by a further mass ratio, cf. . Since the integrals $`{\displaystyle DzH(z_+,z_{})}={\displaystyle _1^{+1}}𝑑z_+{\displaystyle _{1+\|z_+\|}^{+1\|z_+\|}}H(z_+,z_{})=0`$ (4.19) for $`H=F,G`$ vanish, all projections Eqs. (4.44.17) vanish in the Bjorken limit proving current conservation for the twist–2 contributions in the generalized Bjorken region. In the case of forward scattering the evaluation of the current–current operator in terms of the twist–2 operators current conservation is easily obtained, cf. Ref. . In the limit $`p_{},\eta 0`$ the above expressions (4.44.17) vanish for all values of $`\nu `$ by virtue of Eq. (4.19). Unlike the forward case which is characterized by only two kinematic vectors $`q`$ and $`p`$, the presence of the vector $`p_{}`$ in the non–forward case allows for twist–3 terms also for spin–averaged matrix elements. These are of $`O(1/\sqrt{\nu })`$ relative to the twist–2 contributions. The above conclusions hold analogously if $`q_2`$ is a light–like vector. Here the polarization vectors Eqs. (3.223.24, 3.26) have to be used. ## 5 The Helicity Projections of the Compton Amplitude In the following we calculate the helicity projections of the twist–2 contributions to the Compton amplitude (2.40) representing the matrices $`\mathrm{\Gamma }_{\mu \nu }^O`$ (2.412.44) in the generalized Bjorken region. The unpolarized Dirac–type terms are<sup>5</sup><sup>5</sup>5Here and in the following we will use the terms unpolarized and polarized for the symmetric and anti-symmetric contributions to the Compton amplitude, respectively. In the spin-averaged case the anti-symmetric terms cancel. : $`𝖳_{11}^F`$ $`=`$ $`2\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1)\left[F_1(\xi ,\eta )+\epsilon _1^{(2)\mu }\epsilon _1^{(1)\nu }F_{2,\mu \nu }(\xi ,\eta )\right]`$ (5.1) $`𝖳_{22}^F`$ $`=`$ $`2\overline{u}(p_2,S_2)\gamma _\mu q^\mu u(p_1,S_1)\left[F_1(\xi ,\eta )+\epsilon _2^{(2)\mu }\epsilon _2^{(1)\nu }F_{2,\mu \nu }(\xi ,\eta )\right]`$ (5.2) $`𝖳_{kl}^F`$ $``$ $`\left({\displaystyle \frac{1}{\nu }}\right)^{1/2+n}\mathrm{for}\mathrm{the}\mathrm{other}\mathrm{projections}k,lϵ\{1,2,3\}\mathrm{and}n0,`$ (5.3) with $`\epsilon _1^{(2)\mu }\epsilon _1^{(1)\nu }F_{2\mu \nu }(\xi ,\eta )=\epsilon _2^{(2)\mu }\epsilon _2^{(1)\nu }F_{2\mu \nu }(\xi ,\eta )={\displaystyle Dz\frac{q.p_z}{(Q^2+i\epsilon )^2}F(z_+,z_{})},`$ (5.4) leading to $`𝖳_{11}^F=𝖳_{22}^F.`$ (5.5) Correspondingly the polarized Dirac–type terms yield : $`𝖳_{12}^{F5}=𝖳_{21}^{F5}`$ (5.6) $`𝖳_{kl}^{F5}`$ $``$ $`\left({\displaystyle \frac{1}{\nu }}\right)^{1/2+n}\mathrm{for}\mathrm{the}\mathrm{other}\mathrm{projections}k,lϵ\{1,2,3\}\mathrm{and}n0.`$ (5.7) We define the vector $`S_{21}^\sigma :={\displaystyle \frac{1}{2}}\overline{u}(p_2,S_2)\gamma _5\gamma ^\sigma u(p_1,S_1).`$ (5.8) In the case of forward scattering, $`(p_2,S_2)(p_1,S_1)`$, $`S_{21}^\sigma `$ denotes the nucleon spin vector $`S^\sigma `$. $`𝖳_{12}^{F5}`$ is given by $`𝖳_{12}^{F5}`$ $`=`$ $`i\epsilon ^{\mu \lambda \nu \sigma }\epsilon _{1\mu }^{(2)}\epsilon _{2\nu }^{(1)}{\displaystyle Dz\frac{q_\lambda }{Q^2+i\epsilon }\left[S_{21,\sigma }+\frac{q.S_{21}}{Q^2+i\epsilon }p_{z\sigma }\right]F_5(z_+,z_{})}.`$ (5.9) A similar structure is obtained for the Pauli–type terms. The unpolarized contributions read $`𝖳_{11}^G`$ $`=`$ $`2\overline{u}(p_2,S_2)\sigma _{\alpha \beta }q^\alpha p_{}^\beta u(p_1,S_1)\left[G_1(\xi ,\eta )+\epsilon _1^{(2)\mu }\epsilon _1^{(1)\nu }G_{2,\mu \nu }(\xi ,\eta )\right]`$ (5.10) $`𝖳_{22}^G`$ $`=`$ $`2\overline{u}(p_2,S_2)\sigma _{\alpha \beta }q^\alpha p_{}^\beta u(p_1,S_1)\left[G_1(\xi ,\eta )+\epsilon _2^{(2)\mu }\epsilon _2^{(1)\nu }G_{2,\mu \nu }(\xi ,\eta )\right]`$ (5.11) $`𝖳_{kl}^G`$ $``$ $`\left({\displaystyle \frac{1}{\nu }}\right)^{1/2+n}\mathrm{for}\mathrm{the}\mathrm{other}\mathrm{projections}k,lϵ\{1,2,3\}\mathrm{and}n0,`$ (5.12) resulting into $`𝖳_{11}^G=𝖳_{22}^G.`$ (5.13) Here the tensor $`G_{2,\mu \nu }`$ is obtained from Eq. (5.4) substituting $`F`$ into $`G`$. The polarized Pauli–type terms obey : $`𝖳_{12}^{G5}=𝖳_{21}^{G5}`$ (5.14) $`𝖳_{kl}^{G5}`$ $``$ $`\left({\displaystyle \frac{1}{\nu }}\right)^{1/2+n}\mathrm{for}\mathrm{the}\mathrm{other}\mathrm{projections}k,lϵ\{1,2,3\}\mathrm{and}n0,`$ (5.15) with $`𝖳_{12}^{G5}`$ $`=`$ $`i\epsilon ^{\mu \lambda \nu \sigma }\epsilon _1^{(2)\mu }\epsilon _2^{(1)\nu }{\displaystyle Dz\frac{q_\lambda }{Q^2+i\epsilon }\left[\mathrm{\Sigma }_{21,\sigma }+\frac{q.\mathrm{\Sigma }_{21}}{Q^2+i\epsilon }p_{z\sigma }\right]G_5(z_+,z_{})},`$ (5.16) and $`\mathrm{\Sigma }_{21}^\sigma :={\displaystyle \frac{1}{2}}\overline{u}(p_2,S_2)\gamma _5\sigma ^{\sigma \alpha }p_\alpha u(p_1,S_1).`$ (5.17) A comparison of the structure of the projections $`𝖳_{12(21)}^{F5,G5}`$ with the corresponding tensor structure for forward scattering, cf. , shows that the two contributions are related to terms containing the following contributions of structure functions $`q_\lambda S_{\sigma ,21}`$ $``$ $`g_1(x_B)+g_2(x_B)`$ (5.18) $`q_\lambda p_{z\sigma }`$ $``$ $`g_2(x_B).`$ (5.19) In summary, the non–forward Compton amplitude at the level of the twist–2 contributions in lowest order QCD in the generalized Bjorken region is described by two helicity states of both virtual photons. The unpolarized and polarized Dirac– and Pauli–terms are related to two non–perturbative amplitude functions respectively. ## 6 The Integral Relations The amplitude functions describing the twist–2 contributions to the virtual Compton amplitude in the generalized Bjorken region obey relations to which we turn now. ### 6.1 Unpolarized Contributions For the unpolarized contributions to the Compton amplitude the relations $`𝖳_{11}^{F,G}=𝖳_{22}^{F,G}`$ (6.1) hold for the diagonal helicity projections. All other projections vanish at the level of the twist–2 contributions in lowest order QCD in the generalized Bjorken region. Although Eq. (6.1) holds, the partonic interpretation of $`𝖳_{11}^{F,G}`$ and $`𝖳_{22}^{F,G}`$ for non–forward scattering as given in the case of forward scattering by $`F_2(x_B)=2xF_1(x_B){\displaystyle \underset{q}{}}e_q^2x\left[q(x_B)+\overline{q}(x_B)\right]`$ (6.2) has still to be clarified since the representations (5.1, 5.2, 5.10, 5.11) contain yet two types of $`Dz`$–integrals (2.46,5.4), $`𝖧_\mathrm{𝟣}(\xi ,\eta )`$ $`=`$ $`{\displaystyle Dz\frac{\nu }{Q^2+i\epsilon }H(z_+,z_{})}={\displaystyle Dz\frac{H(z_+,z_{})}{\xi +ti\epsilon }}`$ (6.3) $`𝖧_\mathrm{𝟤}(\xi ,\eta )`$ $`=`$ $`{\displaystyle Dz\frac{\nu q.p_z}{(Q^2+i\epsilon )^2}H(z_+,z_{})}={\displaystyle Dz\frac{tH(z_+,z_{})}{(\xi +ti\epsilon )^2}}.`$ (6.4) Changing the integration variables to $`(t,z_{})`$ the integration over $`z_{}`$ can be performed, $`\widehat{H}(t,\eta )`$ $`=`$ $`{\displaystyle _{z_{}^{\mathrm{min}}}^{z_{}^{\mathrm{max}}}}𝑑z_{}H(t\eta z_{},z_{})={\displaystyle _0^1}{\displaystyle \frac{d\lambda }{\lambda ^2}}{\displaystyle _{z_{}^{\mathrm{min}}}^{z_{}^{\mathrm{max}}}}𝑑z_{}h({\displaystyle \frac{t}{\lambda }}\eta {\displaystyle \frac{z_{}}{\lambda }},{\displaystyle \frac{z_{}}{\lambda }})`$ (6.5) $`=`$ $`{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\widehat{h}(z,\eta ),`$ with $`z_{}^{\mathrm{min},\mathrm{max}}={\displaystyle \frac{t\pm 1}{\eta \pm 1}},`$ (6.6) and $`\widehat{h}(z,\eta )={\displaystyle _{\rho _{\mathrm{min}}}^{\rho _{\mathrm{max}}}}𝑑\rho h(z\eta \rho ,\rho ),`$ (6.7) where $`z=t/\lambda ,\rho =z_{}/\lambda `$ and $`\rho _{\mathrm{min}(\mathrm{max})}=z_{,\mathrm{min}(\mathrm{max})}(z/t)`$. As it will turn out below $`\widehat{h}(z,\eta )`$ denotes an amplitude function on the partonic level. This interpretation does not apply to $`\widehat{H}(t,\eta )`$. By partial integration one may rewrite Eq. (6.4). The support of the variable $`t`$ is $`tϵ[1,+1]`$. One obtains $`{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{t}{(\xi +ti\epsilon )^2}}{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\widehat{h}(z,\eta )`$ $`=`$ $`{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\widehat{h}(z,\eta )`$ (6.8) $`{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}\widehat{h}(t,\eta ),`$ which yields $`𝖧_\mathrm{𝟤}(\xi ,\eta )=𝖧_\mathrm{𝟣}(\xi ,\eta ){\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{\widehat{h}(t,\eta )}{\xi +ti\epsilon }}.`$ (6.9) The $`z`$–integral contributions in Eqs. (5.1,5.2,5.10,5.11) cancel and the helicity projections $`T_{11(22)}^{F,G}`$ are $`𝖳_{11(22)}^H(\xi ,\eta ){\displaystyle _1^{+1}}{\displaystyle \frac{\widehat{h}(t,\eta )}{\xi +ti\epsilon }}=𝖯{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{\widehat{h}(t,\eta )}{\xi +t}}i\pi \widehat{h}(\xi ,\eta ),`$ (6.10) which shows that $`𝖳_{11(22)}^H`$ obeys a ‘partonic’ description. For the Dirac–type terms the function $`\widehat{h}(\xi ,\eta )=\widehat{f}(\xi ,\eta )`$ turns into the quark and antiquark densities in the forward limit $`\eta 0,\xi x_B`$. To derive the relations on the Lorentz level, we consider Eq. (2.40). The unpolarized part of $`T_{\mu \nu }`$ contains the following functions : $`H_1(\xi ,\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\nu }}{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}\widehat{𝖧}(t,\eta )`$ (6.11) $`H_k^\sigma (\xi ,\eta )`$ $`=`$ $`{\displaystyle \frac{p_+^\sigma }{(\nu )^k}}{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{t}{(\xi +ti\epsilon )^k}}\widehat{𝖧}(t,\eta )+{\displaystyle \frac{\pi ^\sigma }{(\nu )^k}}{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{(\xi +ti\epsilon )^2}}\stackrel{~}{𝖧}_\mathrm{𝟣}(t,\eta )`$ (6.12) $`H_{2,\mu \nu }(\xi ,\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\nu ^2}}{\displaystyle _1^{+1}}dt{\displaystyle \frac{1}{(\xi +ti\epsilon )^2}}\{[p_{+\mu }p_{+\nu }t^2+(q_\nu p_{+\mu }+q_\mu p_{+\nu })tg_{\mu \nu }q.p_+t]\widehat{𝖧}(t,\eta )`$ (6.13) $`+[(p_{+\nu }\pi _\mu +p_{+\mu }\pi _\nu )t+(q_\nu \pi _\mu +q_\mu \pi _\nu )]\stackrel{~}{𝖧}_\mathrm{𝟣}(t,\eta )\pi _\mu \pi _\nu \stackrel{~}{𝖧}_\mathrm{𝟤}(t,\eta )\},`$ where $`k=1,2`$ and $`\widehat{𝖧}(t,\eta )`$ $`=`$ $`{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\widehat{h}(z,\eta )`$ (6.14) $`\stackrel{~}{𝖧}_𝗄(t,\eta )`$ $`=`$ $`{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\stackrel{~}{h}_k(z,t,\eta ),`$ (6.15) with $`\stackrel{~}{h}_k(z,t,\eta )=\left({\displaystyle \frac{t}{z}}\right)^k{\displaystyle _{\rho _{\mathrm{min}}}^{\rho _{\mathrm{max}}}}𝑑\rho \rho ^kh(z\eta \rho ,\rho ).`$ (6.16) Note that contractions of the vector $`\pi _\sigma `$ in the above equations with one of the vectors $`q^\mu ,p_\pm ^\mu `$ or $`n_2^\mu `$, which span the Minkowski space, yields contributions of at most $`O(\mu ^2)`$ if compared to the large invariants which are of $`O(\nu )`$. The contributions due to these terms are therefore hadronic mass scale corrections or of higher twist. These terms vanish also both for forward scattering and in the case of vacuum–meson transition $`p_+=p_{},\eta =1`$. Due to this we will consider only the terms which do not contain the vectors $`\pi _\sigma `$ in the following. The integrals (6.12,6.13) may be rewritten by partial integration. $`H_2^\sigma (\xi ,\eta )`$ $`=`$ $`{\displaystyle \frac{p_+^\sigma }{\nu ^2}}{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}\left[\widehat{𝖧}(t,\eta )\widehat{h}(t,\eta )\right]+O(\pi ^\sigma )`$ (6.17) $`H_{2,\mu \nu }(\xi ,\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\nu ^2}}{\displaystyle _1^{+1}}dt{\displaystyle \frac{1}{\xi +ti\epsilon }}\{p_{+\mu }p_{+\nu }t[2\widehat{𝖧}(t,\eta )\widehat{h}(t,\eta )]`$ (6.18) $`+[(q_\nu p_{+\mu }+q_\mu p_{+\nu })g_{\mu \nu }q.p_+][\widehat{𝖧}(t,\eta )\widehat{h}(t,\eta )]\}+O(\pi _{\mu (\nu )})`$ Let us define the vector $`P_{21}^\sigma :=\overline{u}(p_2,S_2)\gamma ^\sigma u(p_1,S_1),`$ (6.19) with $`P_{21}^\sigma =p_+^\sigma `$ for forward scattering. For the unpolarized part of $`T^{\mu \nu }`$, Eq. (2.40), $`A^{\mu \nu }`$, one obtains $`A^{\mu \nu }`$ $`=`$ $`2{\displaystyle \frac{q.P_{21}}{\nu }}\left[g^{\mu \nu }{\displaystyle \frac{q^\mu p_+^\nu +q^\nu p_+^\mu }{q.p_+}}\right]{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{𝖥_\mathrm{𝟣}(t,\eta )}{\xi +ti\epsilon }}`$ (6.20) $`+{\displaystyle \frac{2}{\nu }}\left[q^\mu \left(P_{21}^\nu p_+^\nu {\displaystyle \frac{q.P_{21}}{\nu }}\right)+q^\nu \left(P_{21}^\mu p_+^\mu {\displaystyle \frac{q.P_{21}}{\nu }}\right)\right]{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{\widehat{𝖧}(t,\eta )}{\xi +ti\epsilon }},`$ $`{\displaystyle \frac{q.P_{21}}{\nu ^2}}p_+^\mu p_+^\nu {\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{𝖥_\mathrm{𝟤}(t,\eta )}{\xi +ti\epsilon }}`$ $`{\displaystyle \frac{2}{\nu }}\left[p_+^\mu \left(p_{21}^\nu p_+^\nu {\displaystyle \frac{q.P_{21}}{\nu }}\right)+p_+^\nu \left(p_{21}^\mu p_+^\mu {\displaystyle \frac{q.P_{21}}{\nu }}\right)\right]{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{t\widehat{𝖧}(t,\eta )}{\xi +ti\epsilon }}.`$ Here we defined $`𝖥_\mathrm{𝟣}(t,\eta )`$ $`=`$ $`\widehat{h}(t,\eta )`$ (6.21) $`𝖥_\mathrm{𝟤}(t,\eta )`$ $`=`$ $`2t\widehat{h}(t,\eta ).`$ (6.22) The vector $`\mathrm{\Pi }^\mu =P_{21}^\mu p_+^\mu {\displaystyle \frac{q.P_{21}}{\nu }}`$ (6.23) in Eq. (6.20), as also the case for the vector $`\pi ^\mu `$ above, has contractions with the vectors $`q^\mu ,p_\pm ^\mu `$ and $`n_2^\mu `$ which either vanish or are of $`O(\mu ^2)`$ only. Therefore these terms do not contribute at the lowest twist level unlike those due to $`𝖥_\mathrm{𝟣}(t,\eta )`$ and $`𝖥_\mathrm{𝟤}(t,\eta )`$. $`𝖥_\mathrm{𝟤}(t,\eta )=2t𝖥_\mathrm{𝟣}(t,\eta )`$ (6.24) is the generalization of the CallanGross relation for non–forward scattering. These distribution amplitudes have the partonic representation (6.21, 6.22), whereas the other distribution amplitudes of non-leading twist in Eq. (6.20) depend on the function $`\widehat{𝖧}(t,\eta )`$, which is related to $`\widehat{h}(t,\eta )`$ by the integral Eq. (6.14). ### 6.2 Polarized Contributions The matrix element $`𝖳_{12}^{H5}`$, Eq. (5.6), $`𝖳_{12}^{H5}=i\epsilon ^{\mu \lambda \nu \sigma }\epsilon _{1\mu }^{(2)}\epsilon _{2\nu }^{(1)}B_{\lambda \sigma }`$ (6.25) contains the tensor $`B_{\lambda \sigma }={\displaystyle Dz\frac{q_\lambda }{Q^2+i\epsilon }\left[S_{21,\sigma }^H+\frac{q.S_{21}^H}{Q^2+i\epsilon }p_{z\sigma }\right]H_5(z_+,z_{})},`$ with $`S_{21}^H=S_{21}(\mathrm{\Sigma }_{21})`$ for $`H=F(G)`$. It may be rewritten as $`B_{\lambda \sigma }={\displaystyle \frac{1}{\nu }}{\displaystyle Dz\frac{q_\lambda }{\xi +ti\epsilon }\left[S_{21,\sigma }^H\frac{1}{\nu }\frac{tq.S_{21}^H}{\xi +ti\epsilon }p_{+\sigma }+\frac{1}{\nu }\frac{q.S_{21}^H}{\xi +ti\epsilon }z_{}\pi _\sigma \right]H_5(z_+,z_{})}.`$ The latter term in Eq. (6.2) vanishes for forward scattering and in the case of vacuum–meson transition. We perform the integration over $`z_{}`$ and obtain $`B_{\lambda \sigma }`$ $`=`$ $`{\displaystyle \frac{1}{\nu }}q_\lambda S_{21,\sigma }^H{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\widehat{h}_5(z,\eta )`$ (6.27) $`{\displaystyle \frac{1}{\nu ^2}}q_\lambda p_{+\sigma }q.S_{21}^H{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}\left[\widehat{h}_5(t,\eta ){\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\widehat{h}_5(z,\eta )\right]`$ $`{\displaystyle \frac{1}{\nu ^2}}q_\lambda \pi _\sigma q.S_{21}^H{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{1}{\xi +ti\epsilon }}{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\stackrel{~}{h}_5(z,t,\eta ).`$ Here the distribution amplitude $`\stackrel{~}{h}_5`$ denotes the first moment with respect to $`z_{}`$ of the function $`H_5(z_+,z_{})`$ $`\stackrel{~}{h}_5(z,t,\eta )=\left({\displaystyle \frac{t}{z}}\right){\displaystyle _{\rho _{\mathrm{min}}}^{\rho _{\mathrm{max}}}}𝑑\rho \rho h(z\eta \rho ,\rho ),`$ (6.28) whereas $`\widehat{h}(z,\eta )`$, Eq. (6.7), is the corresponding 0th moment. Let us rewrite Eq. (6.27) choosing the notation in analogy to the forward case by $`B_{\lambda \sigma }={\displaystyle \frac{1}{\nu }}q_\lambda {\displaystyle _1^{+1}}{\displaystyle \frac{dt}{\xi +ti\epsilon }}`$ (6.29) $`\times \{S_{21,\sigma }^H[𝖦_\mathrm{𝟣}(t,\eta )+𝖦_\mathrm{𝟤}(t,\eta )]+{\displaystyle \frac{1}{\nu }}p_{+\sigma }q.S_{21}^H𝖦_\mathrm{𝟤}(t,\eta )+{\displaystyle \frac{1}{\nu }}\pi _\sigma q.S_{21}^H𝖦_\mathrm{𝟥}(t,\eta )\}.`$ One obtains $`𝖦_\mathrm{𝟣}(t,\eta )`$ $`:=`$ $`\widehat{h}_5(t,\eta )`$ (6.30) $`𝖦_\mathrm{𝟤}(t,\eta )`$ $`=`$ $`𝖦_\mathrm{𝟣}(t,\eta )+{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}𝖦_\mathrm{𝟣}(z,\eta )`$ (6.31) $`𝖦_\mathrm{𝟥}(t,\eta )`$ $`:=`$ $`{\displaystyle _t^{\mathrm{sign}(t)}}{\displaystyle \frac{dz}{z}}\stackrel{~}{h}_5(z,t,\eta ).`$ (6.32) The $`t`$–integrals in Eq. (6.29) are performed further according to $`{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{A(t)}{\xi +ti\epsilon }}=𝖯{\displaystyle _1^{+1}}𝑑t{\displaystyle \frac{A(t)}{\xi +t}}+i\pi A(\xi ).`$ (6.33) The polarized non–forward distribution amplitude $`𝖦_\mathrm{𝟣}(t,\eta )`$ has a partonic interpretation as in the unpolarized case Eq. (6.10). Eq. (6.31) is the non–forward generalization of the WandzuraWilczek relation. For the special case of vacuum–meson transition according integral relations were discussed in . In Eq. (6.32) a new distribution amplitude $`𝖦_\mathrm{𝟥}(t,\eta )`$ emerges which is, however, not of twist–2 since the contractions of $`\pi _\sigma `$ with the 4–momenta of the scattering process are of $`O(\mu ^2)`$ if compared to the other terms which are of $`O(\nu )`$. The above results show that the well–known results for forward scattering, the Callan-Gross and the WandzuraWilczek relation are not bound to forward scattering only, but are of more general validity. In particular their derivation does not require to use the optical theorem. Furthermore, the above relations hold for quite general functions $`H_{(5)}(z_+,z_{})`$. In deriving the above relations assumptions on their complex structure had not to be made. The $`O(\mu ^2/\nu )`$ corrections are associated with new distribution amplitudes, cf. the contributions $`\pi _\sigma `$ or $`\mathrm{\Pi }_\sigma `$ in Eqs. (2.46, 2.3, 6.12, 6.13, 6.20, 6.29), which are either representable by integral relations of leading twist distribution amplitudes or are higher moments in $`z_{}`$ of the functions $`H(z_+,z_{})`$ and $`H^5(z_+,z_{})`$, respectively. ### 6.3 Forward Scattering After the above considerations the limit to forward scattering is easily performed. Because the Pauli–type terms vanish linearly with $`p_{}`$ only the Dirac–type terms remain. For forward scattering both the functions $`𝖥_\mathrm{𝟣}(\xi ,\eta )`$ and $`𝖦_\mathrm{𝟣}(\xi ,\eta )`$ are real. The absorptive part of the Compton amplitude, the hadronic tensor $`W_{\mu \nu }`$, is given by $`W_{\mu \nu }={\displaystyle \frac{1}{2\pi }}\mathrm{𝖨𝗆}T_{\mu \nu }.`$ (6.34) To derive the forward structure functions we rewrite Eqs. (6.20, 6.27) applying the symmetry relations Eqs. (2.29, 2.30) interchanging $`(z_+,z_{})(z_+,z_{})`$, i.e. $`tt`$, respectively. We define the branches of the functions $`𝖥_\mathrm{𝟣}(t,0)`$ and $`𝖦_\mathrm{𝟣}(t,0)`$ for $`1t<0`$ and $`0<t+1`$ as the antiquark and quark distribution function by $`𝖥_\mathrm{𝟣}(t,0)`$ $`=`$ $`{\displaystyle \underset{q}{}}e_q^2\left[q(t)\theta (t)\overline{q}(t)\theta (t)\right]`$ (6.35) $`𝖦_\mathrm{𝟣}(t,0)`$ $`=`$ $`{\displaystyle \underset{q}{}}e_q^2\left[\mathrm{\Delta }q(t)\theta (t)+\mathrm{\Delta }\overline{q}(t)\theta (t)\right].`$ (6.36) For the unpolarized contributions one obtains for forward scattering, $`q=q_1=q_2,p=p_+/2`$, $`q_\mu A^{\mu \nu }=p^\nu \left[{\displaystyle _1^{+1}}{\displaystyle \frac{2\xi 𝖥_\mathrm{𝟣}(t,0)𝖥_\mathrm{𝟤}(t,0)}{\xi ti\epsilon }}{\displaystyle _1^{+1}}{\displaystyle \frac{2\xi 𝖥_\mathrm{𝟣}(t,0)+𝖥_\mathrm{𝟤}(t,0)}{\xi +ti\epsilon }}\right]=0.`$ (6.37) Taking the absorptive part yields $`\pm 2\xi 𝖥_\mathrm{𝟣}(\pm \xi ,0)=𝖥_\mathrm{𝟤}(\pm \xi ,0).`$ (6.38) The structure functions are now given by $`F_1(x_B)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[𝖥_\mathrm{𝟣}(\xi ,0)𝖥_\mathrm{𝟣}(\xi ,0)\right]={\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}e_q^2\left[q(x_B)+\overline{q}(x_B)\right],`$ (6.39) $`F_2(x_B)`$ $`=`$ $`𝖥_\mathrm{𝟤}(\xi ,0)+𝖥_\mathrm{𝟤}(\xi ,0),`$ (6.40) with the Bjorken variable $`x_B=lim_{\eta 0}\xi `$. The structure functions obey $`F_2(x_B)`$ $`=`$ $`2x_BF_1(x_B)`$ (6.41) the CallanGross relation . Correspondingly, for the polarized part one obtains $`B_{\lambda \sigma }`$ $`=`$ $`{\displaystyle \frac{1}{2\nu }}q_\lambda S_{21,\sigma }^H{\displaystyle _1^{+1}}𝑑t\left[{\displaystyle \frac{𝖦_\mathrm{𝟣}(t,0)+𝖦_\mathrm{𝟤}(t,0)}{\xi +ti\epsilon }}+{\displaystyle \frac{𝖦_\mathrm{𝟣}(t,0)+𝖦_\mathrm{𝟤}(t,0)}{\xi ti\epsilon }}\right]`$ $`{\displaystyle \frac{1}{2\nu ^2}}q_\lambda p_{+\sigma }q.S_{21}^H{\displaystyle _1^{+1}}𝑑t\left[{\displaystyle \frac{𝖦_\mathrm{𝟤}(t,0)}{\xi +ti\epsilon }}+{\displaystyle \frac{𝖦_\mathrm{𝟤}(t,0)}{\xi ti\epsilon }}\right].`$ The absorptive part is described by the structure functions $`g_1(x_B)`$ and $`g_2(x_B)`$ with, cf. (6.27), $`g_1(x_B)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[𝖦_\mathrm{𝟣}(\xi ,0)+𝖦_\mathrm{𝟣}(\xi ,0)\right]={\displaystyle \frac{1}{2}}{\displaystyle \underset{q}{}}e_q^2\left[\mathrm{\Delta }q(x_B)+\mathrm{\Delta }\overline{q}(x_B)\right]`$ (6.43) $`g_2(x_B)`$ $`=`$ $`g_1(x_B)+{\displaystyle _{x_B}^1}{\displaystyle \frac{dz}{z}}g_1(z).`$ (6.44) Eq. (6.44) is the WandzuraWilczek relation<sup>6</sup><sup>6</sup>6 In the presence of electroweak currents five polarized structure functions contribute on the level of twist-2 and for the quark operators also for twist–3. These structure functions are connected by two further relations for the twist–2 contributions, the Dicus-relation and a relation by Blümlein and Kochelev . The three twist–3 relations among the respective contributions to the polarized structure functions due to the quark operators at lowest order QCD have been derived in Ref. recently. One of these relations applies to the case of pure electromagnetic interactions. . ## 7 Conclusions We studied the structure of the virtual Compton amplitude for deep–inelastic non–forward scattering $`\gamma ^{}+p\gamma ^{()^{}}+p^{}`$ in lowest order in QCD in the massless limit. In the generalized Bjorken region $`q.p_+,q^2\mathrm{}`$ the twist–2 contributions to the Compton amplitude were calculated using the non–local operator product expansion for general spin states. In this approximation the Compton amplitude consists of an unpolarized and a polarized Dirac– and Pauli–type amplitude, the latter of which vanishes in the case of forward scattering. The expectation values of the (non–local) twist–2 vector operators in the non–forward case do still contain terms $`1/\nu ^{1/2+k},k0`$, which are contributions of twist–3 and higher to the Compton amplitude. These contributions have to be considered in common with the non–forward expectation values of the higher twist operators. A decomposition of the amplitude was performed with respect to the helicity states of both (virtual) photons. The twist–2 contributions are due to two polarization states only, which are for the unpolarized part $`𝖳_{11}`$ and $`𝖳_{22}`$ and for the polarized part $`𝖳_{12}`$ and $`𝖳_{21}`$. For the twist–2 contributions the gauge invariance of the non–local light cone expansion was proven in the non–forward case in the generalized Bjorken region. The relations between the twist–2 contributions of the unpolarized and polarized amplitude functions were derived. They are the non–forward generalizations of the CallanGross and WandzuraWilczek relations for unpolarized and polarized deep–inelastic forward scattering. The relations for the Dirac and Pauli parts are of the same form. Acknowledgement Our thanks are due to B. Geyer for fruitful discussions and U. Gensch for his constant support. Discussions with A. Tkabladze in an early phase of this project are acknowledged. J.B. would like to thank the Institute of Theoretical Physics at Graz University for their kind hospitality. D.R. likes to thank DESY Zeuthen and the Institute of Theoretical Physics at Graz University for the kind hospitality extended to him, in particular to C.B. Lang, H. Mitter, N. Pucker, and W. Schweiger.
warning/0002/astro-ph0002464.html	ar5iv	text	# A Comparison of the Extra Nuclear X-ray and Radio Features in M87 ## 1 Introduction Eight ROSAT/HRI observations of M87 were made between 1992Jun and 1998Jan to study the X-ray structure and variability of the core and jet (Harris, Biretta, and Junor, 1997 and 1998a). To study the large scale X-ray features of M87, we have used these data to make an image with effective exposure of 230 ksec. A preliminary analysis based on the data then available was presented at the Ringberg Workshop on M87 held in 1997Sep (Harris, Biretta, and Junor 1998b). The radio map of M87 (see Owen, this volume) appears to indicate an exceedingly complex structure, for which it is difficult to make meaningful measurements of isolated features. Generally, we need to define volumes and measure their emission properties. Once you leave the inner (brightest) lobes and jet, this is not easily done. Even for definable filaments and other features, the surface brightness is often a sum of emissivities from various entities along the line of sight. The X-ray map suffers from the same problem although generally there is less fine structure than in the radio (at similar resolutions). An additional complexity for the X-ray analysis is that we cannot be certain that the very large scale X-ray distribution from the hot gas of the Virgo cluster is circularly symmetric. ## 2 M87 as a Wide Angle Tailed Radio Galaxy One possible interpretation of the radio map is that we are viewing a Wide Angle Tailed (WAT) radio galaxy from an angle such that a large fraction of the radio source is well beyond (or in front of) the central part of the galaxy. Thus the observed radio brightness at any given location will usually be the sum of contributions from several different emitting volumes. There are at least three supporting arguments for this hypothesis, although none is conclusive. Without invoking substantial projection effects, it is difficult to maintain continuity between various radio features. In the bright inner region, we see the primary jet bending around to the south, yet the larger scale emission which one might expect to connect to the bright features, is found due West. The same situation occurs on the East side of the source. Without projection, many features display rather strong curvature. Besides the inner lobes, projection effects may be the cause for the apparent curvature in the eastern double ring and in the ‘L’ filament at the south-east edge of the source. The double ring itself is difficult to understand if it were completely in the plane of the sky. The radio galaxy $`0053016`$ in Abell 119 (Feretti et al. 1999; see also Govoni et al. in thses proceedings) shows a helical structure which, if viewed along the arm of the tail, could project to something like the apparent structure of the double ring. Pressure balance between the ambient gas and the non-thermal pressure within radio features is expected for most of the radio structures except those within the very bright central region. We have compared non-thermal pressures for two low brightness regions close to the center with the expected thermal pressure derived from the work of Nulsen and Böhringer (1995). As shown in Table 1, the thermal pressure exceeds the minimum non-thermal pressures for both regions. We integrated the radio spectrum between 10<sup>7</sup> and 10<sup>10</sup> Hz and used spectral index values of 0.8 and 0.9 (from a spectral index map between 74 and 327 MHz provided by N. Kassim). The entries for k=0 represent the minimum pressure for the case of the filling factor, $`\varphi `$=1 and no contribution to the particle energy density from relativistic protons (k=0). The entries with k=100 correspond to either protons having 100 times the energy density of the relativistic electrons or $`\varphi `$=0.01. The last column indicates the radial distance required to reach pressure balance between the thermal gas and the non-thermal radio feature. In all cases, these distances are greater than the projected distance from the center. There are various possible reasons for a relative motion between the ICM and M87 which could produce a WAT structure. Harris et al. (1998b) suggested that the SW X-ray spur (see fig. 1) might be caused by a shock between the ISM of M87 and the ICM disturbed by a merger of the M86 subgroup with the M87 subgroup. We note the relative velocity of M87 and M86, $`1500`$ km s<sup>-1</sup>, is sufficient to shock $`10^7`$ K gas. Other possibilities are that M87 might not be at the center of the potential well of the Virgo cluster or that the cooling flow might be asymmetric if the kinetic energy supplied by the movement of radio structures supplied heat to the local gas in a non symmetric fashion. ## 3 Data The data reduction for the present comparison of radio and X-ray maps involved the following steps. * Five of the eight X-ray observations were corrected for the bug in the standard processing which used incorrect aspect times (see the documents available at the anonymous ftp server: sao-ftp.harvard.edu; cd to pub/rosat/aspfix). The other three observations were made after 1997Jan when this bug had essentially no effect on image quality. * The centroid of the core emission was measured for each observation and each map was then shifted to a common position before stacking. * The radial profile of the total emission was measured in a 90 quadrant towards the north west, centered on the peak of the core emission. * On the basis of this profile, a circularly symmetric King model was constructed to represent the bulk of the cluster emission. The model was then subtracted from the data and the residual was smoothed with Gaussians of various widths. Note that this was more of a procedural process than a true modelling of the emission since the profile does not match a simple King model. It does, however, emphasize the contrast between residual features and the remaining background. * Finally, the radio map was precessed to J2000 and the X-ray map was shifted so as to align the X-ray and radio core emissions. ## 4 Radio/X-ray Coincidence Coincidence between radio and X-ray emissions can be expected when both arise from non-thermal processes: i.e. the X-ray emission is either synchrotron or inverse Compton emission. Although successful models for synchrotron X-ray emission from knot A in the M87 jet have been published (Biretta, Stern, and Harris 1991), the (presumably) older and larger radio structures (fig. 1) under discussion here are not expected to provide the environment to produce the very high energy electrons (Lorentz energy factors $`\gamma `$10<sup>7</sup>) required for X-ray synchrotron emission. Inverse Compton emission however, will be present, both from the 3K background photons (‘IC/3K’) and from other photons such as star light. There are only two regions where a general coincidence of radio and X-ray emissions are found. These are (a) in the southern part of the radio source where a curved feature, the bottom of the (so-called) ‘cobra’ shows low brightness enhancements in both bands (fig. 2); and (b) the ‘eastern radio arm’ extending due East from the core and ending in the region at the center of the double radio rings (fig. 3). Note however that even here, the spatial agreement is not precise. Because most of the radio and residual X-ray emissions are not co-spatial in a detailed sense, we conclude that the bulk of the residual X-ray emission is likely to be thermal. For example, the general association of the radio and X-ray emission might arise from heating of the ISM by the passage of the radio jets. It has become fashionable (at least at this meeting) to suggest that relativistic particles do not necessarily occupy the same volumes defined by radio emitting regions containing the strongest average magnetic field strength. Thus one could have IC/3K emission from regions essentially devoid of radio emission. While this is a valid physical scenario, with our presently available technology, we have no methods to demonstrate that any given X-ray brightness comes from IC processes other than coincidence with known non-thermal emission at other bands and the requisite calculations to demonstrate that the observed intensity and implied field strengths are reasonable. We also note that many of the brighter radio features are devoid of spatially coincident X-ray emission, which tends to argue against the IC process. ## 5 Radio/X-ray Anti-coincidence Anti-coincidence is expected for cases in which radio features exclude the hot gas and thus diminish the X-ray surface brightness for particular lines of sight. Convincing examples of this behavior have been found for Cygnus A (Carilli, Perley, and Harris, 1994) and for NGC 1275 (Böhringer et al. 1993). For simple models of lobe expansion, the change in X-ray surface brightness between lines of sight which intersect the lobe and those that traverse only undisturbed ICM can be either positive or negative. This is because we expect to find a sheath of enhanced density around the inflating lobe. The precise behavior for any given situation depends on the energy range covered by the telescope system and on the density, temperature, and thickness of the sheath (Clarke, Harris, and Carilli, 1997). For M87, most of the source outside the high brightness (inner) radio lobes may not fulfill the simple conditions of an expanding lobe with a well defined wall separating it from the ICM and the unknown projection effects may mask the signatures of cavities in the gas. However, there are two possibilities for this sort of effect. The residual map shows a slight discontinuity in brightness gradient along the northern boundary of the radio source which may indicate a smaller integrated emissivity when looking through the lobe (see fig. 4). The other effect is evidenced by two regions of slightly negative brightness in the residual map. These areas are identified in the color figure by the single contour level with center $``$ 4 to the northeast of the core and a similar feature located $``$ 3 west of the core. It is difficult to evaluate the reality of these features because the subtraction of the model is somewhat arbitrary since there is no suitable pie segment in which to measure the true radial profile of ‘undisturbed’ gas. ## 6 Conclusions * We suspect that projection effects confuse the interpretation of both radio and X-ray features. Harris et al. (1998b) suggested that the SW spur might be caused by a bow shock between the ICM and the ISM, and the resulting change in temperature and density might explain the presence of knot A in the jet. Since the current evidence supports the notion that the spur is a thermal feature, a convincing explanation for the spur will probably await the spectral/spatial capabilities of Chandra and XMM. * Many X-ray and radio features are located in the same general region, but it appears that they do not actually occupy the same volumes. * There is some evidence for edge effects and cavities. ###### Acknowledgements. N. Kassim kindly provided us with an unpublished spectral index map. WJ thanks the National Science Foundation for support under Grant AST-980307. The work at SAO was partially supported by NASA contract 5-99002.
warning/0002/math0002161.html	ar5iv	text	# Geometry without Topology ## 1 Introduction There are several methods of the proper Euclidean geometry<sup>1</sup><sup>1</sup>1We use the term ”Euclidean geometry” as a collective concept with respect to terms ”proper Euclidean geometry” and ”pseudo-Euclidean geometry”. In the first case the eigenvalues of the metric tensor matrix have similar signs, in the second case they have different signs. description. The most old way of description is the axiomatic conception of the Euclidean geometry. The proper Euclidean geometry is described in terms of points, straights and planes, which are determined by their properties in terms of axioms. Some axioms describe properties of natural geometric objects (points, straights and planes), other axioms describe such properties of proper Euclidean geometry as uniformity, isotropy, continuity and degeneracy<sup>2</sup><sup>2</sup>2In general case the set of vectors, parallel to the given vector, forms a cone. Degeneracy of the proper Euclidean geometry means that this cone degenerates into a line.. Real geometry of the space-time is uniform only approximately, and one needs to consider the geometries which should not be uniform, isotropic, continuous and degenerate. In other words, one needs to generalize and modify the proper Euclidean geometry. But it is quite impossible one to modify axioms of the proper Euclidean geometry, and one needs to describe the proper Euclidean geometry in the form, containing numerical characteristics which may be modified rather easily. Such numerical characteristic of the proper Euclidean geometry is the metric $`\rho (P,Q)`$, describing distance between any two points $`P`$ and $`Q`$ of the space. After modification of the Euclidean metric a new geometry appears, which may have other properties than uniformity, isotropy, continuity and degeneracy. Usually for construction of (Riemannian) geometry one uses the following logical scheme $$\begin{array}{ccccccc}\begin{array}{c}\text{coordinate}\\ \text{system}\end{array}& & \begin{array}{c}\text{infinitesimal}\\ \text{distance}\end{array}& & \begin{array}{c}\text{set of }\\ \text{geodesics}\end{array}& \begin{array}{c}\\ \end{array}& \begin{array}{c}\text{geometry}\\ \\ \begin{array}{c}\text{finite}\\ \text{distance}\end{array}\end{array}\end{array}$$ (1.1) As it follows from this scheme for construction of geometry one needs a coordinate system and a system of geodesics. The coordinate system is necessary for introduction of infinitesimal distance. The geodesic is defined as a shortest curve (line), connecting two points. Thus, the considered construction of geometry refers to the concept of a curve. The curve is a topological object, defined as a continuous mapping of the real axis onto the geometrical space of points. As a result the topology is considered usually to be a necessary element of geometry. According to (1.1) one cannot construct geometry without a use of topology (in the form of a curve). Actually the topology is only a mathematical tool, using for construction of geometry. To prove this, it is sufficient to construct geometry without a reference to the topological concept of a curve. We shall make this in the present paper. The geometry is constructed in accord with the following logical scheme $$\begin{array}{cc}\begin{array}{c}\text{finite}\\ \text{distance}\end{array}& \begin{array}{cc}\begin{array}{c}\\ \\ \end{array}& \begin{array}{c}\text{geometry}\\ \\ \begin{array}{c}\text{set of }\\ \text{geodesics}\end{array}\end{array}\end{array}\end{array}$$ (1.2) where geometry is constructed independently of a possibility of the geodesics construction. It is possible such a situation, when the geometry can be constructed, whereas geodesics (the shortest curves) cannot. Such a situation is not exotic, because the real space-time geometry appears to be of such a kind. Timelike geodesics of the space-time are substituted by thin hallow tubes. Thickness of the tubes is microscopic. Describing macroscopic phenomena, one may neglect the tube thickness and substitute the tubes by lines. Then geometry may be considered as a degenerate one (the tubes degenerate into lines). Describing microscopic phenomena, one may not neglect the thickness of tubes, because the thickness of tubes (nondegeneracy of geometry) is a reason of quantum effects. Besides the geometry constructed in accord with the scheme (1.2) is free from such constraints as continuity and degeneracy, imposed by a use of the concept of a curve. To carry out the idea of nondegenerate geometry, let us give some definitions which help us to formulate the problem of generalization and modification of the proper Euclidean geometry. ###### Definition 1.1 The metric space $`M=\{\rho ,\mathrm{\Omega }\}`$ is a set $`\mathrm{\Omega }`$ of points $`P\mathrm{\Omega }`$ with the metric $`\rho `$ given on $`\mathrm{\Omega }\times \mathrm{\Omega }`$ $$\rho :\mathrm{\Omega }\times \mathrm{\Omega }D_+$$ (1.3) $$\rho (P,P)=0,\rho (P,Q)=\rho (Q,P),P,Q\mathrm{\Omega }$$ (1.4) $$D_+=[0,\mathrm{}),\rho (P,Q)=0,\text{if and only if}P=Q,P,Q\mathrm{\Omega }$$ (1.5) $$\rho (P,Q)+\rho (Q,R)\rho (P,R),P,Q,R\mathrm{\Omega }$$ (1.6) ###### Definition 1.2 Any subset $`\mathrm{\Omega }^{}\mathrm{\Omega }`$ of points of the metric space $`M=\{\rho ,\mathrm{\Omega }\}`$, equipped with the metric $`\rho ^{}`$ which is a contraction $`\rho \|_{\mathrm{\Omega }^{}\times \mathrm{\Omega }^{}}`$ of the mapping (1.3). on the set $`\mathrm{\Omega }^{}\times \mathrm{\Omega }^{}`$ is called the metric subspace $`M^{}=\{\rho ^{},\mathrm{\Omega }^{}\}`$ of the metric space $`M=\{\rho ,\mathrm{\Omega }\}`$. It is easy to see that the metric subspace $`M^{}=\{\rho ^{},\mathrm{\Omega }^{}\}`$ is a metric space. ###### Definition 1.3 The metric space $`M=\{\rho ,\mathrm{\Omega }\}`$ is called finite, if the set $`\mathrm{\Omega }`$ contains a finite number of points. The finite metric subspace $`M(𝒫^n)=\{\rho ,𝒫^n\}`$ of $`M=\{\rho ,\mathrm{\Omega }\}`$, consisting of $`n+1`$ points $`𝒫^n\{P_0,P_1,\mathrm{}P_n\}\mathrm{\Omega },`$ $`n=0,1,\mathrm{}`$is called the $`n`$th order metric subspace. The proper Euclidean space may be considered to be a kind of metric space $`E=\{\rho _E,\mathrm{\Omega }\}`$. Being a metric space, the proper Euclidean space and geometry on this space can be described in terms of only metric $`\rho `$ and of finite metric subspaces. The finite metric subspaces $`M(𝒫^n)`$ are the simplest constituents of the metric space. Some properties of finite metric subspaces $`M(𝒫^n)`$ were investigated by Blumenthal , but he did not consider them to be primitive fundamental objects of metric space as we do. Metric space $`M(𝒫^n)`$, consisting of $`n+1`$ points and having nonvanishing length (concept of the length will be defined further), generates in the proper Euclidean space $`n`$-dimensional plane $`_n(𝒫^n)`$, which appears to be an attribute of $`M(𝒫^n)`$ and can be defined in terms of $`M(𝒫^n)`$. ###### Definition 1.4 Elementary geometrical object is a set of points having some metric property. ###### Definition 1.5 Geometrical object is a set of points derived as joins and intersections of elementary geometrical objects. In other words, a geometrical object is a metric subspace $`M_G=\{\rho ,G\}`$, $`G\mathrm{\Omega }`$ of metric space $`M=\{\rho ,\mathrm{\Omega }\}`$. ###### Definition 1.6 Geometry is a totality of all propositions (definitions, axioms and theorems) on properties of geometrical objects. In other words, the geometry is a totality of all propositions on properties of all metric subspaces of the metric space $`M=\{\rho ,\mathrm{\Omega }\}`$. Let us consider some examples of elementary geometrical objects. ###### Definition 1.7 The sphere $`𝒮(O;P)`$, having its center at the point $`O`$ and passing through the point $`P`$, is the set of points $`R\mathrm{\Omega }`$ of the metric space $`M=\{\rho ,\mathrm{\Omega }\},`$ defined by the relation $$𝒮(O;P)=\left\{R\right\|\rho (O,R)=\rho (O,P)\},O,P,R\mathrm{\Omega }$$ The basic points $`O`$ and $`P`$, determining the sphere $`𝒮(O;R)`$, are not equivalent, because $`𝒮(O;R)`$ and $`𝒮(R;O)`$ are different elementary geometrical objects (different spheres). In particular, $`P𝒮(O;R)`$, but $`O𝒮(O;R)`$. The sphere $`𝒮(O;P)`$ is an attribute of zeroth order metric subspaces $`M(O)`$ and $`M(P)`$ (or two points $`O,P`$). ###### Definition 1.8 The circle cylinder $`𝒞(P_1,P_2;P)`$, passing through the point $`P`$, with axis, determined by the basic points $`P_1,P_2`$, is the set of points $`R\mathrm{\Omega }`$ of the metric space $`M=\{\rho ,\mathrm{\Omega }\}`$, defined by the relation $`𝒞(P_1,P_2;P)=\left\{R\right\|S_2(P_1,P_2,R)=S_2(P_1,P_2,P)\},`$ $`P_1,P_2,P,R\mathrm{\Omega }`$ where $`S_2(P_1,P_2,R)`$ is the area of the triangle with vertices at the points $`P_1,P_2,R`$. If the areas of triangles $`\mathrm{}P_1P_2R`$ and $`\mathrm{}P_1P_2P`$ are equal, the heights (radii) dropped from the vertices $`R`$ and $`P`$ of these triangles onto their common base $`P_1P_2`$ (axis of the cylinder) are also equal. The triangle area $`S_2(P_1,P_2,R)`$ can be expressed via metric by the Hero’s formula $$S_2(A,B,C)=\sqrt{p\left(pa\right)\left(pb\right)\left(pc\right)},$$ $$a=\rho (B,C),b=\rho (A,C),c=\rho (A,B),p=\left(a+b+c\right)/2$$ The circle cylinder $`𝒞(P_1,P_2;P)`$ is an attribute of two finite metric subspaces $`M(P_1,P_2)`$ and $`M(P)`$. ###### Definition 1.9 The ellipsoid $`(P_1,P_2;P)`$, having its focuses at the basic points $`P_1,P_2`$ and passing through the point $`P`$ is the set of points $`R\mathrm{\Omega }`$ of the metric space $`M=\{\rho ,\mathrm{\Omega }\}`$, defined by the relation $`(P_1,P_2;P)=\left\{R\right\|\rho (P_1,R)+\rho (P_2,R)=\rho (P_1,P)+\rho (P_2,P)\},`$ $`P_1,P_2,P,R\mathrm{\Omega }`$ The ellipsoid $`(P_1,P_2;P)`$ is an attribute of two finite metric subspaces $`M(P_1,P_2)`$ and $`M(P)`$. If $`P_1P`$ and $`P_2P`$, the points $`P_1,P_2(P_1,P_2;P),`$ but the point $`P(P_1,P_2;P)`$. If the point $`P=P_1`$, the ellipsoid $`(P_1,P_2;P)`$ degenerates into segment $`𝒯_{[P_1P_2]}`$ between the points $`P_1`$ and $`P_2`$ of the straight line $`𝒯_{P_1P_2}`$, passing through the points $`P_1`$ and $`P_2`$. The segment $`𝒯_{[P_1P_2]}`$ is defined as follows. ###### Definition 1.10 The segment $`𝒯_{[P_1P_2]}`$ of the straight between the basic points $`P_1,P_2`$ is the set of points $`R\mathrm{\Omega }`$ of the metric space $`M=\{\rho ,\mathrm{\Omega }\},`$ defined by the relation $$𝒯_{[P_1P_2]}=\left\{R\right\|\rho (P_1,R)+\rho (P_2,R)\rho (P_1,P_2)=0\},P_1,P_2,R\mathrm{\Omega }$$ (1.7) The segment $`𝒯_{[P_1P_2]}`$ is an elementary geometrical object which does not depend on the order of points $`P_1,P_2`$. Besides both basic points $`P_1,P_2𝒯_{[P_1P_2]}.`$ The segment $`𝒯_{[P_1P_2]}`$ is an attribute of the first order metric subspace $`M(P_1,P_2)`$ in the sense that $`𝒯_{[P_1P_2]}`$ is determined by the metric subspace $`M(P_1,P_2)`$ in itself. For instance, the sphere $`𝒮(O;P)`$ is determined by the points $`O,P`$ of the metric subspace $`M(O,P),`$ but not by the metric subspace $`M(O,P)`$ in itself, and the sphere $`𝒮(O;P)`$ is not an attribute of the metric subspace $`M(O,P)`$, but it is an attribute of two zeroth order metric subspaces $`M(O)`$ and $`M(P)`$ (or two points $`O,P`$). ###### Definition 1.11 The elemetary geometrical object which is an attribute of the $`n`$th order metric subspace $`M\left(𝒫^n\right)`$ is the $`n`$th order natural geometric object (the $`n`$th order NGO). Such geometrical objects as a point, an Euclidean straight, and an Euclidean plane are NGOs of the proper Euclidean geometry. The point $`P_0`$ is the zeroth order NGO $`𝒯_{P_0}`$ of the proper Euclidean geometry which is determined by the zeroth order metric subspace $`M\left(P_0\right)=P_0`$. The straight $`𝒯_{P_0P_1}`$ of the proper Euclidean geometry is the first order NGO which is determined by the first order metric subspace $`M(P_0,P_1)`$. It means, in particular, that $`𝒯_{P_0P_1}=𝒯_{P_1P_0}`$. The two-dimensional plane $`𝒯_{P_0P_1P_2}`$ of the proper Euclidean geometry is the second order NGO, determined by the second order metric subspace $`M(P_0,P_1,P_2)`$. It means that the NGO $`𝒯_{P_0P_1P_2}`$ does not depend on the order of basic points $`P_0,P_1,P_2`$, which determine $`𝒯_{P_0P_1P_2}`$. It does not always happen that the second order metric subspace $`M(P_0,P_1,P_2)`$ determines $`𝒯_{P_0P_1P_2}`$. Only $`M(P_0,P_1,P_2)𝒯_{P_0P_1}`$ enables to determine $`𝒯_{P_0P_1P_2}`$. For explicit determination of the $`n`$th order NGO one needs to attribute a length $`\|M\left(𝒫^n\right)\|`$ to any $`n`$th order metric subspace $`M\left(𝒫^n\right)`$ ###### Definition 1.12 The squared length $`\left\|M\left(𝒫^n\right)\right\|^2`$of the $`n`$th order metric subspace $`M\left(𝒫^n\right)\mathrm{\Omega }`$ of the proper Euclidean space $`E=\{\rho _E,\mathrm{\Omega }\}`$ is the real number. $$\left\|M\left(𝒫^n\right)\right\|^2=\left(n!S_n(𝒫^n)\right)^2=F_n\left(𝒫^n\right)$$ where $`S_n(𝒫^n)`$ is the Euclidean volume of the $`(n+1)`$-edr with vertices at points $`𝒫^n\{P_0,P_1,\mathrm{}P_n\}\mathrm{\Omega }`$. In the proper Euclidean geometry the volume $`S_n(𝒫^n)`$ of the $`(n+1)`$-edr and the value $`F_n\left(𝒫^n\right)`$ of the function $`F_n`$, connected with it, can be expressed in terms of metric $`\rho `$ by means of relations $$F_n:\mathrm{\Omega }^{n+1},\mathrm{\Omega }^{n+1}=\underset{k=1}{\overset{n+1}{}}\mathrm{\Omega },n=1,2,\mathrm{}$$ (1.8) $$F_n\left(𝒫^n\right)=det\|\|(𝐏_0𝐏_i.𝐏_0𝐏_k)\|\|,P_0,P_i,P_k\mathrm{\Omega },i,k=1,2,\mathrm{}n$$ (1.9) $`\left(𝐏_0𝐏_i.𝐏_0𝐏_k\right)`$ $``$ $`\mathrm{\Gamma }(P_0,P_i,P_k)\sigma (P_0,P_i)+\sigma (P_0,P_k)\sigma (P_i,P_k),`$ (1.10) $`i,k`$ $`=`$ $`1,2,\mathrm{}n.`$ where the function $`\sigma `$ is defined via metric $`\rho `$ by the relation $$\sigma (P,Q)\frac{1}{2}\rho ^2(P,Q),P,Q\mathrm{\Omega }.$$ (1.11) and $`𝒫^n`$ denotes $`n+1`$ points $`P_0,P_1,\mathrm{},P_n`$ of $`\mathrm{\Omega }`$ $$𝒫^n=\{P_0,P_1,\mathrm{},P_n\}\mathrm{\Omega }$$ (1.12) The function $`\sigma ,`$ called world function , is very important quantity which may be used instead of metric $`\rho `$. In many cases a use of the function $`\sigma `$ appears to be more convenient than a usage of metric $`\rho `$. The squared length $`\left\|M\left(𝒫^n\right)\right\|^2=F_n\left(𝒫^n\right)`$ is calculated for the proper Euclidean space, but the expression (1.9) - (1.11) may be used for any finite subspaces of any metric space, because it contains only world function $`\sigma `$ (metric $`\rho )`$ and may be calculated for any metric space. ###### Definition 1.13 A description is called $`\sigma `$-immanent, if it does not contain any references to objects or concepts other than finite subspaces of the metric space and its metric. Prefix $`\sigma `$ in the term ”$`\sigma `$-immanent” associates with the world function $`\sigma `$. Concept of $`\sigma `$-immanent description is very important for modification of the proper Euclidean geometry. Considering the proper Euclidean geometry to be a standard geometry and defining a geometrical object there in a $`\sigma `$-immanent way, one can use this definition in any metric space. Note that definition of geometrical objects is a principal problem of the metric geometry, i.e. the geometry, generated by the metric space. The shortest (line), connecting two arbitrary points $`P,Q\mathrm{\Omega }`$ of the metric space $`\{\rho ,\mathrm{\Omega }\}`$, is the basic geometrical object which is constructed usually in the metric space . One can construct an angle, triangle, different polygons from segments of the shortest. Construction of two-dimensional and three-dimensional planes in the metric space is rather problematic. At any rate it is unclear how one could construct these planes, using the shortest as the main geometrical object. A possibility of the metric space description in terms of only the shortest is restricted. Although exhibiting ingenuity, such a description may be constructed. For instance, A.D. Alexandrov showed that internal geometry of two-dimensional boundaries of convex three-dimensional bodies may be represented in terms of metric . Apparently, without introducing geometric objects which are analogs of two-dimensional plane, the solution of similar problem for three-dimensional boundaries of four-dimensional bodies is very difficult. Note that constraints (1.5), (1.6), imposed on metric, are necessary only for constructing the shortest. The shortest, determined by two points $`P_1,P_2`$, may be replaced by the $`\sigma `$-immanent definition (1.7) of segment $`𝒯_{[P_1P_2]}`$, which coincides with the shortest in the metric space, described by the definition 1.1. This definition in itself does not need constraint (1.5), describing definiteness of the metric space, and constraint (1.6), describing one-dimensionality of the segment $`𝒯_{[P_1P_2]}.`$ If the metric is not restricted by constraint (1.6), the segment $`𝒯_{[P_1P_2]}`$ takes the shape of a hallow tube, reminding ellipsoid, described by definition 1.9. If the constraint (1.6) is strengthened ($``$ is replaced by $`<`$), the segment $`𝒯_{[P_1P_2]}`$ degenerates into two points $`P_1,P_2`$. The case of the one-dimensional shortest is intermediate between the two cases. In the case of the proper Euclidean space, considered to be a metric space, the first order NGO, defined by (1.7) is one-dimensional line. It is not clear whether one-dimensionality is a special property of the Euclidean geometry, or it is a property of any geometry in itself. We do not see, why one should insist on the one-dimensionality of the first order NGO $`𝒯_{[P_1P_2]}`$ in the case of an arbitrary modification of the proper Euclidean geometry. First, it is useful to consider the most general modification of the proper Euclidean geometry. Second, at the end of investigation, if it appears to be necessary, one can always reduce a degree of generalization, imposing additional constraints. In the proper Euclidean space the $`n`$-dimensional plane ($`n`$th order NGO) $`n=1,2,\mathrm{}`$ is defined as follows ###### Definition 1.14 The $`n`$th order metric subspace $`M\left(𝒫^n\right)`$ of unvanishing length $`\left\|M\left(𝒫^n\right)\right\|^2=F_n\left(𝒫^n\right)0`$ determines the $`n`$th order tube (the $`n`$th order NGO) $`𝒯\left(𝒫^n\right)`$ by means of the relation $$𝒯\left(𝒫^n\right)𝒯_{𝒫^n}=\left\{P_{n+1}\right\|F_{n+1}\left(𝒫^{n+1}\right)=0\},P_i\mathrm{\Omega },i=0,1\mathrm{}n+1,$$ (1.13) where the function $`F_n`$ is defined by the relations (1.8), (1.10) The $`n`$th order tube $`𝒯_{𝒫^n}`$ which is an analog of the $`n`$-dimensional Euclidean plane may be constructed in any metric space, as far as its definition 1.14 is $`\sigma `$-immanent. It may be defined also in the metric space with omitted constraints (1.5), (1.6), imposed usually on the metric. We shall refer to such a generalized metric space as the $`\sigma `$-space. The geometry, generated by the $`\sigma `$-space, will be referred to as T-geometry (tubular geometry). ###### Definition 1.15 $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$ is nonempty set $`\mathrm{\Omega }`$ of points $`P`$ with given on $`\mathrm{\Omega }\times \mathrm{\Omega }`$ real function $`\sigma `$ $$\sigma :\mathrm{\Omega }\times \mathrm{\Omega },\sigma (P,P)=0,\sigma (P,Q)=\sigma (Q,P)P,Q\mathrm{\Omega }.$$ (1.14) The function $`\sigma `$ is called world function, or $`\sigma `$-function. The metric $`\rho `$ may be introduced in the $`\sigma `$-space by means of the relation (1.11). If $`\sigma `$ is positive, metric $`\rho `$ is also positive, but if $`\sigma `$ is negative, the metric is imaginary. ###### Definition 1.16 . Nonempty subset $`\mathrm{\Omega }^{}\mathrm{\Omega }`$ of points of the $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$ with the world function $`\sigma ^{}=\sigma \|_{\mathrm{\Omega }^{}\times \mathrm{\Omega }^{}}`$, which is a contraction of $`\sigma `$ on $`\mathrm{\Omega }^{}\times \mathrm{\Omega }^{}`$, is called $`\sigma `$-subspace $`V^{}=\{\sigma ^{},\mathrm{\Omega }^{}\}`$ of $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$. Further the world function $`\sigma ^{}=\sigma \|_{\mathrm{\Omega }^{}\times \mathrm{\Omega }^{}}`$, which is a contraction of $`\sigma `$ will be designed by means of $`\sigma `$. Any $`\sigma `$-subspace of $`\sigma `$-space is a $`\sigma `$-space. ###### Definition 1.17 . $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$ is called isometrically embeddable in $`\sigma `$-space $`V^{}=\{\sigma ^{},\mathrm{\Omega }^{}\}`$, if there exists such a monomorphism $`f:\mathrm{\Omega }\mathrm{\Omega }^{}`$, that $`\sigma (P,Q)=\sigma ^{}(f(P),f(Q))`$, $`P,Q\mathrm{\Omega },f(P),f(Q)\mathrm{\Omega }^{}`$, Any $`\sigma `$-subspace $`V^{}`$ of $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$ is isometrically embeddable in it. ###### Definition 1.18 . Two $`\sigma `$-spaces $`V=\{\sigma ,\mathrm{\Omega }\}`$ and $`V^{}=\{\sigma ^{},\mathrm{\Omega }^{}\}`$ are called to be isometric (equivalent), if $`V`$ is isometrically embeddable in $`V^{}`$, and $`V^{}`$ is isometrically embeddable in $`V`$. ###### Definition 1.19 The $`\sigma `$-space $`M=\{\rho ,\mathrm{\Omega }\}`$ is called a finite $`\sigma `$-space, if the set $`\mathrm{\Omega }`$ contains a finite number of points. ###### Definition 1.20 . The $`\sigma `$-subspace $`M_n(𝒫^n)=\{\sigma ,𝒫^n\}`$of the $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$, consisting of $`n+1`$ points $`𝒫^n=\{P_0,P_1,\mathrm{},P_n\}`$ is called the $`n`$th order $`\sigma `$-subspace . All geometrical objects of T-geometry are obtained as follows. Geometrical objects of the proper Euclidean geometry are defined in the $`\sigma `$-immanent form. Then they may be considered to be definitions of corresponding geometrical objects in T-geometry. The world function $`\sigma `$ of the proper Euclidean space satisfies some $`\sigma `$-immanent relations, describing special properties of the proper Euclidean geometry. Metric side of these relations had been formulated and proved by Menger . Using our designations, we present this result in the form of theorem. ###### Theorem 1.1 The $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$ is isomerically embeddable in $`n`$-dimensional Euclidean space $`E_n`$, if and only if any $`(n+2)`$th order $`\sigma `$-subspace $`M(𝒫^{n+2})\mathrm{\Omega }`$ is isometrically embeddable in $`E_n`$. Unfortunately, the formulation of this theorem is not $`\sigma `$-immanent, as far as it contains a reference to $`n`$-dimensional Euclidean space $`E_n`$ which is not defined $`\sigma `$-immanently. A more constructive version of the $`\sigma `$-space Euclideness conditions is formulated in the form of the following theorem. ###### Theorem 1.2 The $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$ is the $`n`$-dimensional Euclidean space, if and only if the following three $`\sigma `$-immanent conditions are fulfilled. I. $$𝒫^n\mathrm{\Omega },F_n(𝒫^n)0,F_{n+1}(\mathrm{\Omega }^{n+2})=0,$$ (1.15) II. $$\sigma (P,Q)=\frac{1}{2}\underset{i,k=1}{\overset{n}{}}g^{ik}(𝒫^n)[\mathrm{\Gamma }(P_0,P_i,P)\mathrm{\Gamma }(P_0,P_i,Q)]$$ $$\times [\mathrm{\Gamma }(P_0,P_k,P)\mathrm{\Gamma }(P_0,P_k,Q)],P,Q\mathrm{\Omega }$$ (1.16) where $`\mathrm{\Gamma }(P_0,P_k,P)`$ are defined by the relations (1.10). The quantities $`g^{ik}(𝒫^n),`$ $`(i,k=1,2,\mathrm{}n)`$ are defined by the relations $$\underset{k=1}{\overset{n}{}}g_{ik}(𝒫^n)g^{kl}(𝒫^n)=\delta _i^l,i,l=1,2,\mathrm{}n$$ (1.17) where $$g_{ik}(𝒫^n)=\mathrm{\Gamma }(P_0,P_i,P_k),i,k=1,2,\mathrm{}n$$ (1.18) III. The relations $$\mathrm{\Gamma }(P_0,P_i,P)=x_i,x_i,i=1,2,\mathrm{}n,$$ (1.19) considered to be equations for determination of $`P\mathrm{\Omega }`$, have always one and only one solution. ###### Remark 1.1 For the Euclidean space to be the proper Euclidean the eigenvalues of the matrix $`g_{ik}(𝒫^n)=\mathrm{\Gamma }(P_0,P_i,P_k),i,k=1,2,\mathrm{}n`$ are to be of the same sign, otherwise the Euclidean space is pseudo-Euclidean. ###### Remark 1.2 The condition (1.15) is a corollary of condition (1.16). It is formulated as a separate condition in order to separate definition of dimension and that of the coordinate system. Let us note that all three conditions are written in $`\sigma `$-immanent form. Proof of this theorem can be found in . Now we consider how results of this theorem can be used for construction of conventional description of the proper Euclidean space in some rectilinear coordinate system, starting from an abstract $`\sigma `$-space, satisfying conditions I - III of the theorem. Let there be $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\},`$ and it is known that conditions I - III of the theorem are fulfilled. Then the $`\sigma `$-space $`V`$ is an Euclidean space, but the dimension $`n`$ of the space is unknown. To determine the dimension $`n`$, let us take two different points $`P_0,P_1\mathrm{\Omega },F_1(𝒫^1)=2\sigma (P_0,P_1)0`$. 1. Let us construct the first order tube $`𝒯\left(𝒫^1\right)`$. If $`𝒯\left(𝒫^1\right)=\mathrm{\Omega }`$, then dimension of the $`\sigma `$-space $`V`$ $`n=1`$. If $`\mathrm{\Omega }\backslash 𝒯\left(𝒫^1\right)\mathrm{},P_2\mathrm{\Omega },P_2𝒯\left(𝒫^1\right),`$ and hence, $`F_2(𝒫^2)0`$. 2. Let us construct the second order tube $`𝒯\left(𝒫^2\right)`$. If $`𝒯\left(𝒫^2\right)=\mathrm{\Omega }`$, then $`n=2`$, otherwise$`P_3\mathrm{\Omega },P_3𝒯\left(𝒫^2\right),`$ and hence, $`F_3(𝒫^3)0`$. 3. Let us construct the third order tube $`𝒯\left(𝒫^3\right)`$. If $`𝒯\left(𝒫^3\right)=\mathrm{\Omega }`$, then $`n=3`$, otherwise$`P_4\mathrm{\Omega },P_4𝒯\left(𝒫^3\right),`$ and hence, $`F_4(𝒫^4)0`$. 4. Etc. Continuing this process, one determines such $`n+1`$ points $`𝒫^n`$, that the condition $`𝒯\left(𝒫^n\right)=\mathrm{\Omega }`$ and, hence, conditions (1.15) are fulfilled. Then by means of relations $$x_i\left(P\right)=\mathrm{\Gamma }(P_0,P_i,P),i=1,2,\mathrm{}n,$$ (1.20) one attributes covariant coordinates $`x\left(P\right)=\left\{x_i(P)\right\},i=1,2,\mathrm{}n`$ to $`P\mathrm{\Omega }`$. Let $`x=x\left(P\right)^n`$ and $`x^{}=x\left(P^{}\right)^n.`$ Substituting $`\mathrm{\Gamma }(P_0,P_i,P)=x`$ and $`\mathrm{\Gamma }(P_0,P_i,P^{})=x_i^{}`$ in (1.16), one obtains the conventional expression for the world function of the Euclidean space in the rectilinear coordinate system $$\sigma (P,P^{})=\sigma _E(x,x^{})=\frac{1}{2}\underset{i,k=1}{\overset{n}{}}g^{ik}(𝒫^n)\left(x_ix_i^{}\right)\left(x_kx_k^{}\right)$$ (1.21) where $`g^{ik}(𝒫^n)`$, defined by relations (1.18) and (1.17), is the contravariant metric tensor in this coordinate system. Condition III of the theorem states that the mapping $$x:\mathrm{\Omega }^n$$ described by the relation (1.20) is a bijection, i.e. for $`y^n`$ there exists such one and only one point $`Q\mathrm{\Omega },`$ that $`y=x\left(Q\right)`$. Thus, on the base of the world function, given on abstract set $`\mathrm{\Omega }\times \mathrm{\Omega }`$, one can determine the dimension $`n`$ of the Euclidean space, construct rectilinear coordinate system with the metric tensor $`g_{ik}(𝒫^n)=\mathrm{\Gamma }(P_0,P_i,P_k),i,k=1,2,\mathrm{}n`$ and describe all geometrical objects which are determined in terms of coordinates. The Euclidean space and Euclidean geometry is described in terms and only in terms of world function (metric). Changing the world function, one obtains another $`\sigma `$-space and another (non-Euclidean) geometry. One should expect that another geometry is also described completely in terms of the world function. The properties of geometrical objects may appear other than the properties of these objects in the proper Euclidean geometry. For instance, in the Euclidean geometry $`𝒯_{P_0P_1}𝒯_{P_0P_1P_2},`$ i.e. the straight, passing through the points $`P_0`$ and $`P_1`$, belongs to any two-dimensional plane, passing through these points. To prove these statement, one needs to use the relations (1.16). In the case of non-Euclidean geometry the relation $`𝒯_{P_0P_1}𝒯_{P_0P_1P_2}`$ is invalid, in general. Another example. Two circle cylinders $`𝒞(P_0,P_1;P)`$ and $`𝒞(P_0,P_1^{};P),`$ $`P_1^{}𝒯_{[P_0P_1]},P_1^{}P_1,`$ $`P_1^{}P_0`$ coincide in the proper Euclidean geometry, but they are different geometrical objects in non-Euclidean geometry. In the proper Euclidean geometry there exists geometrical object called line. ###### Definition 1.21 The broken line $`𝒯_{\mathrm{br}}`$ is the set of connected straight segments $`𝒯_{[P_iP_{i+1}]}`$ $$𝒯_{\mathrm{br}}=\underset{i}{}𝒯_{[P_iP_{i+1}]}$$ (1.22) The continuous line (or curve) is defined as a limit of the broken line $`𝒯_{\mathrm{br}}`$ at $`P_iP_{i+1},`$ $`(i=0,\pm 1,\pm 2,\mathrm{})`$. The smooth line is defined as a limit of (1.22) at $`P_iP_{i+1},`$ $`(i=0,\pm 1,\pm 2,\mathrm{})`$ under the constraint that $`\mathrm{cos}\mathrm{}P_{i1}P_iP_{i+1}1`$. Defining the segment $`𝒯_{[P_iP_{i+1}]}`$ by means of definition 1.10, one obtains metric definition of broken line (1.22). To obtain metric definition of the continuous line and that of smooth line, one needs to go to corresponding limits in (1.22). According to this definition the line is many-point geometrical object. This object is very complicated, because their points are given independently (i.e. there are many degrees of freedom). On the other hand, in the proper Euclidean geometry there exists another (nonmetric) definition of continuous line. The continuous line $``$ is defined as a continuous mapping $$L:I\mathrm{\Omega },I=[0,1].$$ (1.23) Strictly, the geometrical object is a set $`=L(I)\mathrm{\Omega }`$ of points of the $`\sigma `$-space $`V=\{\sigma ,\mathrm{\Omega }\}`$, but not the mapping (1.23) in itself. However, as far as the number set $`I`$ is fixed and the same in all cases, then with some stipulations one can consider the correspondence between the mapping $`L`$ and the set of images $`=L(I)`$ to be one-to-one. Then one can label the geometrical objects (considered as $`\sigma `$-subspaces) by means of mappings (1.23) and identify the mapping (1.23) with the geometrical object $``$, called curve (line). In the proper Euclidean geometry the definition of line (1.23) agrees with the definition 1.22. But in non-Euclidean geometry definitions (1.23) and (1.22) do not agree, in general. Already in the Riemannian geometry an application of definition (1.23) as one of basic definitions poses problems. In the Riemannian space the world function $`\sigma _R(x,x^{})`$ between the points $`x`$ and $`x^{}`$ is determined by the relation $$\sigma _R(x,x^{})=\frac{1}{2}\left(\underset{_{_{[xx^{}]}}}{}\sqrt{g_{ik}dx^idx^k}\right)^2$$ (1.24) where $`_{[xx^{}]}`$ denotes segment of geodesic connecting points $`x`$ and $`x^{}`$. Let us use the world function (1.24) instead of the Euclidean world function (1.21) in the $`\sigma `$-immanent description of geometry. In other words, let us use for construction of geometry the logical scheme (1.2), but not (1.1). One obtains the $`\sigma `$-Riemannian geometry which is expected to be equivalent to the Riemannian geometry, because both the Riemannian geometry and the $`\sigma `$-Riemannian one are two generalizations of the Euclidean geometry, using the same world function which has to describe any geometry completely. In reality, using for geometry construction different logical schemes, the $`\sigma `$-Riemannian geometry and the Riemannian one coincide, but not at all points,. The point is that the world function is a fundamental object of the $`\sigma `$-Riemannian geometry, whereas it is a derivative object in the Riemannian geometry, where the infinitesimal distance and the curve (line) are fundamental objects. The line $``$, defined by nonmetric definition (1.23), is a complicated and fundamental structure of Riemannian geometry, which is absent in such a form in the $`\sigma `$-Riemannian geometry. The continuous line $``$ in the $`\sigma `$-Riemannian geometry may be defined as a limit of the broken tube (1.22). But it is a derivative (not fundamental) geometrical object. As a whole the situation looks as follows. The $`\sigma `$-Riemannian geometry is constructed $`\sigma `$-immanently, i.e. on the base of metric and does not need the nonmetric definition of line (1.23). The Riemannian geometry is constructed on the base of infinitesimal metric $`dS=\sqrt{g_{ik}dx^idx^k}`$ (which coincide with the infinitesimal metric of the $`\sigma `$-Riemannian geometry) and uses the nonmetric definition of line (1.23) for definition of finite metric. As a result the finite metric of both geometries coincide, but only in the whole domain $`D=\mathrm{\Omega },`$ where both geometries are defined. If one considers $`\sigma `$-Riemannian and Riemannian geometries in some subdomain $`D^{}D`$, the finite metrics are defined in $`D^{}`$ in different ways for these geometries. For $`\sigma `$-Riemannian geometry the finite metric in $`D^{}`$ is defined as a cotraction of the finite metric in $`D`$, whereas for Riemannian geometry the finite metric is defined on the basis of system of geodesics inside $`D^{}`$ which does not coincide, in general, with the system of geodesics in $`D`$. The geodesic segment $`_{[xx^{}]}`$ which determines $`\sigma _R(x,x^{})`$ is a lengthy geometrical object, depending on the shape of the region $`D^{}`$, where the Riemannian geometry is defined. As a result the finite metrics of both geometries may be different in $`D^{}D`$, although they coincide in $`D`$. Note that the nonmetric definition of line (1.23) needs additional constraints to be rather definite. Let us discuss these problems. ## 2 Riemannian space and convexity problem The Riemannian space and the Riemannian geometry are introduced as follows. $`n`$-dimensional Riemannian space can be derived as a result of a generalization of the $`n`$-dimensional proper Euclidean space, written in a covariant form. Indeed, the $`n`$-dimensional Euclidean space $`E_n=\{𝐠_E,K,^n\}`$ is described by the infinitesimal distance written in the rectilinear coordinate system $`K`$ $$dS^2=g_{ik}dx^idx^k,g_{ik}=\mathrm{diag}\{1,1,\mathrm{}1\}$$ (2.1) $`𝐠_E`$ denotes the matrix $`g_{ik}=\mathrm{diag}\{1,1,\mathrm{}1\}`$ of the metric tensor. In the arbitrary curvilinear coordinate system $`\stackrel{~}{K}`$ the same distance have the form $$dS^2=\stackrel{~}{g}_{ik}\left(\stackrel{~}{x}\right)d\stackrel{~}{x}^id\stackrel{~}{x}^k,det\stackrel{~}{g}_{ik}0,$$ (2.2) Here $`\stackrel{~}{g}_{ik}\left(x\right)`$ is constrained by the relation $$\stackrel{~}{g}_{ik}\left(x\right)=\underset{l=1}{\overset{n}{}}\frac{f_l\left(x\right)}{x^i}\frac{f_l\left(x\right)}{x^k},$$ (2.3) where $`f_l:^n`$, $`l=1,2,\mathrm{}n`$ are $`n`$ functions restricted by one condition $`detf_i/x^k0,i,k=1,2,\mathrm{}n`$. If $`\stackrel{~}{g}_{ik}\left(x\right)`$ does not satisfy the relation (2.3) the space stops to be Euclidean and becomes a Riemannian space $`R_n=\{\stackrel{~}{𝐠},\stackrel{~}{K},^n\}`$. Constraint (2.3) is a condition of the Euclideness of the space. Eliminating (2.3) one obtains a Riemannian space $`R_n=\{𝐠,K,^n\}`$, which is determined by the form of the metric tensor $`g_{ik}(x)`$. The world function is determined by the relation (1.24), where $`_{[xx^{}]}^n`$ is the geodesic segment of the geodesic $`_{xx^{}}^n`$. This geodesic is an extremal of (1.24), considered as a functional of the curve $`:x=x(\tau )`$, written in the form $$\sigma [x(\tau )]=\frac{1}{2}\left(\underset{_{}}{}\sqrt{g_{ik}(x)\frac{dx^i}{d\tau }\frac{dx^k}{d\tau }}𝑑\tau \right)^2$$ (2.4) The geodesic $`_{xx^{}}:x=x(\tau )`$ is described by the equations $$_{xx^{}}:\frac{d^2x^i}{d\tau ^2}+\gamma _{kl}^i\frac{dx^k}{d\tau }\frac{dx^l}{d\tau }=0,i=1,2,\mathrm{}n$$ (2.5) where $$\gamma _{kl}^i=\gamma _{kl}^i(x)=\frac{1}{2}g^{ij}\left(g_{kj,l}+g_{lj,k}g_{kl,j}\right)$$ (2.6) is the Christoffel symbol, and comma before index $`l`$ denotes differentiation with respect to $`x^l`$. In particular, if $`g_{ik}=`$const, $`i,k=1,2,\mathrm{}n`$, $`g=detg_{ik}0`$, the world function is described by the relation (1.21), and the Riemannian space $`R_n=\{𝐠,K,^n\}`$ is the Euclidean space. Let us consider now the Riemannian space $`R_n=\{𝐠_E,K,D\}`$, where $`D^n`$ is some region of the Euclidean space $`E_n=\{𝐠_E,K,^n\}.`$ If the region $`D`$ is convex, i.e. any segment $`_{[xx^{}]}`$ of the straight $`_{xx^{}}`$, connecting the points $`x,x^{}D`$ belongs to $`D`$ ($`_{[xx^{}]}D)`$, then the world function of the Riemannian space $`R_n=\{𝐠_E,K,D\}`$ has the form (1.21) and the Riemannian space $`R_n=\{𝐠_E,K,D\}`$ can be embedded isometrically into the Euclidean space $`E_n=\{𝐠_E,K,^n\}`$. If the region $`D`$ is noncovex, then the system of geodesics of $`R_n=\{𝐠_E,K,D\}`$ is not a system of straight lines, and the world function (1.24) is not described by the relation (1.21). Example. Let us consider two-dimensional proper Euclidean space, and rectilinear orthogonal coordinates on it. Let us consider the region $`D:`$ $`\left(x^1\right)^2+\left(x^2\right)^21.`$ Geodesics of the Riemannian space $`R_2^{}=\{𝐠_E,K,D\}`$ looks as it is shown in Figure 1. After cutting a hole in the Euclidean plane the shape and length of geodesic segment between the points $`P`$ and $`P^{}`$ changes. World function $`\sigma (P,P^{})`$ between the points $`P`$ and $`P^{}`$ changes, and the part $`R_2^{}=\{𝐠_E,K,D\}`$ of the Euclidean plane $`R_2=\{𝐠_E,K,^2\}`$ stops to be embeddable isometrically in $`R_2=\{𝐠_E,K,^2\}`$. It seems to be rather strange, when part of the Euclidean plane cannot be embedded isometrically in the plane. The problem of convexity is rather strong, and most of geometricians prefer to get around this problem, considering convex regions . In the T-geometry the convexity problem is absent. Indeed, according to definition 1.16 any subset of a $`\sigma `$-space is always embeddable isometrically into the $`\sigma `$-space. From viewpoint of T-geometry, cutting a hole in the Euclidean plane $`R_2=\{𝐠_E,K,^2\}`$, one does not change the system of geodesics (the first order NGOs), one cuts only holes in geodesics, making them discontinuous. Continuity is a property of coordinate systems, used in Riemannian geometry as the main tool of description. From viewpoint of T-geometry the convexity problem is a problem made artificially. Insisting on continuity of geodesics, one overestimates importance of the continuity for geometry and attributes the continuity of geodesics (the first order NGOs) to any Riemannian geometry, whereas the continuity of geodesics is a special property of the proper Euclidean geometry. ## 3 Riemannian geometry and one-dimensionality of the first order NGOs Let us consider the $`n`$-dimensional pseudo-Euclidean space $`E_n=\{𝐠_1,K,^n\}`$ of the index $`1`$, $`𝐠_1=`$diag$`\{1,1,1\mathrm{}1\}`$ to be a kind of $`n`$-dimensional Riemannian space<sup>3</sup><sup>3</sup>3The term ”Riemannian space” is considered to be a collective term with respect to concepts ”proper Riemannian” and ”pseudo-Riemannian”. Matrix $`𝐠`$ of the metric tensor has eigenvalues of the same sign in the case of proper Riemannian space and of different signs in the case of pseudo-Riemannian one.. The world function is defined by the relation (1.21) $$\sigma _1(x,x^{})=\frac{1}{2}\underset{i,k=1}{\overset{n}{}}g^{ik}\left(x_ix_i^{}\right)\left(x_kx_k^{}\right),g^{ik}=\mathrm{diag}\{1,1,1\mathrm{}1\}$$ (3.1) Geodesic $`_{yy^{}}`$ is a straight line, and it is considered in pseudo-Euclidean geometry to be the first order NGOs, determined by two points $`y`$ and $`y^{}`$ $$_{yy^{}}:x^i=(y^iy^i)\tau ,i=1,2,\mathrm{}n,\tau $$ (3.2) The geodesic $`_{yy^{}}`$ is called timelike, if $`\sigma _1(y,y^{})>0`$, and it is called spacelike if $`\sigma _1(y,y^{})<0`$. The geodesic $`_{yy^{}}`$ is called null, if $`\sigma _1(y,y^{})=0`$. The pseudo-Euclidean space $`E_n=\{𝐠_1,K,^n\}`$ generates the $`\sigma `$-space $`V=\{\sigma _1,^n\}`$, where the world function $`\sigma _1`$ is defined by the relation (3.1). The first order tube (NGO) $`𝒯(x,x^{})`$ in the $`\sigma `$-Riemannian space $`V=\{\sigma _1,^n\}`$ is defined by the relation (1.13) $$𝒯(x,x^{})𝒯_{xx^{}}=\left\{r\right\|F_2(x,x^{},r)=0\},\sigma _1(x,x^{})0,x,x^{},r^n,$$ (3.3) $$F_2(x,x^{},r)=\left\|\begin{array}{cc}(x_i^{}x_i)(x^ix^i)& (x_i^{}x_i)(r^ix^i)\\ (r_ix_i)(x^ix^i)& (r_ix_i)(r^ix^i)\end{array}\right\|$$ (3.4) Solution of equations (3.3), (3.4) gives the following result $$𝒯_{xx^{}}=\left\{r\right\|\underset{y^n}{}\underset{\tau }{}r=\left(x^{}x\right)\tau +yx\mathrm{\Gamma }(x,x^{},y)=0\mathrm{\Gamma }(x,y,y)=0\},$$ (3.5) $$x,x^{},y,r^n$$ where $`\mathrm{\Gamma }(x,x^{},y)=(x_i^{}x_i)(y^ix^i)`$ is the scalar product of vectors $`\stackrel{}{xy}`$ and $`\stackrel{}{xx^{}}`$ defined by the relation (1.10). In the case of timelike vector $`\stackrel{}{xx^{}}`$, when $`\sigma _1(x,x^{})>0`$, there is a unique null vector $`\stackrel{}{xy}=\stackrel{}{xx}=\stackrel{}{0}`$ which is orthogonal to the vector $`\stackrel{}{xx^{}}`$. In this case the ($`n1)`$-dimensional surface $`𝒯_{xx^{}}`$ degenerates into the one-dimensional straight $$𝒯_{xx^{}}=\left\{r\right\|\underset{\tau }{}r=\left(x^{}x\right)\tau \},\sigma _1(x,x^{})>0,x,x^{},r^n,$$ (3.6) Thus, for timelike vector $`\stackrel{}{xx^{}}`$ the first order tube $`𝒯_{xx^{}}`$ coincides with the geodesic $`_{xx^{}}`$. In the case of spacelike vector $`\stackrel{}{xx^{}}`$ the $`(n1)`$-dimensional tube $`𝒯_{xx^{}}`$ contains the one-dimensional geodesic $`_{xx^{}}`$ of the pseudo-Euclidean space $`E_n=\{𝐠_1,K,^n\}`$. This difference poses the question what is the reason of this difference and what of the two generalization of the proper Euclidean geometry is more reasonable. Note that four-dimensional pseudo-Euclidean geometry is used for description of the real space-time. One can try to resolve this problem from experimental viewpoint. Free classical particles are described by means of timelike straight lines. At this point the pseudo-Euclidean geometry and the $`\sigma `$-pseudo-Euclidean geometry (T-geometry) lead to the same result. The spacelike straights are believed to describe the particles moving with superlight speed (so-called taxyons). Experimental attempts of taxyons discovery were failed. Of course, trying to discover taxyons, one considered them to be described by spacelike straights. On the other hand, the physicists believe that all what can exist does exist and may be discovered. From this viewpoint the failure of discovery of taxyons in the form of spacelike line justifies in favour of taxyons in the form of three-dimensional surfaces. To interpret the structure of the set (3.5), describing the first order tube, let us take into account the zeroth order tube $`𝒯_x`$, determined by the point $`x`$ in the $`\sigma `$-pseudo-Euclidean space is the light cone with the vertex at the point $`x`$ (not the point $`x`$). Practically the first order tube consists of such sections of the light cones with their vertex $`y_{xx^{}}`$ that all vectors $`\stackrel{}{yr}`$ of these sections are orthogonal to the vector $`\stackrel{}{xx^{}}`$. In other words, the first order tube $`𝒯_{xx^{}}`$ consists of the zeroth order tubes $`𝒯_y`$ sections at $`y`$, orthogonal to $`\stackrel{}{xx^{}}`$, with $`y_{xx^{}}`$. For timelike $`\stackrel{}{xx^{}}`$ this section consists of one point, but for the spacelike $`\stackrel{}{xx^{}}`$ it is two-dimensional section of the light cone. ## 4 Collinearity in Riemannian and $`\sigma `$-Riemannian geometry Let us return to the Riemannian space $`R_n=\{𝐠,K,D\},D^n`$, which generates the world function $`\sigma (x,x^{})`$ defined by the relation (1.24). Then the $`\sigma `$-space $`V=\{\sigma ,D\}`$ appears. it will be referred to as $`\sigma `$-Riemannian space. We are going to compare concept of collinearity (parallelism) of two vectors in the two spaces. The world function $`\sigma =\sigma (x,x^{})`$ of both $`\sigma `$-Riemannian and Riemannian spaces satisfies the system of equations <sup>4</sup><sup>4</sup>4The paper is hardly available for English speaking reader. Survey of main results of in English may be found in . See also $$\begin{array}{cc}(1)\sigma _l\sigma ^{lj^{}}\sigma _j^{}=2\sigma & (4)det\sigma _{i\|\|k}0\\ (2)\sigma (x,x^{})=\sigma (x^{},x)& (5)det\sigma _{ik^{}}0\\ (3)\sigma (x,x)=0& (6)\sigma _{ikl}=0\end{array}$$ (4.1) where the following designations are used $$\sigma _i\frac{\sigma }{x^i},\sigma _i^{}\frac{\sigma }{x^i},\sigma _{ik^{}}\frac{^2\sigma }{x^ix^k},\sigma ^{ik^{}}\sigma _{lk^{}}=\delta _l^i$$ Here the primed index corresponds to the point $`x,`$ and unprimed index corresponds to the point $`x`$. Two parallel vertical strokes mean covariant derivative $`\stackrel{~}{}_i^x^{}`$ with respect to $`x^i`$ with the Christoffel symbol $$\mathrm{\Gamma }_{kl}^i\mathrm{\Gamma }_{kl}^i(x,x^{})\sigma ^{is^{}}\sigma _{kls^{}},\sigma _{kls^{}}\frac{^3\sigma }{x^kx^lx^s}$$ For instance, $$G_{ik}G_{ik}(x,x^{})\sigma _{i\|\|k}\frac{\sigma _i}{x^k}\mathrm{\Gamma }_{ik}^l(x,x^{})\sigma _l\frac{\sigma _i}{x^k}\sigma _{iks^{}}\sigma ^{ls^{}}\sigma _l$$ (4.2) $$G_{ik\|\|l}\frac{G_{ik}}{x^l}\sigma _{ils^{}}\sigma ^{js^{}}G_{jk}\sigma _{kls^{}}\sigma ^{js^{}}G_{ij}$$ Summation from $`1`$ to $`n`$ is produced over repeated indices. The covariant derivative $`\stackrel{~}{}_i^x^{}`$ with respect to $`x^i`$ with the Christoffel symbol $`\mathrm{\Gamma }_{kl}^i(x,x^{})`$ acts only on the point $`x`$ and on unprimed indices. It is called the tangent derivative, because it is a covariant derivative in the Euclidean space $`E_x^{}`$ which is tangent to the Riemannian space $`R_n`$ at the point $`x^{}`$. The covariant derivative $`\stackrel{~}{}_i^{}^x`$ with respect to $`x^i`$ with the Christoffel symbol $`\mathrm{\Gamma }_{k^{}l^{}}^i^{}(x,x^{})`$ acts only on the point $`x^{}`$ and on primed indices. It is a covariant derivative in the Euclidean space $`E_x`$ which is tangent to the $`\sigma `$-Riemannian space $`V`$ at the point $`x`$ . In general, the world function $`\sigma `$ carries out the geodesic mapping $`G_x^{}:`$ $`R_nE_x^{}`$ of the Riemannian space $`R_n=\{𝐠,K,D\}`$ on the Euclidean space $`E_x^{}=\{𝐠,K_x^{},D\}`$, tangent to $`R_n=\{𝐠,K,D\}`$ at the point $`x^{}`$ . This mapping transforms the coordinate system $`K`$ in $`R_n`$ into the coordinate system $`K_x^{}`$ in $`E_x^{}`$. The mapping is geodesic in the sense that it conserves the lengths of segments of all geodesics, passing through the tangent point $`x^{}`$ and angles between them at this point. The tensor $`G_{ik}`$, defined by (4.2) is the metric tensor at the point $`x`$ in the tangent Euclidean space $`E_x^{}`$. The covariant derivatives $`\stackrel{~}{}_i^x^{}`$ and $`\stackrel{~}{}_k^x^{}`$ commute identically, i.e. $`(\stackrel{~}{}_i^x^{}\stackrel{~}{}_k^x^{}\stackrel{~}{}_k^x^{}\stackrel{~}{}_i^x^{})A_{ls}0,`$ for any tensor $`A_{ls}`$ . This shows that they are covariant derivatives in the flat space $`E_x^{}`$. The system of equations (4.1) contains only world function $`\sigma `$ and its derivatives, nevertheless the system of equations (4.1) is not $`\sigma `$-immanent, because it contains a reference to a coordinate system. It does not contain the metric tensor explicitly. Hence, it is valid for any Riemannian space $`R_n=\{𝐠,K,D\}`$. All relations written above are valid also for the $`\sigma `$-space $`V=\{\sigma ,D\}`$, provided the world function $`\sigma `$ is coupled with the metric tensor by relation (1.24). $`\sigma `$-immanent expression for scalar product $`\left(𝐏_0𝐏_1.𝐐_0𝐐_1\right)`$ of two vectors $`𝐏_0𝐏_1`$ and $`𝐐_0𝐐_1`$ in the proper Euclidean space has the form $$(𝐏_0𝐏_1.𝐐_0𝐐_1)\sigma (P_0,Q_1)+\sigma (Q_0,P_1)\sigma (P_0,Q_0)\sigma (P_1,Q_1)$$ (4.3) This relation can be easily proved as follows. In the proper Euclidean space three vectors $`𝐏_0𝐏_1`$, $`𝐏_0𝐐_1`$, and $`𝐏_1𝐐_1`$ are coupled by the relation $$𝐏_1𝐐_1^2=𝐏_0𝐐_1𝐏_0𝐏_1^2=𝐏_0𝐏_1^2+𝐏_0𝐐_1^22(𝐏_0𝐏_1.𝐏_0𝐐_1)$$ (4.4) where $`(𝐏_0𝐏_1.𝐏_0𝐐_1)`$ denotes the scalar product of two vectors $`𝐏_0𝐏_1`$ and $`𝐏_0𝐐_1`$ in the proper Euclidean space. It follows from (4.4) $$(𝐏_0𝐏_1.𝐏_0𝐐_1)=\frac{1}{2}\{𝐏_0𝐐_1^2+𝐏_0𝐏_1^2𝐏_1𝐐_1^2\}$$ (4.5) Substituting the point $`Q_1`$ by $`Q_0`$ in (4.5), one obtains $$(𝐏_0𝐏_1.𝐏_0𝐐_0)=\frac{1}{2}\{𝐏_0𝐐_0^2+𝐏_0𝐏_1^2𝐏_1𝐐_0^2\}$$ (4.6) Subtracting (4.6) from (4.5) and using the properties of the scalar product in the Euclidean space, one obtains $$(𝐏_0𝐏_1.𝐐_0𝐐_1)=\frac{1}{2}\{𝐏_0𝐐_1^2+𝐐_0𝐏_1^2𝐏_0𝐐_0^2𝐏_1𝐐_1^2\}$$ (4.7) Taking into account that $`𝐏_0𝐐_1^2=2\sigma (P_0,Q_1)`$, one obtains the relation (4.3) from the relation (4.7). Two vectors $`𝐏_0𝐏_1`$ and $`𝐐_0𝐐_1`$ are collinear $`𝐏_0𝐏_1\|\|𝐐_0𝐐_1`$ (parallel or antiparallel), provided $`\mathrm{cos}^2\theta =1,`$ where $`\theta `$ is the angle between the vectors $`𝐏_0𝐏_1`$ and $`𝐐_0𝐐_1`$. Taking into account that $$\mathrm{cos}^2\theta =\frac{(𝐏_0𝐏_1.𝐐_0𝐐_1)^2}{(𝐏_0𝐏_1.𝐏_0𝐏_1)(𝐐_0𝐐_1.𝐐_0𝐐_1)}=\frac{(𝐏_0𝐏_1.𝐐_0𝐐_1)^2}{\|𝐏_0𝐏_1\|^2\|𝐐_0𝐐_1\|^2}$$ (4.8) one obtains the following $`\sigma `$-immanent condition of the two vectors collinearity $$𝐏_0𝐏_1\|\|𝐐_0𝐐_1:(𝐏_0𝐏_1.𝐐_0𝐐_1)^2=\|𝐏_0𝐏_1\|^2\|𝐐_0𝐐_1\|^2$$ (4.9) The collinearity condition (4.9) is $`\sigma `$-immanent, because by means of (4.3) it can be written in terms of the $`\sigma `$-function only. Thus, this relation describes the vectors collinearity in the case of arbitrary $`\sigma `$-space. Let us describe this relation for the case of $`\sigma `$-Riemannian geometry. Let coordinates of the points $`P_0,P_1,Q_0,Q_1`$ be respectively $`x,`$ $`x+dx,`$ $`x^{}`$ and $`x^{}+dx^{}`$. Then writing (4.3) and expanding it over $`dx`$ and $`dx^{}`$, one obtains $`\left(𝐏_0𝐏_1.𝐐_0𝐐_1\right)`$ $``$ $`\sigma (x,x^{}+dx^{})+\sigma (x^{},x+dx)\sigma (x,x^{})\sigma (x+dx,x^{}+dx^{})=`$ $`\sigma +\sigma _l^{}dx^l^{}+{\displaystyle \frac{1}{2}}\sigma _{l^{},s^{}}dx^l^{}dx^s^{}+\sigma +\sigma _idx^i+{\displaystyle \frac{1}{2}}\sigma _{i,k}dx^idx^k\sigma `$ $`\sigma \sigma _idx^i\sigma _l^{}dx^l^{}{\displaystyle \frac{1}{2}}\sigma _{i,k}dx^idx^k\sigma _{i,l^{}}dx^idx^l^{}{\displaystyle \frac{1}{2}}\sigma _{l^{},s^{}}dx^l^{}dx^s^{}`$ $$(𝐏_0𝐏_1.𝐐_0𝐐_1)=\sigma _{i,l^{}}dx^idx^l^{}=\sigma _{il^{}}dx^idx^l^{}$$ (4.10) Here comma means differentiation. For instance, $`\sigma _{i,k}\sigma _i/x^k`$. One obtains for $`\|𝐏_0𝐏_1\|^2`$ and $`\|𝐐_0𝐐_1\|^2`$ $$\|𝐏_0𝐏_1\|^2=g_{ik}dx^idx^k,\|𝐐_0𝐐_1\|^2=g_{l^{}s^{}}dx^l^{}dx^s^{}$$ (4.11) where $`g_{ik}=g_{ik}(x)`$ and $`g_{l^{}s^{}}=g_{l^{}s^{}}(x^{})`$. Then the collinearity condition (4.9) is written in the form $$\left(\sigma _{il^{}}\sigma _{ks^{}}g_{ik}g_{l^{}s^{}}\right)dx^idx^kdx^l^{}dx^s^{}=0$$ (4.12) Let us take into account that in the Riemannian space the metric tensor $`g_{l^{}s^{}}`$ at the point $`x^{}`$ can be expressed via the world function $`\sigma `$ of points $`x,x^{}`$ by means of the relation $$g_{l^{}s^{}}=\sigma _{il^{}}G^{ik}\sigma _{ks^{}},g^{l^{}s^{}}=\sigma ^{il^{}}G_{ik}\sigma ^{ks^{}}$$ (4.13) where the tensor $`G_{ik}`$ is defined by the relation (4.2), and $`G^{ik}`$ is defined by the relation $$G^{il}G_{lk}=\delta _k^i$$ (4.14) Substituting the first relation (4.2) in (4.12) and using designation $$u_i=\sigma _{il^{}}dx^l^{},u^i=G^{ik}u_k=\sigma ^{il^{}}g_{l^{}s^{}}dx^{l^{}s^{}}$$ (4.15) one obtains $$\left(\delta _i^l\delta _k^sg_{ik}G^{ls}\right)u_lu_sdx^idx^k=0$$ (4.16) The vector $`u_i`$ is the vector $`dx_i^{}^{}=g_{i^{}k^{}}dx^k^{}`$ transported parallelly from the point $`x^{}`$ to the point $`x`$ in the Euclidean space $`E_x^{}`$ tangent to the Riemannian space $`R_n`$. Indeed, $$u_i=\sigma _{il^{}}g^{l^{}s^{}}dx_s^{},\stackrel{~}{}_k^x^{}\left(\sigma _{il^{}}g^{l^{}s^{}}\right)0,i,k=1,2,\mathrm{}n$$ (4.17) and tensor $`\sigma _{il^{}}g^{l^{}s^{}}`$ is the operator of the parallel transport in $`E_x^{}`$, because $$\left[\sigma _{il^{}}g^{l^{}s^{}}\right]_{x=x^{}}=\delta _i^{}^s^{}$$ and the tangent derivative of this operator is equal to zero identically. For the same reason, i.e. because of $$\left[\sigma ^{il^{}}g_{l^{}s^{}}\sigma ^{ks^{}}\right]_{x=x^{}}=g^{i^{}k^{}},\stackrel{~}{}_s^x^{}(\sigma ^{il^{}}g_{l^{}s^{}}\sigma ^{ks^{}})0$$ $`G^{ik}=\sigma ^{il^{}}g_{l^{}s^{}}\sigma ^{ks^{}}`$ is the contravariant metric tensor in $`E_x^{}`$, at the point $`x`$. The relation (4.16) contains vectors at the point $`x`$ only . At fixed $`u_i=\sigma _{il^{}}dx^l^{}`$ it describes a collinearity cone, i.e. a cone of infinitesimal vectors $`dx^i`$ at the point $`x`$ parallel to the vector $`dx^i^{}`$ at the point $`x^{}`$. Under some condition the collinearity cone can degenerates into a line. In this case there is only one direction, parallel to the fixed vector $`u^i`$. Let us investigate, when this situation takes place. At the point $`x`$ two metric tensors $`g_{ik}`$ and $`G_{ik}`$ are connected by the relation $$G_{ik}(x,x^{})=g_{ik}(x)+\underset{x}{\overset{x^{}}{}}F_{ikj^{\prime \prime }s^{\prime \prime }}(x,x^{\prime \prime })\sigma ^{j^{\prime \prime }}(x,x^{\prime \prime })𝑑x_{}^{\prime \prime }{}_{}{}^{s^{\prime \prime }},$$ (4.18) where according to $$\sigma ^i^{}=\sigma ^{li^{}}\sigma _l=G^{l^{}i^{}}\sigma _l^{}=g^{l^{}i^{}}\sigma _l^{}$$ (4.19) Integration does not depend on the path, because it is produced in the Euclidean space $`E_x^{}`$. The two-point tensor $`F_{ilk^{}j^{}}=F_{ilk^{}j^{}}(x,x^{})`$ is the two-point curvature tensor, defined by the relation $$F_{ilk^{}j^{}}=\sigma _{ilj^{}k^{}}=\sigma _{ilj^{},k^{}}\sigma _{sj^{}k^{}}\sigma ^{sm^{}}\sigma _{ilm^{}}=\sigma _{il\|\|k^{}\|\|j^{}}$$ (4.20) where one vertical stroke denotes usual covariant derivative and two vertical strokes denote tangent derivative. The two-point curvature tensor $`F_{ilk^{}j^{}}`$ has the following symmetry properties $$F_{ilk^{}j^{}}=F_{lik^{}j^{}}=F_{ilj^{}k^{}},F_{ilk^{}j^{}}(x,x^{})=F_{k^{}j^{}il}(x^{},x)$$ (4.21) It is connected with the one-point Riemann-Ghristoffel curvature tensor $`r_{iljk}`$ by means of relations $$r_{iljk}=\left[F_{ikj^{}l^{}}F_{ijk^{}l^{}}\right]_{x^{}=x}=f_{ikjl}f_{ijkl},f_{iklj}=\left[F_{ikj^{}l^{}}\right]_{x^{}=x}$$ (4.22) In the Euclidean space the two-point curvature tensor $`F_{ilk^{}j^{}}`$ vanishes as well as the Riemann-Ghristoffel curvature tensor $`r_{iljk}`$. Let us introduce designation $$\mathrm{\Delta }_{ik}=\mathrm{\Delta }_{ik}(x,x^{})=\underset{x}{\overset{x^{}}{}}F_{ikj^{\prime \prime }s^{\prime \prime }}(x,x^{\prime \prime })\sigma ^{j^{\prime \prime }}(x,x^{\prime \prime })𝑑x_{}^{\prime \prime }{}_{}{}^{s^{\prime \prime }}$$ (4.23) and choose the geodesic $`_{xx^{}}`$ as the path of integration. It is described by the relation $$\sigma _i(x,x^{\prime \prime })=\tau \sigma _i(x,x^{})$$ (4.24) which determines $`x^{\prime \prime }`$ as a function of parameter $`\tau `$. Differentiating with respect to $`\tau `$, one obtains $$\sigma _{ik^{\prime \prime }}(x,x^{\prime \prime })dx^{\prime \prime k^{\prime \prime }}=\sigma _i(x,x^{})d\tau $$ (4.25) Resolving equations (4.25) with respect to $`dx^{\prime \prime }`$ and substituting in (4.23), one obtains $$\mathrm{\Delta }_{ik}(x,x^{})=\sigma _l(x,x^{})\sigma _p(x,x^{})\underset{0}{\overset{1}{}}F_{ikj^{\prime \prime }s^{\prime \prime }}(x,x^{\prime \prime })\sigma ^{lj^{\prime \prime }}(x,x^{\prime \prime })\sigma ^{ps^{\prime \prime }}(x,x^{\prime \prime })\tau 𝑑\tau $$ (4.26) where $`x^{\prime \prime }`$ is determined from (4.24) as a function of $`\tau `$. Let us set $$F_{ik}^{..lp}(x,x^{})=F_{ikj^{}s^{}}(x,x^{})\sigma ^{lj^{}}(x,x^{})\sigma ^{ps^{}}(x,x^{})$$ (4.27) then $$G_{ik}(x,x^{})=g_{ik}(x)+\mathrm{\Delta }_{ik}(x,x^{})$$ (4.28) $$\mathrm{\Delta }_{ik}(x,x^{})=\sigma _l(x,x^{})\sigma _p(x,x^{})\underset{0}{\overset{1}{}}F_{ik}^{..lp}(x,x^{\prime \prime })\tau 𝑑\tau $$ (4.29) Substituting $`g_{ik}`$ from (4.28) in (4.16), one obtains $$\left(\delta _i^l\delta _k^sG^{ls}\left(G_{ik}\mathrm{\Delta }_{ik}\right)\right)u_lu_sdx^idx^k=0$$ (4.30) Let us look for solutions of equation in the form of expansion $$dx^i=\alpha u^i+v^i,G_{ik}u^iv^k=0$$ (4.31) Substituting (4.31) in (4.30), one obtains equation for $`v^i`$ $$G_{ls}u^lu^s\left[G_{ik}v^iv^k\mathrm{\Delta }_{ik}\left(\alpha u^i+v^i\right)\left(\alpha u^k+v^k\right)\right]=0$$ (4.32) If the $`\sigma `$-Riemannian space $`V=\{\sigma ,D\}`$ is $`\sigma `$-Euclidean, then as it follows from (4.29) $`\mathrm{\Delta }_{ik}=0`$. If $`V=\{\sigma ,D\}`$ is the proper $`\sigma `$-Euclidean space, $`G_{ls}u^lu^s0`$, and one obtains two equations for determination of $`v^i`$ $$G_{ik}v^iv^k=0,G_{ik}u^iv^k=0$$ (4.33) The only solution $$v^i=0,dx^i=\alpha u^i,i=1,2,\mathrm{}n$$ (4.34) of (4.32) is a solution of the equation (4.30), where $`\alpha `$ is an arbitrary constant. In the proper Euclidean geometry the collinearity cone always degenerates into a line. Let now the space $`V=\{\sigma ,D\}`$ be the $`\sigma `$-pseudo-Euclidean space of index $`1`$, and the vector $`u^i`$ be timelike, i.e. $`G_{ik}u^iu^k>0`$. Then equations (4.33) also have the solution (4.34). If the vector $`u^i`$ is spacelike, $`G_{ik}u^iu^k<0,`$ then two equations (4.33) have non-trivial solution, and the collinearity cone does not degenerate into a line. The collinearity cone is a section of the light cone $`G_{ik}v^iv^k=0`$ by the plane $`G_{ik}u^iv^k=0`$. If the vector $`u^i`$ is null, $`G_{ik}u^iu^k=0,`$ then equation (4.32) reduces to the form $$G_{ik}u^iu^k=0,G_{ik}u^iv^k=0$$ (4.35) In this case (4.34) is a solution, but besides there are spacelike vectors $`v^i`$ which are orthogonal to null vector $`u^i`$ and the collinearity cone does not degenerate into a line. In the case of the proper $`\sigma `$-Riemannian space $`G_{ik}u^iu^k>0,`$ and equation (4.32) reduces to the form $$G_{ik}v^iv^k\mathrm{\Delta }_{ik}\left(\alpha u^i+v^i\right)\left(\alpha u^k+v^k\right)=0$$ (4.36) In this case $`\mathrm{\Delta }_{ik}0`$ in general, and the collinearity cone does not degenerate. $`\mathrm{\Delta }_{ik}`$ depends on the curvature an on the distance between the points $`x`$ and $`x^{}`$. The more space curvature and the distance $`\rho (x,x^{}),`$ the more the collinearity cone aperture. In the curved proper $`\sigma `$-Riemannian space there is an interesting special case, when the collinearity cone degenerates . In any $`\sigma `$-Riemannian space the following equality takes place $$G_{ik}\sigma ^k=g_{ik}\sigma ^k,\sigma ^kg^{kl}\sigma _l$$ (4.37) Then it follows from (4.28) that $$\mathrm{\Delta }_{ik}\sigma ^k=0$$ (4.38) It means that in the case, when the vector $`u^i`$ is directed along the geodesic, connecting points $`x`$ and $`x^{}`$, i.e. $`u^i=\beta \sigma ^i`$, the equation (4.36) reduces to the form $$\left(G_{ik}\mathrm{\Delta }_{ik}\right)v^iv^k=0,u^i=\beta \sigma ^i$$ (4.39) If $`\mathrm{\Delta }_{ik}`$ is small enough as compared with $`G_{ik},`$ then eigenvalues of the matrix $`G_{ik}\mathrm{\Delta }_{ik}`$ have the same sign, as those of the matrix $`G_{ik}.`$ In this case equation (4.39) has the only solution (4.34), and the collinearity cone degenerates. ## 5 Discussion Thus, we see that in the $`\sigma `$-Riemannian geometry at the point $`x`$ there are many vectors parallel to given vector at the point $`x^{}`$. This set of parallel vectors is described by the collinearity cone. Degeneration of the collinearity cone into a line, when there is only one direction, parallel to the given direction, is an exception rather than a rule, although in the proper Euclidean geometry this degeneration takes place always. Nonuniformity of space destroys the collinearity cone degeneration. In the proper Riemannian geometry, where the world function satisfies the system (4.1), one succeeded in conserving this degeneration for direction along the geodesic, connecting points $`x`$ and $`x^{}`$. This circumstance is very important for degeneration of the first order NGOs into geodesic, because degeneration of NGOs is connected closely with the collinearity cone degeneration. Indeed, definition of the first order tube (1.13), or (3.3) may be written also in the form $$𝒯\left(𝒫^1\right)𝒯_{P_0P_1}=\left\{R\|𝐏_0𝐏_1\|\|𝐏_0𝐑\right\},P_0,P_1,R\mathrm{\Omega },$$ (5.1) where collinearity $`𝐏_0𝐏_1\|\|𝐏_0𝐑`$ of two vectors $`𝐏_0𝐏_1`$ and $`𝐏_0𝐑`$ is defined by the $`\sigma `$-immanent relation (4.9), which can be written in the form $$𝐏_0𝐏_1\|\|𝐏_0𝐑:F_2(P_0,P_1,R)=\left\|\begin{array}{cc}\left(𝐏_0𝐏_1.𝐏_0𝐏_1\right)& \left(𝐏_0𝐏_1.𝐏_0𝐑\right)\\ \left(𝐏_0𝐑.𝐏_0𝐏_1\right)& \left(𝐏_0𝐑.𝐏_0𝐑\right)\end{array}\right\|=0$$ (5.2) The form (5.1) of the first order tube definition allows one to define the first order tube $`𝒯(P_0,P_1;Q_0)`$, passing through the point $`Q_0`$ collinear to the given vector $`𝐏_0𝐏_1`$. This definition has the $`\sigma `$-immanent form $$𝒯(P_0,P_1;Q_0)=\left\{R\|𝐏_0𝐏_1\|\|𝐐_0𝐑\right\},P_0,P_1,Q_0,R\mathrm{\Omega },$$ (5.3) where collinearity $`𝐏_0𝐏_1\|\|𝐐_0𝐑`$ of two vectors $`𝐏_0𝐏_1`$ and $`𝐐_0𝐑`$ is defined by the $`\sigma `$-immanent relations (4.9), (4.7). In the proper Euclidean space the tube (5.3) degenerates into the straight line, passing through the point $`Q_0`$ collinear to the given vector $`𝐏_0𝐏_1`$. Let us define the set $`\omega _{Q_0}=\{𝐐_0𝐐)\|Q\mathrm{\Omega }\}`$ of vectors $`𝐐_0𝐐`$. Then $$𝒞(P_0,P_1;Q_0)=\left\{𝐐_0𝐐\right\|Q𝒯(P_0,P_1;Q_0)\}\omega _{Q_0}$$ (5.4) is the collinearity cone of vectors $`𝐐_0𝐐`$ collinear to vector $`𝐏_0𝐏_1`$. Thus, the one-dimensionality of the first order tubes and the collinearity cone degeneration are connected phenomena. In the Riemannian geometry the very special property of the proper Euclidean geometry (the collinearity cone degeneration) is considered to be a property of any geometry and extended to the case of Riemannian geometry. The line $``$, defined as a continuous mapping (1.23) is considered to be the most important geometrical object. This object is considered to be more important, than the metric, and metric in the Riemannian geometry is defined in terms of the shortest lines. Use of line as a basic concept of geometry is inadequate for description of geometry and poses problems, which appears to be artificial. For instance, the convexity problem, when elimination of part of the point set $`\mathrm{\Omega }`$ generates variation of properties of other regions is a result of the metric definition via concept of the line. Although choosing the world function in the proper way (satisfying equations (4.1)), one succeeded in conserving the collinearity cone degeneration for geodesic lines, but for distant points $`x`$ and $`x^{}`$ the collinearity cone does not degenerate, and the absolute parallelism is absent in the Riemannian geometry. Instead of the cone of collinear vectors one introduces concept of parallel transport of a vector, where the result depends on the path of the transport. Practically, it means that one vector of the vector cone is chosen and it is attributed to some curve connecting the points $`x`$ and $`x^{}`$. Being a special case of T-geometry, the $`\sigma `$-Riemannian geometry does not use the nonmetric concept of line at all. Here the nonmetric line is a special geometrical object characteristic for the proper Euclidean geometry which is a result of the collinearity cone degeneration. Instead of the continuous mapping (1.23) one uses the mapping $$m_n:I_n\mathrm{\Omega },I_n=\{0,1,\mathrm{}n\}$$ (5.5) which determines geometrical object $`m_n`$, called the $`n`$th order multivector. . The $`n`$th order multivector may be considered to be some generalization of the $`n`$th order $`\sigma `$-subspace $`M(𝒫^n)`$, and definition (5.5) of multivector appears to be $`\sigma `$-immanent. Application of mappings (5.5) is sufficient for description of any geometry, because all geometric objects are determined as subsets of the space $`\mathrm{\Omega }`$ (not as mappings). Use of such complicated mappings as (1.23) is not necessary. For instance, to investigate the properties of the first order tube $`𝒯_{P_0P_1}\mathrm{\Omega }`$ (geodesic), one needs to investigate the set $`𝒯=\left\{P_0\right\}\left\{P_1\right\}𝒯_{P_0P_1}\mathrm{\Omega }^3,`$ satisfying the condition $`F_3(𝒯)=0.`$ Here the mapping $`F_3`$ is known and fixed. Only zeros of the function $`F_3`$, having the form $`𝒯=\left\{P_0\right\}\left\{P_1\right\}𝒯_{P_0P_1}`$, are investigated. Power of the set $`𝒯`$ is much less than the power of the set of all mappings (1.23), and investigation of $`𝒯`$ is not so complicated as investigation of mappings (1.23). One can reduce the power of the set of all mappings (1.23), imposing some additional restrictions on mapping (1.23), but nothing can change the fact that the mapping (1.23) is an attribute of the proper Euclidean geometry and is not an attribute of a geometry in itself. The convexity problem confirms this. The real space-time may appear not to have property of the collinearity cone degeneration . Insisting on the mapping (1.23) as the main tool of geometry investigation, one closes the door for real investigations of geometry and shows a wrong way for them. Besides purely logical arguments in favour of the T-geometry approach there are arguments of applied character. The fact is that application of T-geometry to the space-time model construction leads to new encouraging results . Consideration of uniform isotropic continuous model with zeroth curvature leads to a class of models, distinguishing by the shape of the tube. This class contains the well known Minkowski model, for which the timelike tubes degenerate into lines and which is not optimal, because it does not enable to describe quantum phenomena without using the quantum principles. Other (nondegenerate) models of this class have the following properties: (1) geometrization of mass of a particle described by the broken tube (1.22), (2) stochasticity of the world tube of a free particle which is conditioned by the collinearity cone non-degeneracy. It turns out that it is possible one to choose optimal space-time model, for which the statistical description of stochastic free particle tubes coincides with the quantum description in terms of the Schrödinger equation. The quantum constant $`\mathrm{}`$ appears to be a space-time property, introducing some ”elementary length” (it is connected with the thickness of the particle world tube). As a result one does not need the quantum principles, and the quantum theory looks as a conception, created for compensation of our incorrect ideas on the space-time geometry at small distances.
warning/0002/math0002040.html	ar5iv	text	# The LMO-invariant of 3-manifolds of rank one and the Alexander polynomial ## Introduction In analogy with the theory of Vassiliev invariants of links, different notions of finite type invariants of $`3`$-manifolds have been introduced. For integral homology spheres these different notions coincide with the original definition of Ohtsuki (\[Oht\]). The LMO-invariant $`Z^{LMO}`$ assembles all $``$-valued finite type invariants of integral homology spheres in a formal series and is therefore called a universal finite type invariant (\[LMO\], \[Le1\]). For connected closed manifolds $`M`$ the following is known about $`Z^{LMO}`$: $$\begin{array}{cc}& \\ \text{For }M\text{ with …}\hfill & Z^{LMO}\text{ is determined by and determines … }\hfill \\ & \\ & \\ \mathrm{rank}H_1(M)=0\hfill & \text{all }\text{-valued invariants of Goussarov and Habiro (}\text{[Habi]}\text{)}\hfill \\ & \\ H_1(M)=\hfill & \text{the Alexander polynomial (Theorem 1 of }\text{[GaH]}\text{)}\hfill \\ & \\ \mathrm{rank}H_1(M)2\hfill & \text{the Casson-Walker-Lescop invariant (}\text{[HaB]}\text{}\text{[Hab]}\text{}\text{[Les]}\text{)}\hfill \end{array}$$ In this article we fill in the missing puzzle piece for the interpretation of the LMO-invariant of manifolds of rank $`1`$ in terms of classical invariants. We prove the following generalization of Theorem 1 of \[GaH\]. ###### Theorem 1. Let $`M`$ be a closed oriented $`3`$-manifold of rank $`1`$. Then the LMO-invariant $`Z^{LMO}(M)`$ is determined by the Alexander polynomial $`(M)`$, and conversely, $`(M)`$ is determined by $`Z^{LMO}(M)`$. In the proof of Theorem 1 of \[GaH\] it was used that the Alexander polynomial $``$ of links $`L`$ in $`S^3`$ can be obtained from the universal Vassiliev invariant $`Z`$ of links in $`S^3`$ via a map $`W_{}`$ as follows: (1) $$\frac{h}{e^{h/2}e^{h/2}}(L)_{\|t^{1/2}=e^{h/2}}=W_{}Z(L).$$ We generalize Equation (1) by replacing $`Z`$ by the Århus invariant Å (\[BGRT1\]) of knots in a rational homology sphere: ###### Theorem 2. Let $`K`$ be a null-homologous knot in a rational homology $`3`$-sphere. Then $$\frac{h}{e^{h/2}e^{h/2}}(K)_{\|t^{1/2}=e^{h/2}}=W_{}\text{Å}(K).$$ Theorem 2 is an important ingredient in the proof of Theorem 1. Theorems 1 and 2 will be proven in Section 4. In Sections 13 we prepare these proofs by recalling definitions and properties of the Alexander polynomial of links and manifolds, of unitrivalent diagrams and the map $`W_{}`$, and of the universal finite type invariants $`Z`$$`Z^{LMO}`$ and Å. ## 1. The Alexander polynomial In this section we make preliminary definitions and recall some facts about the Alexander polynomial $``$ from \[Les\]. All manifolds and submanifolds in this paper are oriented. Let $`M`$ be a rational homology $`3`$-sphere (meaning that $`M`$ is a connected closed manifold of dimension $`3`$ with $`H_1(M,)=0`$). Let $`K`$ be a knot in $`M`$. Choose a tubular neighborhood $`T`$ of $`K`$. A meridian of $`K`$ is a simple closed curve $`m`$ on the boundary $`T`$ of $`T`$ that is null-homologous in $`T`$. The curve $`m`$ is oriented by the right-hand rule. There exists a unique isomorphism $`i_K:H_1(MK,)`$ that sends a meridian of $`K`$ to $`1`$. As a $``$-linear map $`i_K`$ is uniquely determined by the property that for any oriented surface $`\mathrm{\Sigma }M`$ with $`\mathrm{\Sigma }K=\mathrm{}`$ the value $`i_K(\mathrm{\Sigma })`$ is the intersection number of $`K`$ with $`\mathrm{\Sigma }`$. For disjoint knots $`K_1,K_2`$ the linking number $`\mathrm{lk}(K_1,K_2)`$ is defined as $`i_{K_1}(K_2)`$. The linking number $`\mathrm{lk}(.,.)`$ is symmetric. Denote the number of components of a link $`L`$ by $`\|L\|`$. A framed link $`L`$ is a link with a simple closed curve $`\mu _i`$ on the boundary $`T_i`$ of a tubular neighborhood $`T_i`$ of each component $`K_i`$ ($`i=1,\mathrm{},\|L\|`$). Inside of $`T_i`$, $`\mu _i`$ is homologous to $`q_iK_i`$ for some $`q_i`$. The linking matrix $`(l_{ij})`$ of $`L`$ is defined by $`l_{ij}=\mathrm{lk}(K_i,\mu _j)/q_j`$. The link $`L`$ has integral framing if all $`q_i`$ are $`1`$. The values $`l_{ii}`$ are called framing of $`K_i`$. We denote by $`M_L`$ the manifold obtained by surgery on $`LM`$. Let $`LM`$ be a null-homologous link. Then there exists an oriented connected surface $`\mathrm{\Sigma }`$ embedded in $`M`$ such that $`\mathrm{\Sigma }=L`$. Any surface with this property is called a Seifert surface of $`L`$. Let $`\mathrm{\Sigma }^\pm =\mathrm{\Sigma }^+\mathrm{\Sigma }^{}`$ be a tubular neighborhood of $`\mathrm{\Sigma }`$ such that $`\mathrm{\Sigma }=\mathrm{\Sigma }^+\mathrm{\Sigma }^{}`$ and $`\mathrm{\Sigma }^+`$ lies on the positive side of $`\mathrm{\Sigma }`$. The Seifert form of $`\mathrm{\Sigma }M`$ is the $``$-bilinear form $`s:H_1(\mathrm{\Sigma })\times H_1(\mathrm{\Sigma })`$ defined by sending homology classes $`a`$, $`b`$ to $`\mathrm{lk}(A^{},B^+)`$ where $`A^{}`$ is a knot in $`\mathrm{\Sigma }^{}`$ representing $`a`$ and $`B^+`$ is a knot in $`\mathrm{\Sigma }^+`$ representing $`b`$. In this section a matrix of $`s`$ with respect to an arbitrary basis of $`H_1(\mathrm{\Sigma })`$ is called a Seifert matrix of $`\mathrm{\Sigma }`$ (later we will choose a particular basis of $`H_1(\mathrm{\Sigma })`$). Define the bilinear form $`s^{}`$ by $`s^{}(a,b)=s(b,a)`$. Then $`ss^{}`$ is the intersection form of $`\mathrm{\Sigma }`$. Denote the transpose of a matrix $`V`$ by $`V^{}`$. ###### Proposition 3. Let $`L`$ be a null-homologous link in a rational homology sphere $`M`$. Choose a Seifert surface $`\mathrm{\Sigma }`$ of $`L`$. Let $`V`$ be a Seifert matrix of $`\mathrm{\Sigma }`$. Then $$(L)=det(t^{1/2}Vt^{1/2}V^{})(t^{1/2}t^{1/2})^{\|L\|1}[(t^{1/2}t^{1/2})^2][t^{\pm 1/2}]$$ is an invariant of the pair $`LM`$ up to homeomorphism; in particular it is an isotopy invariant of $`L`$. Proposition 3 can be proven by using sign-determined Reidemeister torsion (see Proposition 2.3.13 of \[Les\], \[Tur\]). Up to sign the invariant $`(L)`$ can be described as follows. Let $`N`$ be a connected $`3`$-manifold and let $`\phi :H_1(N)𝒵=`$ be a homomorphism. Let $`\stackrel{~}{N}`$ be the connected cover of $`N`$ corresponding to $`\mathrm{Ker}(\phi )`$. Then $`H_1(\stackrel{~}{N})`$ is a module over the group ring $`[𝒵][t^{\pm 1}]`$. Let $`J[t^{\pm 1}]`$ be the order ideal of $`H_1(\stackrel{~}{N})`$. Let $`\mathrm{\Delta }_\phi (N)`$ be a generator of the smallest principal ideal containing $`J`$. Then $`\mathrm{\Delta }_\phi (N)`$ is unique up to multiplication by $`\pm t^i`$. For a link in a rational homology sphere $`M`$ we denote $`\mathrm{\Delta }_\phi (ML)`$ by $`\mathrm{\Delta }(L)`$, where $`\phi :H_1(ML)`$ is given by the sum of the linking numbers with the components of $`L`$. The following lemma (see Proposition 2.3.13 of \[Les\]) relates $`(L)`$ and $`\mathrm{\Delta }(L)`$. ###### Lemma 4. Let $`L`$ be a null-homologous link in a rational homology sphere $`M`$. Then there exists a unique $`i`$ such that $`t^{i/2}\mathrm{\Delta }(L)`$ is invariant under the replacement of $`t^{1/2}`$ by $`t^{1/2}`$. For some $`ϵ\{\pm 1\}`$ we have $`ϵt^{i/2}\mathrm{\Delta }(L)=\|H_1(M)\|(L)`$. Now consider a connected closed $`3`$-manifold $`N`$ of rank $`1`$. Denote the quotient of $`H_1(N)`$ by its torsion subgroup $`\mathrm{Tor}(H_1(N))`$ by $`H_1^\mathrm{\#}(N)`$. Choose an isomorphism $`\psi :H_1^\mathrm{\#}(N)`$. Denote the composition of the canonical projection $`H_1(N)H_1^\mathrm{\#}(N)`$ with $`\psi `$ by $`\overline{\psi }`$. The following two lemmas (see \[Les\], Section 5.1) allow to compare $`\mathrm{\Delta }_{\overline{\psi }}(N)`$ with a knot invariant. ###### Lemma 5. Every connected closed $`3`$-manifold $`N`$ of rank 1 can be obtained by $`0`$-framed surgery on a null-homologous knot $`K`$ in a rational homology sphere $`M`$. We then have $`\mathrm{Tor}(H_1(N))H_1(M)`$. ###### Lemma 6. Let $`K`$ be a null-homologous $`0`$-framed knot in a rational homology $`3`$-sphere $`M`$. Then $`\mathrm{\Delta }(K)`$ is equal to $`\mathrm{\Delta }_{\overline{\psi }}(M_K)`$ up to multiplication by $`\pm t^i`$. We see by Lemmas 4, 5 and 6 that there exists $`j`$ such that $`t^{j/2}\mathrm{\Delta }_{\overline{\psi }}(N)`$ is invariant under the replacement of $`t^{1/2}`$ by $`t^{1/2}`$. Furthermore, we can choose $`ϵ\{\pm 1\}`$ such that $`ϵ\mathrm{\Delta }_{\overline{\psi }}(N)_{\|t=1}=\|H_1(\mathrm{Tor}(N))\|>0`$. Denote $`(ϵt^{j/2}/\|H_1(\mathrm{Tor}(N))\|)\mathrm{\Delta }_{\overline{\psi }}(N)`$ by $`(N)`$. The definition of $`(N)`$ does not depend on the choice of the isomorphism $`\psi `$ because $`(N)[(t^{1/2}t^{1/2})^2]`$. The invariant $``$ satisfies (2) $$(K)=(M_K).$$ for all null-homologous $`0`$-framed knots $`K`$ in a rational homology sphere $`M`$. ## 2. Unitrivalent diagrams and $`W_{}`$ In this section we briefly recall facts about unitrivalent diagrams and use them to state properties of the Vassiliev invariants in the Alexander polynomial $``$. Let $`\mathrm{\Gamma }`$ be a compact oriented $`1`$-manifold whose boundary $`\mathrm{\Gamma }`$ is partitioned into two ordered sets called upper and lower boundary. Let $`X`$ be a set. A unitrivalent diagram with skeleton $`\mathrm{\Gamma }`$ is a graph $`D`$ with distinguished subgraph $`\mathrm{\Gamma }`$ such that all vertices of $`D`$ are either univalent or trivalent. Trivalent vertices not lying on $`\mathrm{\Gamma }`$ are called internal and are oriented by a cyclic order of the incident edges. Univalent vertices are also called legs. Each leg of a unitrivalent diagram is labeled by an element of $`X`$. We allow connected components in $`D`$ that do not intersect $`\mathrm{\Gamma }`$ whenever these components contain at least one trivalent vertex. Recall the definition of a $``$-vector space $`𝒜(\mathrm{\Gamma },X)`$ generated by unitrivalent diagrams modulo relations called $`(STU)`$$`(IHX)`$, and $`(AS)`$ (\[BN1\]). When $`\mathrm{\Gamma }`$ is equipped with additional information (for example: dots on circle-components of $`\mathrm{\Gamma }`$, a set $`Y`$ in bijection with circle-components of $`\mathrm{\Gamma }`$, a distinguished subset of the components of $`\mathrm{\Gamma }`$,…), we require in the definition of $`𝒜(\mathrm{\Gamma },X)`$ that homeomorphisms between unitrivalent diagrams also preserve this additional data. The space $`𝒜(\mathrm{\Gamma },X)`$ is graded by half of the number of vertices of unitrivalent diagrams. Denote $`𝒜(\mathrm{\Gamma },\mathrm{})`$ by $`𝒜(\mathrm{\Gamma })`$. The invariants of $`\mathrm{}`$-component links in $`S^3`$ that are coefficients of $`z^i=(t^{1/2}t^{1/2})^i`$ in $`(L)`$ induce linear forms $`W_i:𝒜_i^{\mathrm{}}`$ on the degree-$`i`$ part $`𝒜_i^{\mathrm{}}`$ of $`𝒜^{\mathrm{}}:=𝒜(S_{}^{1}{}_{}{}^{\mathrm{}})`$ (see Section 3 of \[BNG\]) . For $`a`$ in the completion of $`𝒜^{\mathrm{}}`$ by the degree, we define $$W_{}(a)=\underset{i}{}W_i(a)h^i[[h]].$$ It will follow from Theorem 2 and can also be seen directly that the Alexander polynomial of links in a rational homology sphere induces the same map $`W_{}`$ (see skein relation 2.3.16 of \[Les\], or Exercise 3.10 of \[BNG\]). The map $`W_{}`$ and its extensions to $`𝒜(\mathrm{\Gamma },X)`$ obtained from representations of the Lie superalgebra $`\mathrm{gl}(1\|1)`$ have the following property (see Proposition 7.1 of \[Vai\], consider the element of $`𝒜(\mathrm{})`$ seperately). ###### Lemma 7. Let $`D𝒜(\mathrm{\Gamma },X)`$ be a unitrivalent diagram. Assume that $`D`$ has an internal vertex $`u`$ such that all edges incident to $`u`$ are connected to internal vertices. Then we have $`W_{}(D)=0`$. Let $`I_xI:=[0,1]`$ ($`xX`$). Denote the disjoint union by $``$. For every partition of $`(\mathrm{\Gamma }_{xX}I_x)`$ into two ordered sets called upper and lower boundary there exists an isomorphism (3) $$\chi _X:𝒜(\mathrm{\Gamma },XY)𝒜(\mathrm{\Gamma }\underset{xX}{}I_x,Y)$$ given by the average over all permutations of putting $`x`$-labeled univalent vertices of a diagram on the corresponding skeleton component $`I_xI`$ of $`\mathrm{\Gamma }_{xX}I_x`$. The inverse of $`\chi _X`$ will be denoted by $`\sigma _X`$ and the set $`X`$ will not be specified when it is clear from the context. Obviously, there exists an isomorphism of $`𝒜(\mathrm{\Gamma }_{xX}I_x,Y)`$ with a space $`𝒜(\mathrm{\Gamma }_{xX}S_{x}^{1}{}_{}{}^{},Y)`$, where the circles $`S_{x}^{1}{}_{}{}^{}`$ have a dot and are in bijection with the set $`X`$. Similarly, we have a surjective map from $`𝒜(_{xX}I_x_{yY}I_y,Z)`$ to $`𝒜(_{xX}I_x_{yY}S_y^1,Z)`$ given by closing the intervals $`I_y`$ to form the circles $`S_y^1`$. Denote the composition of $`\chi _Y`$ with this surjective map by $`\overline{\chi }_Y`$. An important special case is $`𝒜(S^1)𝒜(S_{}^{1}{}_{}{}^{},\mathrm{})𝒜(I)=:𝒜`$ (see \[BN1\]). The space $`𝒜`$ is a commutative algebra with multiplication $`\mathrm{\#}`$ induced by the connected sum of the skeletons $`S^1`$ of diagrams (resp. by the concatenation of skeletons $`I`$ of diagrams). More generally, the connected sum of $`S^1`$ with any distinguished skeleton component $`C`$ of a unitrivalent diagram turns $`𝒜(\mathrm{\Gamma }C,X)`$ into an $`𝒜`$-module. Let $`\overline{𝒜}`$ be the quotient of $`𝒜`$ by the ideal generated by the element and let $`\pi :𝒜\overline{𝒜}`$ be the canonical projection. There exists a unique inclusion of algebras $`i:\overline{𝒜}𝒜`$ with the property that $`i(D)=D`$ for all diagrams $`D`$ such that $`DS^1`$ is connected and $`D`$ contains an internal vertex (\[BN1\], Equation (5), Exercise 3.16). The map $`P_{\mathrm{defr}}=i\pi :𝒜𝒜`$ is called deframing projection. The disjoint union of unitrivalent diagrams turns $`𝒜(\mathrm{},X)`$ into a commutative algebra and $`𝒜(\mathrm{\Gamma },X)`$ into an $`𝒜(\mathrm{},X)`$-module. Important examples of diagrams in $`𝒜(\mathrm{},X)`$ are so-called struts $`{}_{}{}^{i}_{}^{j}`$ with labels $`i,jX`$, and so-called wheels $`\omega _n=\text{}`$ having $`n`$ internal vertices lying on a circle and $`n`$ univalent vertices with the same label ($`n=4`$ in this example). Let $`𝒜(\mathrm{},X)_{\mathrm{strut}}𝒜(\mathrm{},X)`$ be the subalgebra generated by struts and $`𝒜(\mathrm{},X)_{\mathrm{wh}}`$ be the subalgebra generated by wheels. It is known that $`𝒜(\mathrm{},X)_{\mathrm{strut}}`$ is a polynomial algebra in the $`n(n+1)/2`$ different struts ($`n=\|X\|`$) and $`𝒜(\mathrm{},X)_{\mathrm{wh}}`$ is a polynomial algebra in wheels with an even number of univalent vertices. There exist unique projections from $`𝒜(\mathrm{\Gamma },X)`$ to $`𝒜(\mathrm{},X)_{\mathrm{strut}}`$ (resp. $`𝒜(\mathrm{},X)_{\mathrm{wh}}`$) that send all diagrams to $`0`$ that have a connected component that is not a strut (resp. a wheel). Define $`P_{\mathrm{strut}}:𝒜(\mathrm{\Gamma }_{xX}I_x,\mathrm{})𝒜(\mathrm{\Gamma },X)`$ as the composition of $`\sigma `$ with the projection to $`𝒜(\mathrm{},X)_{\mathrm{strut}}𝒜(\mathrm{\Gamma },X)`$. The map $`P_{\mathrm{strut}}`$ descends to $`𝒜(\mathrm{\Gamma }_{xX}S_x^1,\mathrm{})`$ where the circle-components $`S_x^1`$ are in bijection with $`X`$. Define $`P_{\mathrm{wh}}:𝒜𝒜(\mathrm{},\{x\})`$ as the composition of $`\sigma P_{\mathrm{defr}}`$ with the projection to $`𝒜(\mathrm{},\{x\})_{\mathrm{wh}}`$. We have $`P_{\mathrm{wh}}(a\mathrm{\#}b)=P_{\mathrm{wh}}(a)P_{\mathrm{wh}}(b)`$ for all $`a,b𝒜`$. The map $`P_{\mathrm{wh}}`$ is related to $`W_{}`$ as follows (see \[Vai\]\[Kri\]). ###### Lemma 8. For $`D𝒜`$ the value $`W_{}(D)`$ depends only on $`P_{\mathrm{wh}}(D)`$ and is determined by $$W_{}(D_1\mathrm{\#}D_2)=W_{}(D_1)W_{}(D_2)\text{and}W_{}(\overline{\chi }(\omega _{2n}))=2h^{2n}.$$ Lemma 8 was used in proofs of the Melvin-Morton-Rozansky conjecture (\[BNG\]). ## 3. Universal finite type invariants Recall from Section 3 of \[LM2\] that a non-associative framed tangle (or q-tangle) $`T`$ is a usual tangle with integral framing, except that $`\mathrm{source}(T)`$ and $`\mathrm{target}(T)`$ are equipped with parentheses on the sequences of $`\pm `$-symbols associated with the lower and upper boundary points of $`T`$. We denote by $`Z`$ the universal Vassiliev invariant of non-associative framed tangles (see \[LM2\]). Denote the underlying $`1`$-manifold of a tangle $`T`$ (together with the partition of $`T`$ into two ordered sets and possibly together with a decoration of $`T`$ such as dots, distinguished components, …) by $`\mathrm{\Gamma }(T)`$. Then the values $`Z(T)`$ lie in the completion of $`𝒜(\mathrm{\Gamma }(T))`$ by the degree. Let $`\nu =Z(O)`$ be the invariant of the trivial knot with $`0`$-framing. Let $`T=L^{}T^{\prime \prime }`$ be a diagram of a framed non-associative tangle where the components of the sublink $`L^{}`$ of $`T`$ are in bijection with a set $`X^{}`$ and each component of $`L^{}`$ has a dot on its circle. Define $`\stackrel{ˇ}{Z}(T)`$ as the connected sum of $`Z(T)`$ with $`\nu ^{\|L^{}\|}`$ along the components of $`\mathrm{\Gamma }(L^{})`$. Cut the chord diagrams in $`\stackrel{ˇ}{Z}(T)`$ at the dots and apply the isomorphism $`\sigma _X^{}`$. The result lies in the completion of $`𝒜(\mathrm{\Gamma }(T^{\prime \prime }),X^{})`$ and is called $`\stackrel{ˇ}{Z}^\sigma (T)`$. The value $`\stackrel{ˇ}{Z}^\sigma (T)`$ is not invariant under isotopies of the tangle represented by the diagram $`T`$. For tangles $`T`$ with dotted circles $`L^{}`$ invariants $`Z_0^{LMO}(T)`$ and $`\text{Å}_0(T)`$ of isotopy (that are also invariant under second Kirby moves along $`L^{}`$) are obtained from $`C=\stackrel{ˇ}{Z}^\sigma (T)`$ as follows (see \[LMO\], \[Le2\], \[BGRT2\]). Definition of $`Z_0^{LMO}`$: The degree-$`n`$ part of $`Z_0^{LMO}(T):=<C>`$ is obtained from the degree $`n+\|L^{}\|n`$ part of $`C`$ by forgetting the diagrams in $`C`$ that do not have exactly $`2n`$ legs of each color $`xX^{}`$, by summing over all the $`((2n1)!!)^{\|L^{}\|}=((2n)!/2^nn!)^{\|L^{}\|}`$ possible ways of gluing pairs of legs of diagrams in $`C`$ with the same label and by replacing circles that do not belong to $`\mathrm{\Gamma }(T^{\prime \prime })`$ by $`2n`$. Definition of $`\text{Å}_0`$: $`\text{Å}_0`$ is only defined when the linking matrix $`(l_{ij})`$ of $`L^{}`$ is invertible (or equivalently, when $`S_L^{}^3`$ is a rational homology sphere). Write $`C`$ in the form $$C=P\mathrm{exp}\left(\frac{1}{2}\underset{i,jX^{}}{}l_{ij}^i^j\right)$$ where $`P`$ contains no struts. Let $`(l^{ij})`$ be the inverse matrix of $`(l_{ij})`$. Then $$\text{Å}_0(T):=<P,\mathrm{exp}(\frac{1}{2}\underset{i,jX^{}}{}l_i^{ij}_j)>,$$ where $`<D_1,D_2>`$ is $`0`$ if for some $`i`$ the number of $`i`$-labeled legs of $`D_1`$ is not equal to the number of $`i`$-labeled legs of $`D_2`$, and is given by the sum of all ways of gluing all legs with $`i`$-labels to legs with $`i`$-labels in the remaining case. Let $`LM`$ be a link in a $`3`$-manifold. Represent $`LM`$ by a diagram $`L^{}L^{\prime \prime }`$ of a link in $`S^3`$, such that $`S_L^{}^3M`$ and the image of $`L^{\prime \prime }`$ in $`S_L^{}^3`$ is mapped to $`L`$ by this homeomorphism. Put a dot on each component of $`L^{}`$. Two invariants $`Z^{LMO}`$ and Å of homeomorphisms of the pair $`(M,L)`$ are obtained from $`Z_0^{LMO}(L^{}L^{\prime \prime })`$ and $`\text{Å}_0(L^{}L^{\prime \prime })`$ by normalization (making it invariant under the first Kirby move) as follows: (4) $`Z^{LMO}(L)`$ $`=`$ $`Z_0^{LMO}(U_+)^{\sigma _+}Z_0^{LMO}(U_{})^\sigma _{}Z_0^{LMO}(L^{}L^{\prime \prime }),`$ (5) $`\text{Å}(L)`$ $`=`$ $`\text{Å}_0(U_+)^{\sigma _+}\text{Å}_0(U_{})^\sigma _{}\text{Å}_0(L^{}L^{\prime \prime }),`$ where $`U_\pm `$ is the trivial knot with a dot and framing $`\pm 1`$ and $`\sigma _+`$ (resp. $`\sigma _{}`$) is the number of positive (resp. negative) eigenvalues of the linking matrix $`(l_{ij})`$ of $`L^{}`$. The invariants of the empty link $`Z^{LMO}(\mathrm{})`$ and $`\text{Å}(\mathrm{})`$ are also denoted by $`Z^{LMO}(M)`$ and $`\text{Å}(M)`$, respectively. The series $`\text{Å}_0(U_\pm )`$ have degree-$`0`$ term $`1`$. Therefore Lemma 7 implies (6) $$W_{}\text{Å}_0(L^{}L^{\prime \prime })=W_{}\text{Å}(L).$$ We will make use of the following result of \[BGRT3\] (Equation (7) follows from Proposition 1.2 of \[BGRT3\] in the same way as Theorem 1 of \[BGRT3\]): (7) $$\text{Å}(L)=\|H_1(M)\|^{deg}Z^{LMO}(L),$$ where $`\|H_1(M)\|^{deg}`$ denotes the operation of multiplying diagrams of degree $`m`$ by $`\|H_1(M)\|^m`$. Let us recall some notation used in Lemma 9 below. Let $`T`$ be a non-associative framed tangle $`T`$ with a distinguished subset $`\stackrel{~}{T}`$ of its components. Denote by $`d(T)`$ the non-associative framed tangle given by replacing each component in $`\stackrel{~}{T}`$ by two copies that are parallel with respect to the framing. The symbols $`a\{+,\}`$ in $`\mathrm{source}(T)`$ (resp. $`\mathrm{target}(T)`$) that belong to $`\stackrel{~}{T}`$ are replace by $`(aa)`$ in $`\mathrm{source}(d(T))`$ (resp. $`\mathrm{target}(d(T))`$). Define $`s(T)`$ by reversing the orientation of each component in $`\stackrel{~}{T}`$. Define $`ϵ(T)`$ by deleting $`\stackrel{~}{T}`$. Now let $`D`$ be a unitrivalent diagram $`D`$ with a distinguished subset $`\stackrel{~}{\mathrm{\Gamma }}`$ of its skeleton components. Define $`d(D)`$ by replacing each skeleton component in $`\stackrel{~}{\mathrm{\Gamma }}`$ by two copies, and by summing over all ways of lifting vertices of $`D`$ that lie on $`\stackrel{~}{\mathrm{\Gamma }}`$ to the new skeleton. Define $`s(D)`$ by reversing the orientation of the components in $`\stackrel{~}{\mathrm{\Gamma }}`$ and by multiplying with $`_{C\stackrel{~}{\mathrm{\Gamma }}}(1)^{n_C}`$ where $`n_C`$ is the number of vertices lying on the skeleton component $`C`$ of $`D`$. If $`n_C>0`$ for some component $`C`$ of $`\stackrel{~}{\mathrm{\Gamma }}`$, then define $`ϵ(D)=0`$. Define $`ϵ(D)`$ by deleting the components in $`\stackrel{~}{\mathrm{\Gamma }}`$ in the remaining case. The composition $`T_1T_2`$ of non-associative tangles $`T_1,T_2`$ with $`\mathrm{source}(T_1)=\mathrm{target}(T_2)`$ is defined by placing $`T_1`$ on the top of $`T_2`$. For diagrams $`D_i`$ in $`𝒜(\mathrm{\Gamma }(T_i))`$ a composition $`D_1D_2`$ is defined similarly. In the following lemma we state generalizations of well-known properties of $`Z`$. ###### Lemma 9. Let $`T`$, $`T_1`$, $`T_2`$ be non-associative tangles with dotted circles. (1) Assume that some of the components of $`T`$ without dots are distinguished. Then we have<sup>1</sup><sup>1</sup>1As in \[LM3\] we must assume for the first property of $`\text{Å}_0`$ that an even associator is used in the definition of $`Z`$. This causes no restrictions in Theorems 1 and 2 because for links $`L`$, the invariants $`\text{Å}(L)`$ and $`Z^{LMO}(L)`$ do not depend on the choice of an associator. $$d(\text{Å}_0(T))=\text{Å}_0(d(T)),s(\text{Å}_0(T))=\text{Å}_0(s(T)),ϵ(\text{Å}_0(T))=\text{Å}_0(ϵ(T)).$$ (2) Assume that $`\mathrm{source}(T_1)=\mathrm{target}(T_2)`$. Then $$\text{Å}_0(T_1T_2)=\text{Å}_0(T_1)\text{Å}_0(T_2).$$ (3) We have $$\text{Å}_0(T)=\overline{\chi }_Y((\mathrm{exp}(P)),$$ where $`P`$ is a series of connected diagrams in $`𝒜(\mathrm{},Y)`$ and $`Y`$ is a set in bijection with the components of $`T`$ without dots. The proof of Lemma 9 is straightforward. Statements similar to Lemma 9 hold for $`\text{Å}(L)`$. ## 4. Proofs of Theorems 1 and 2 Recall from Equation (1) that for links $`L`$ in $`S^3`$ we have $`c(L)_{\|t^{1/2}=e^{h/2}}=W_{}Z(L)`$ with $`c=h/(e^{h/2}e^{h/2})`$. Equation (1) is proven in \[LM1\] and \[BNG\] by showing that $`W_{}Z`$ satisfies a skein relation and $`W_{}Z(O)=c`$. With the methods of this proof one can show directly that $`W_{}\text{Å}`$ satisfies the same skein relation for links in a rational homology sphere and $`W_{}\text{Å}(O)=c`$, but this does not imply Theorem 2. In this section we present a proof of Theorem 2 based on Equation (1). Then we prove Theorem 1 by using Theorem 2. Let $`L=L^{}L^{\prime \prime }`$ be a framed link in a rational homology sphere $`M`$. Denote the components of $`L^{}`$ (resp. $`L^{\prime \prime }`$) by $`K_x`$ with $`xX^{}`$ (resp. $`xX^{\prime \prime }`$) and their framings by $`\mu _x`$. For $`x,yX^{}X^{\prime \prime }`$ let $`l_{xy}=\mathrm{lk}(\mu _x,K_y)`$ be linking numbers in $`M`$, let the submatrix corresponding to $`L^{}`$ be invertible and denote its inverse by $`(l^{xy})_{x,yX^{}}`$. In the following lemma we recall how the linking numbers transform under surgery. ###### Lemma 10. For $`i,jX^{\prime \prime }`$ the linking numbers $`\stackrel{~}{l}_{ij}=\mathrm{lk}(\mu _i,K_j)`$ of $`L^{\prime \prime }M_L^{}`$ are given by $$\stackrel{~}{l}_{ij}=l_{ij}\underset{x,yX^{}}{}l_{ix}l^{xy}l_{yj}.$$ ###### Proof. Denote the meridians of the components of $`L`$ by $`m_x`$. In $`H_1(M(L^{}L^{\prime \prime }),)`$ the framings $`\mu _y`$ can uniquely be expressed as $`\mu _y=_{jX^{}X^{\prime \prime }}l_{yj}m_j`$. This implies for $`xX^{}`$ that $$\underset{yX^{}}{}l^{xy}\mu _y=m_x+\underset{yX^{},jX^{\prime \prime }}{}l^{xy}l_{yj}m_j.$$ In $`H_1(M_L^{}L^{\prime \prime },)=H_1(M(L^{}L^{\prime \prime }),)/(\mu _x)_{xX^{}}`$ we obtain the following unique expression of $`\mu _i`$ ($`iX^{\prime \prime }`$) in terms of the meridians $`m_j`$ ($`jX^{\prime \prime }`$) of $`L^{\prime \prime }M_L^{}`$: $$\mu _i=\underset{jX^{}X^{\prime \prime }}{}l_{ij}m_j=\underset{jX^{\prime \prime }}{}l_{ij}m_j\underset{xX^{},yX^{},jX^{\prime \prime }}{}l_{ix}l^{xy}l_{yj}m_j.$$ This implies the lemma. ∎ The following lemma tells us that the linking numbers of a link $`LM`$ can be recovered from $`P_{\mathrm{strut}}(\text{Å}(L))`$. ###### Lemma 11. Let $`L`$ be a link with integral framing in a rational homology sphere $`M`$. Let the components of $`L`$ be in bijection with a set $`X`$. Let $`(\stackrel{~}{l}_{ij})_{i,jX}`$ be the linking matrix of $`L`$. Then $$P_{\mathrm{strut}}(\text{Å}(L))=\mathrm{exp}\left(\frac{1}{2}\underset{i,jX}{}\stackrel{~}{l}_{ij}^i^j\right).$$ ###### Proof. Choose a diagram of $`L^{}L^{\prime \prime }S^3`$ such that $`(S_L^{}^3,L^{\prime \prime })(M,L)`$ and put dots on the components of $`L^{}`$. Let $`(l_{xy})_{x,yX^{}X^{\prime \prime }}`$ be the linking matrix of $`L^{}L^{\prime \prime }`$ and let $`(l^{xy})_{x,yX^{}}`$ be the inverse of the linking matrix of $`L^{}`$. Then for a series $`P`$ (resp. $`\stackrel{~}{P}`$) of diagrams in $`𝒜(\mathrm{},X^{}X^{\prime \prime })`$ (resp. in $`𝒜(\mathrm{},X^{\prime \prime })`$) that contains no struts and has degree-$`0`$-term $`1`$, we have $`\stackrel{ˇ}{Z}^\sigma (L^{}L^{\prime \prime })`$ $`=`$ $`\overline{\chi }_{X^{\prime \prime }}(P\mathrm{exp}\left({\displaystyle \frac{1}{2}}{\displaystyle \underset{x,yX^{}X^{\prime \prime }}{}}l_{xy}^x^y\right))\text{and}`$ $`\text{Å}_0(L^{}L^{\prime \prime })`$ $`=`$ $`\overline{\chi }_{X^{\prime \prime }}(<P\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,jX^{\prime \prime }}{}}l_{ij}^i^j+{\displaystyle \underset{iX^{\prime \prime },xX^{}}{}}l_{ix}^i^x),`$ $`\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle \underset{x,yX^{}}{}}l_x^{xy}_y)>)`$ $`=`$ $`\overline{\chi }_{X^{\prime \prime }}(\stackrel{~}{P}\mathrm{exp}({\displaystyle \frac{1}{2}}{\displaystyle \underset{i,jX^{\prime \prime }}{}}l_{ij}^i^j{\displaystyle \frac{1}{2}}{\displaystyle \underset{i,jX^{\prime \prime },x,yX^{}}{}}l_{ix}l^{xy}l_{yj}^i^j)).`$ Since $`P_{\mathrm{strut}}(\text{Å}(L))=P_{\mathrm{strut}}(\text{Å}_0(L^{}L^{\prime \prime }))`$, Lemma 11 follows from Lemma 10. ∎ For technical reasons we fix a representative of each homeomorphism-class of connected compact surfaces with boundary. We call this representative $`\mathrm{\Sigma }`$ a standard surface and equip it with a decomposition into a single vertex $`vI\times I`$ (also called coupon) with bands $`B_iI\times I`$ that are glued along $`I\times \{0,1\}`$ to the upper boundary $`I\times \{1\}`$ of $`v`$. Call the part $`I\times \{0\}`$ of $`v`$ its distinguished lower boundary. We orient the cores $`I\times \{1/2\}`$ of the bands $`B_i`$ ($`i=1,\mathrm{},\mathrm{rank}H_1(\mathrm{\Sigma })`$). An example is shown on the left side of Figure 1. We associate a basis of $`H_1(\mathrm{\Sigma })`$ to the ribbon graph decomposition of $`\mathrm{\Sigma }`$ as shown in Figure 1 by an example. The orientation of $`b_i`$ is determined by the orientation of the core of the band $`B_i`$. An embedding of a standard surface into $`R^2\times I`$ is an example of a ribbon graph in the sense of Section 8 of \[KaT\]. Ribbon graphs without vertices can canonically be identified with framed tangles. We will use this identification in the following. From now on we use the term Seifert matrix of a Seifert surface $`\mathrm{\Sigma }M`$ always with respect to a basis of $`H_1(\mathrm{\Sigma })`$ obtained by identifying $`\mathrm{\Sigma }`$ with a standard surface in some freely chosen way. We use the same basis for a matrix of the intersection form of $`\mathrm{\Sigma }`$. ###### Keylemma 12. Let $`K`$ be a knot in a rational homology sphere $`M`$ bounding a Seifert surface $`\mathrm{\Sigma }`$. Let $`V`$ be a Seifert matrix of $`\mathrm{\Sigma }`$. Then the power series $`W_{}\text{Å}(K)`$ depends only on $`V`$. The coefficient of $`h^i`$ in this series is a polynomial in the entries of $`V`$. Let us prepare the proof of Keylemma 12. We will make some statements more generally for links instead of knots. Let $`V=(v_{ij})`$ be a Seifert matrix. Choose a null-homologous link $`L`$ in a rational homology sphere $`M`$ with Seifert surface $`\mathrm{\Sigma }`$ and Seifert matrix $`V`$. The homeomorphism type of $`\mathrm{\Sigma }`$ is determined by the similarity type of $`VV^{}`$. There exists a link with integral framing $`\stackrel{~}{L}M`$ such that $`M_{\stackrel{~}{L}}=S^3`$. The link $`\stackrel{~}{L}`$ can be chosen to be disjoint from $`\mathrm{\Sigma }`$ because changing crossings between $`L=\mathrm{\Sigma }`$ and $`\stackrel{~}{L}`$ preserves the property that $`M_{\stackrel{~}{L}}=S^3`$. Therefore $`\mathrm{\Sigma }M`$ can be obtained from a surface $`\mathrm{\Sigma }^{\prime \prime }S^3`$ by surgery along a link $`L^{}S^3\mathrm{\Sigma }^{\prime \prime }`$. The identification of $`\mathrm{\Sigma }`$ with a standard surface induces an identification of $`\mathrm{\Sigma }^{\prime \prime }`$ with a standard surface. In a diagram of $`L^{}\mathrm{\Sigma }^{\prime \prime }`$ the vertex $`v`$ of $`\mathrm{\Sigma }^{\prime \prime }`$ can be pulled downwards, such that the diagram $`L^{}\mathrm{\Sigma }^{\prime \prime }`$ is of the form $`(L^{}T_1^{\prime \prime })T_2^{\prime \prime }`$ where $`T_2^{\prime \prime }`$ is a planar diagram of a neighborhood of $`v`$ and the distinguished lower boundary of $`v`$ is the lowest part of the diagram. Put dots on the components of $`L^{}`$. An example is shown in Figure 2. The components of $`T_1^{\prime \prime }`$ are in bijection with the set $`X^{\prime \prime }=\{1,\mathrm{},\mathrm{rank}H_1(\mathrm{\Sigma })\}`$. Regard $`T_1^{\prime \prime }`$ as a non-associative framed tangle with parentheses of the form $`((((\mathrm{}).).).)`$ on $`\mathrm{source}(T_1^{\prime \prime })`$. Let $`F=(f_{ij})=VV^{}`$ be the matrix of the intersection form of $`\mathrm{\Sigma }`$ and let $`U=(u_{ij})=1/2(V+V^{})=V1/2F`$. ###### Lemma 13. With the notation from above we have $$P_{\mathrm{strut}}(\text{Å}_0(L^{}T_1^{\prime \prime }))=\mathrm{exp}\left(\frac{1}{2}\underset{i,jX^{\prime \prime }}{}u_{ij}^i^j\right).$$ ###### Proof. Let $`K_i^+`$ (resp. $`K_i^{}`$) be a knot in the upper part $`\mathrm{\Sigma }^+M`$ (resp. in the lower part $`\mathrm{\Sigma }^{}M`$) of a tubular neighborhood of $`\mathrm{\Sigma }`$ representing the $`i`$-th basis element of $`H_1(\mathrm{\Sigma })`$. Let the knot $`K_i^{}K_i^+`$ have the framing $`\mathrm{lk}(K_i^{},K_i^+)=v_{ii}`$ induced by the surface $`\mathrm{\Sigma }`$. First consider $`ijX^{\prime \prime }`$. Define $`P_{ij}`$ as the composition of $`P_{\mathrm{strut}}`$ with the projection to the part containing only powers of the strut $`{}_{}{}^{i}_{}^{j}`$. Lemma 11 implies that $`P_{ij}(\text{Å}(K_i^{}K_j^+))=\mathrm{exp}(v_{ij}^i^j)`$. Represent $`K_i^{}K_j^+M`$ by a surgery diagram $`(L^{}S_1^{\prime \prime })S_2^{\prime \prime }`$ where the tangle $`S_1^{\prime \prime }`$ consists of the $`i`$-th and $`j`$-th framed strands of $`T_1^{\prime \prime }`$ and $`S_2^{\prime \prime }`$ is a $`0`$-framed tangle consisting of two intervals close to $`T_2^{\prime \prime }`$. See Figure 3 for an example (compare Figures 1 and 2). In this figure the dotted line separates $`S_2^{\prime \prime }`$ from $`L^{}S_1^{\prime \prime }`$ and is not a part of the diagram. We have $`\text{Å}_0(S_2^{\prime \prime })=Z(S_2^{\prime \prime })`$ and the explicit description of $`Z`$ (see \[LM2\]) implies $$P_{\mathrm{strut}}(\text{Å}_0(S_2^{\prime \prime }))=\mathrm{exp}((1/2)f_{ij}^i^j).$$ Observe the following property of $`P_{ij}`$: $$P_{ij}\left(\text{Å}_0(L^{}S_1^{\prime \prime })\text{Å}_0(S_2^{\prime \prime })\right)=P_{ij}(\text{Å}_0(L^{}S_1^{\prime \prime }))P_{ij}(\text{Å}_0(S_2^{\prime \prime })).$$ The last two formulas and Part (2) of Lemma 9 imply $$P_{ij}(\text{Å}_0(L^{}S_1^{\prime \prime }))=\mathrm{exp}((v_{ij}f_{ij}/2)^i^j)=\mathrm{exp}(u_{ij}^i^j).$$ Using Lemma 9 for $`ϵ`$ we see that $`P_{ij}(\text{Å}(L^{}T_1^{\prime \prime }))=\mathrm{exp}(u_{ij}^i^j)`$. For $`i=j`$ Lemma 11 implies $`P_{ii}(\text{Å}(K_i^\pm ))=\mathrm{exp}((1/2)v_{ii}^i^i)=\mathrm{exp}((1/2)u_{ii}^i^i)`$. We apply Lemma 9 for $`ϵ`$ as above and obtain $`P_{ii}(\text{Å}_0(L^{}T_1^{\prime \prime }))=\mathrm{exp}((1/2)u_{ii}^i^i)`$. By Part (3) of Lemma 9 we have $`P_{\mathrm{strut}}(\text{Å}_0(L^{}T_1^{\prime \prime }))=_{ij}P_{ij}(\text{Å}_0(L^{}T_1^{\prime \prime }))`$ which completes the proof. ∎ Starting from $`V`$ we made a lot of choices in the definition of $`T_1^{\prime \prime }`$. Since $`\text{Å}(L)`$ is an invariant only the choice of $`LM`$ can influence $`W_{}\text{Å}(L)`$. Now we are ready to show that for knots $`L`$ the invariant $`W_{}\text{Å}(L)`$ depends only on $`V`$. ###### Proof of Keylemma 12. We use the notation from above. For suitable distinguished components of $`T_1^{\prime \prime }`$ and of $`d(T_1^{\prime \prime })`$ the tangle $`s(d(T_1^{\prime \prime }))`$ coincides with the part of the framed oriented boundary of $`\mathrm{\Sigma }^{\prime \prime }`$ that belongs to $`T_1^{\prime \prime }`$. Let $`T_3^{\prime \prime }`$ be the part of the framed oriented boundary of $`\mathrm{\Sigma }^{\prime \prime }`$ that belongs to $`T_2^{\prime \prime }`$. We regard $`T_3^{\prime \prime }`$ as a non-associative tangle with $`\mathrm{target}(T_3^{\prime \prime })=\mathrm{source}(s(d(T_1^{\prime \prime })))`$. The invariant $`Z(T_3^{\prime \prime })=\text{Å}_0(T_3^{\prime \prime })`$ depends only on $`\mathrm{rank}H_1(\mathrm{\Sigma })`$. Since we know that the Seifert matrix $`V`$ is chosen with respect to a basis induced by a standard surface, the definition of the map $`\chi _{X^{\prime \prime }}:𝒜(\mathrm{},X^{\prime \prime })𝒜(\mathrm{\Gamma }(T_1^{\prime \prime }),\mathrm{})`$ depends only on $`V`$ (see Equation (3)). We will show below that for knots $`L`$ all terms in $`\text{Å}_0(L^{}T_1^{\prime \prime })`$ that contain an internal vertex do not contribute to $`W_{}(\text{Å}(L))`$. Equation (6), Lemma 9 and Lemma 13 then imply $`W_{}(\text{Å}(L))`$ $`=`$ $`W_{}(\text{Å}_0(L^{}\mathrm{\Sigma }^{\prime \prime }))`$ $`=`$ $`W_{}(s(d(\text{Å}_0(L^{}T_1^{\prime \prime })))Z(T_3^{\prime \prime }))`$ $`=`$ $`W_{}\left(s(d(\chi _{X^{\prime \prime }}(P_{\mathrm{strut}}(\text{Å}_0(L^{}T_1^{\prime \prime })))))Z(T_3^{\prime \prime })\right)`$ $`=`$ $`W_{}(sd\chi _{X^{\prime \prime }}\left(\mathrm{exp}\left({\displaystyle \frac{1}{4}}{\displaystyle \underset{i,jX^{\prime \prime }}{}}(v_{ij}+v_{ji})^i^j\right)\right)Z(T_3^{\prime \prime })).`$ This will show that $`W_{}(\text{Å}(L))`$ is determined by the Seifert matrix $`V`$. Obviously, the coefficients of $`h^i`$ in $`W_{}(\text{Å}(L))`$ are polynomials of degree $`i`$ in the entries of $`V`$. This will prove the keylemma. It remains to consider diagrams $`D`$ in $`\text{Å}_0(L^{}T_1^{\prime \prime })`$ with an internal vertex $`u`$. In $`s(d(D))`$ each of the edges incident to $`u`$ is either connected to another internal vertex or appears twice, namely as the difference of the two ways of lifting it to the skeleton $`\mathrm{\Gamma }(s(d(T_1^{\prime \prime })))`$. We represent this difference by a box in Figure 4. A neighborhood of the internal vertex $`u`$ looks like in one of the possibilities (a)-(f) in Figure 4. When we push a lifted vertex in the box along the circle $`\mathrm{\Gamma }(\mathrm{\Sigma }^{\prime \prime })`$, then it will finally cancel with the second lifted vertex. By the (STU)-relation we can replace a box in Figure 4 by a sum of diagrams with an additional internal vertex. More precisely, a part of the diagram looking like in (a), (b), (c), or (d) in Figure 4 is replaced by a sum of diagrams where a neighborhood of $`u`$ looks like in diagrams that can be reached by following a directed arrow in Figure 4. When we apply this procedure to all boxes, we will finally end up with possibilities (e) and (f). By Lemma 7 all diagrams that have a subdiagram as in (e) or (f) are sent to $`0`$ by $`W_{}`$. ∎ Let us recall a fact about knots (and links) in $`S^3`$ (see Proposition 8.7 of \[BuZ\]). ###### Fact 14. Let $`V`$ be a $`n\times n`$-matrix over $``$ such that $`VV^{}`$ is a matrix of the intersection form of a surface. Then $`V`$ is a Seifert matrix of a link in $`S^3`$. Since for all Seifert forms $`s`$ the intersection form of $`\mathrm{\Sigma }`$ is equal to $`ss^{}`$, we see that Seifert forms of a fixed surface $`\mathrm{\Sigma }`$ are a subset of an affine space whose associated $``$-vector space are symmetric $``$-bilinear forms on $`H_1(\mathrm{\Sigma })`$ with values in $``$. By Fact 14 Seifert forms of Seifert surfaces $`\mathrm{\Sigma }`$ in $`S^3`$ are a lattice of full rank in this affine space. ###### Proof of Theorem 2. By Proposition 3 and Keylemma 12 the coefficients of $`h^i`$ in the two power series $`\frac{h}{e^{h/2}e^{h/2}}(K)_{\|t^{1/2}=e^{h/2}}`$ and $`W_{}\text{Å}(K)`$ only depend on a Seifert matrix $`V`$ of a knot $`K`$ and are polynomials $`p_i`$ and $`q_i`$ in the entries of $`V`$. By Equation (1) we have $`p_i(V)=q_i(V)`$ for all Seifert matrices of knots in $`S^3`$. Fact 14 implies that $`p_i=q_i`$ for all $`i`$. ∎ The following lemma is a straightforward extension of a result of \[GaH\]. ###### Lemma 15. Let $`K`$ be a $`0`$-framed knot in a rational homology sphere $`M`$. Then any of the series $`Z^{LMO}(M_K)`$, $`W_{}Z^{LMO}(K)`$, $`W_{}\text{Å}(K)`$ can be computed from any other of these series. ###### Sketch of proof. The invariants $`Z^{LMO}`$ and Å of $`KM`$ differ only by normalization (see Equation (7)). Let $`C=Z^{LMO}(K)`$. Then we have $`Z^{LMO}(M_K)=<\sigma (\nu \mathrm{\#}C)>`$ with $`\nu =Z(O)`$. The following four steps show that $`W_{}(C)`$ can be calculated from $`Z^{LMO}(M_K)`$ and vice versa. This will complete the proof. 1) $`W_{}(C)`$ depends only on the wheel-part $`P_{\mathrm{wh}}(C)`$ of $`C`$ (Lemma 8). 2) $`P_{\mathrm{wh}}(C)`$ can be calculated from $`W_{}(C)`$ because $`P_{\mathrm{wh}}(C)=\mathrm{exp}(P)`$ where $`P`$ is a formal series of connected wheels (see Part (3) of Lemma 9), $`W_{}(C)=\mathrm{exp}(W_{}(P))`$, and $`W_{}`$ is injective on connected wheels (Lemma 8). 3) $`\sigma (\nu \mathrm{\#}C)`$ contains no struts because $`K`$ is $`0`$-framed (see Lemma 11 and Equation (7)). All remaining non-vanishing diagrams in $`𝒜(\mathrm{},\{x\})`$ have at least as many internal vertices as univalent vertices. This implies that $`Z^{LMO}(M_K)`$ depends only on $`P_{\mathrm{wh}}(\nu \mathrm{\#}C)=P_{\mathrm{wh}}(\nu )P_{\mathrm{wh}}(C)`$. 4) The map $`<>`$ is injective on wheels (see \[GaH\], Lemma 3.1, use the $`\mathrm{sl}_2`$-weight system on $`𝒜(\mathrm{})`$ to see that $`<>`$ is injective on connected wheels). Therefore $`P_{\mathrm{wh}}(\nu )P_{\mathrm{wh}}(C)`$ can be calculated from $`Z^{LMO}(M_K)=<P_{\mathrm{wh}}(\nu )P_{\mathrm{wh}}(C)>`$. $`P_{\mathrm{wh}}(\nu )`$ is invertible. ∎ Now we prove the main result of this paper. ###### Proof of Theorem 1. By Lemmas 5 and 15 and by Equation (2) it is sufficient to show that for a null-homotopic knot $`K`$ in a rational homology sphere each of the invariants $`(K)`$ and $`W_{}\text{Å}(K)`$ can be computed from the other one. This statement follows from Theorem 2. ∎
warning/0002/cond-mat0002115.html	ar5iv	text	# Delocalization in the Anderson model due to a local measurement ## Abstract We study a one-dimensional Anderson model in which one site interacts with a detector monitoring the occupation of that site. We demonstrate that such an interaction, no matter how weak, leads to total delocalization of the Anderson model, and we discuss the experimental consequences. PACS: 03.65.Bz, 73.20.Fz, 73.20.Jc Consider an electron in a one-dimensional array of $`N`$ coupled wells. The system is described by the Anderson tunneling Hamiltonian $$H_A=\underset{j=1}{\overset{N}{}}E_jc_j^{}c_j+\underset{j=1}{\overset{N1}{}}(\mathrm{\Omega }_jc_{j+1}^{}c_j+H.c.),$$ (1) where the operator $`c_j^{}(c_j)`$ corresponds to the creation (annihilation) of an electron in the well $`j`$. We assume for simplicity that each of the wells contains one bound state $`E_j`$ and is coupled only to its nearest neighbors with couplings $`\mathrm{\Omega }_j`$ and $`\mathrm{\Omega }_{j1}`$. (We choose $`\mathrm{\Omega }_j`$ real without loss of generality.) The electron-wave function in this system can be written as $`\|\mathrm{\Psi }(t)=_jb_j(t)c_j^{}\|0`$, where $`b_j(t)`$ is the probability amplitude of finding the electron in the well $`j`$ at time $`t`$. These amplitudes are obtained from the time-dependent Schrödinger equation $`i_t\|\mathrm{\Psi }(t)=H_A\|\mathrm{\Psi }(t)`$. It is well known that for randomly distributed levels $`E_j`$ (or random couplings $`\mathrm{\Omega }_j`$) all electronic states in this structure are localized. Hence, if the electron initially occupies the first well, $`b_j(0)=\delta _{j1}`$, the probability of finding it in the last well, $`P_N(t)=\|b_N(t)\|^2`$ drops exponentially with $`N`$: $`P_N(t\mathrm{})_{ensemble}\mathrm{exp}(\alpha N)`$. Anderson localization is usually associated with destructive quantum-mechanical interference between different probability amplitudes $`b_j(t)`$. This interference, however, can be affected by measuring the electron’s position in the system due to interaction of the electron with a macroscopic detector. For instance, the continuous monitoring of one of the wells of a double-well system ($`N=2`$ in Eq. (1)) destroys the off-diagonal elements (coherences) of the electron density matrix. As a result, the latter become the statistical mixture: $`\sigma _{jj^{}}(t)(1/2)\delta _{jj^{}}`$ for $`t\mathrm{}`$. In the case of the $`N`$-well structure, however, the monitoring of one of the wells cannot determine the electron’s position in the entire system. One might suppose therefore that such a local measurement cannot totally destroy the interference inside the entire system and hence, the electron localization. We demonstrate in this letter the contrary: any interaction, no matter how weak, of the electron with a macroscopic detector placed on only one well leads to total delocalization of the electron state: $`P_N(t\mathrm{})=1/N`$, even when $`N\mathrm{}`$. As a physical realization we consider a mesoscopic system of coupled quantum dots (Fig. 1), where a point contact is placed near the first dot. The point contact is coupled to two reservoirs, emitter and collector, at different chemical potentials, $`\mu _L`$ and $`\mu _R`$. A current $`I=eT(\mu _L\mu _R)/(2\pi )`$ flows through the point contact, where $`T`$ is its transmission coefficient. If the electron occupies the first dot, the transmission coefficient of the point contact decreases, $`T^{}<T`$, due to the electrostatic repulsion generated by the electron. As a result, the current $`I^{}<I`$ (Fig. 1a). The current returns to its previous value $`I`$ whenever the electron occupies any other dot, since then it is far away from the contact (Fig. 1b). The entire system can be described by the tunneling Hamiltonian $`H=H_A+H_{PC}+H_{int}`$, where $`H_A`$ is given by Eq. (1) and $`H_{PC}`$ $`=`$ $`{\displaystyle \underset{l}{}}E_la_l^{}a_l+{\displaystyle \underset{r}{}}E_ra_r^{}a_r+{\displaystyle \underset{l,r}{}}(\mathrm{\Omega }_{lr}a_r^{}a_l+H.c.)`$ (2) $`H_{int}`$ $`=`$ $`{\displaystyle \underset{l,r}{}}\delta \mathrm{\Omega }_{lr}c_1^{}c_1(a_r^{}a_l+H.c.),`$ (3) where $`a_l^{}(a_l)`$ and $`a_r^{}(a_r)`$ are the creation (annihilation) operators in the left and the right reservoirs, and $`\mathrm{\Omega }_{lr}`$ is the hopping amplitude between the states $`l`$ and $`r`$ of the reservoirs. Consider an initial state where all the levels in the emitter and the collector are filled up to the Fermi energies $`\mu _L`$ and $`\mu _R`$, respectively, and the electron occupies the first well. The many-body wave function describing the entire system can be written in the occupation number representation as $`\|\mathrm{\Psi }(t)={\displaystyle }_{j=1}^N[b_j(t)c_j^{}+{\displaystyle \underset{l,r}{}}b_{jlr}(t)c_j^{}a_r^{}a_l`$ (4) $`+{\displaystyle \underset{l<l^{},r<r^{}}{}}b_{jll^{}rr^{}}(t)c_j^{}a_r^{}a_r^{}^{}a_la_l^{}+\mathrm{}]\|0,`$ (5) where $`b(t)`$ are the probability amplitudes of finding the system in the states defined by the corresponding creation and annihilation operators. Using these amplitudes one defines the reduced density matrices $`\sigma _{jj^{}}^{(m)}(t)`$ that describe the electron and the detector, $`\sigma _{jj^{}}^{(0)}(t)=b_j(t)b_j^{}^{}(t),\sigma _{jj^{}}^{(1)}(t)={\displaystyle \underset{l,r}{}}b_{jlr}(t)b_{j^{}lr}^{}(t),`$ (6) $`\sigma _{jj^{}}^{(2)}(t)={\displaystyle \underset{ll^{},rr^{}}{}}b_{jll^{}rr^{}}(t)b_{j^{}ll^{}rr^{}}^{}(t),\mathrm{}`$ (7) Here $`j,j^{}=\{1,2,\mathrm{},N\}`$ denote the occupation states of the $`N`$-dot system. The index $`m`$ denotes the number of electrons that have reached the right-hand reservoir by time $`t`$. The total probability for the electron to occupy the dot $`j`$ is $`\sigma _{jj}(t)=_m\sigma _{jj}^{(m)}(t)`$. The off-diagonal density-matrix element $`\sigma _{jj^{}}(t)=_m\sigma _{jj^{}}^{(m)}(t)`$ describes interference between the states $`E_j`$ and $`E_j^{}`$. In order to find the amplitudes $`b(t)`$, we substitute Eq. (5) into the time-dependent Schrödinger equation $`i_t\|\mathrm{\Psi }(t)=H\|\mathrm{\Psi }(t)`$, and use the Laplace transform $`\stackrel{~}{b}(E)=_0^{\mathrm{}}b(t)\mathrm{exp}(iEt)𝑑t`$. Then we find an infinite set of algebraic equations for the amplitudes $`\stackrel{~}{b}(E)`$, given by $`(EE_1)\stackrel{~}{b}_1\mathrm{\Omega }_1\stackrel{~}{b}_2{\displaystyle \underset{l,r}{}}\mathrm{\Omega }_{lr}^{}\stackrel{~}{b}_{1lr}=i`$ (9) $`(EE_2)\stackrel{~}{b}_2\mathrm{\Omega }_1\stackrel{~}{b}_1\mathrm{\Omega }_2\stackrel{~}{b}_3{\displaystyle \underset{l,r}{}}\mathrm{\Omega }_{lr}\stackrel{~}{b}_{2lr}=0`$ (10) $`(E+E_lE_1E_r)\stackrel{~}{b}_{1lr}\mathrm{\Omega }_{lr}^{}\stackrel{~}{b}_1\mathrm{\Omega }_1\stackrel{~}{b}_{2lr}`$ (11) $`{\displaystyle \underset{l^{},r^{}}{}}\mathrm{\Omega }_{l^{}r^{}}^{}\stackrel{~}{b}_{1ll^{}rr^{}}=0`$ (12) $`(E+E_lE_2E_r)\stackrel{~}{b}_{2lr}\mathrm{\Omega }_{lr}\stackrel{~}{b}_2\mathrm{\Omega }_1\stackrel{~}{b}_{1lr}`$ (13) $`\mathrm{\Omega }_2\stackrel{~}{b}_{3lr}{\displaystyle \underset{l^{},r^{}}{}}\mathrm{\Omega }_{l^{}r^{}}\stackrel{~}{b}_{2ll^{}rr^{}}=0`$ (14) $`\mathrm{},`$ (15) where $`\mathrm{\Omega }_{lr}^{}=\mathrm{\Omega }_{lr}+\delta \mathrm{\Omega }_{lr}`$. Eqs. (Delocalization in the Anderson model due to a local measurement) can be converted to Bloch-type equations for the density matrix $`\sigma _{jj^{}}(t)`$ without their explicit solution. This technique has been derived in . We explain below only the main points of this procedure and the conditions for its validity. Consider, for example, Eq. (9). In order to perform the summation in the term $`_{l,r}\mathrm{\Omega }_{lr}^{}\stackrel{~}{b}_{1lr}`$, we solve for $`\stackrel{~}{b}_{1lr}`$ in Eq. (12). Then substituting the result into the sum, we can rewrite Eq. (9) as $`\left(EE_1{\displaystyle \frac{\mathrm{\Omega }_{}^{}{}_{lr}{}^{2}\rho _L(E_l)\rho _R(E_r)dE_ldE_r}{E+E_lE_1E_r}}\right)\stackrel{~}{b}_1`$ (16) $`\mathrm{\Omega }_1\stackrel{~}{b}_2+=i,`$ (17) where we have replaced the sum in Eq. (9) by an integral $`_{l,r}\rho _L(E_l)\rho _R(E_r)𝑑E_l𝑑E_r`$, with $`\rho _{L,R}`$ the density of states in the emitter and collector. We split this integral into its principle value and singular part. The singular part yields $`iD^{}/2`$, where $`D^{}=2\pi \mathrm{\Omega }_{}^{}{}_{}{}^{2}\rho _L\rho _R(\mu _L\mu _R)`$, and the principal part is zero, providing $`\mathrm{\Omega }_{lr}^{}`$ and $`\rho _{L,R}`$ are weakly dependent on the energies $`E_{l,r}`$. Note that $`(2\pi )^2\mathrm{\Omega }^2\rho _L\rho _R=T`$, where $`T`$ is the tunneling transmission coefficient of the point contact. Thus, $`eD^{}=I^{}`$ is the current flowing through the point contact whenever the electron occupies the first dot. The quantity $``$ in Eq. (17) denotes the terms in which the amplitudes $`\stackrel{~}{b}`$ cannot be factored out of the integrals. These terms vanish in the large-bias limit, $`(\mu _L\mu _R)\mathrm{\Omega }^2\rho `$. Indeed, all the singularities of the amplitude $`\stackrel{~}{b}(E,E_l,E_l^{},E_r,E_r^{})`$ in the $`E_l,E_l^{}`$ variables lie below the real axis. This can be seen directly from Eqs. (Delocalization in the Anderson model due to a local measurement) by noting that $`E`$ lies above the real axis in the Laplace transform. Assuming that the transition amplitudes $`\mathrm{\Omega }`$ as well as the densities of states $`\rho _{L,R}`$ are independent of $`E_{l,r}`$, one can close the integration contour in the upper $`E_{l,r}`$-plane. Since the integrand decreases faster than $`1/E_{l,r}`$, the resulting integrals are zero. Applying analogous considerations to the other equations of the system (Delocalization in the Anderson model due to a local measurement) we convert Eqs. (Delocalization in the Anderson model due to a local measurement) directly into rate equations via the inverse Laplace transform. The details can be found in . Here we present only the final equations for the electron density matrix $`\sigma _{jj^{}}(t)`$: $`\dot{\sigma }_{jj}=i\mathrm{\Omega }_{j1}(\sigma _{j,j1}\sigma _{j1,j})`$ (19) $`+i\mathrm{\Omega }_j(\sigma _{j,j+1}\sigma _{j+1,j}),`$ (20) $`\dot{\sigma }_{jj^{}}=iϵ_{j^{}j}\sigma _{jj^{}}+i\mathrm{\Omega }_{j^{}1}\sigma _{j,j^{}1}+i\mathrm{\Omega }_j^{}\sigma _{j,j^{}+1}`$ (21) $`i\mathrm{\Omega }_{j1}\sigma _{j1,j^{}}i\mathrm{\Omega }_j\sigma _{j+1,j^{}}{\displaystyle \frac{\mathrm{\Gamma }}{2}}\sigma _{jj^{}}(\delta _{1j}+\delta _{1j^{}}),`$ (22) where $`ϵ_{j^{}j}=E_j^{}E_j`$ and $`\mathrm{\Gamma }=(\sqrt{I/e}\sqrt{I^{}/e})^2`$ is the decoherence rate, generated by interaction with the detector. Note that these equations have been obtained from the many-body Schrödinger equation for the entire system. No stochastic assumptions have been made in their derivation. Eqs. (Delocalization in the Anderson model due to a local measurement) are analogous to the well-known optical Bloch equations used to describe a multilevel atom interacting with the quantized electromagnetic field. To our knowledge, this is their first appearance in connection with the Anderson model. The equations can be rewritten in Lindblad form as $$\dot{\sigma }=i[H_A,\sigma ]\frac{\mathrm{\Gamma }}{2}(Q\sigma +\sigma Q2\stackrel{~}{Q}\sigma \stackrel{~}{Q}^{}),$$ (23) where $`H_A`$ is given by Eq. (1) and $`Q_{jj^{}}=\stackrel{~}{Q}_{jj^{}}=\delta _{1j}\delta _{1j^{}}`$. If $`\mathrm{\Gamma }=0`$, Eq. (23) is equivalent to the Schrödinger equation $`i_t\|\mathrm{\Psi }(t)=H_A\|\mathrm{\Psi }(t)`$. In this case the electron density matrix $`\sigma (t)`$ displays Anderson localization, i.e., $`\sigma _{NN}(t\mathrm{})\mathrm{exp}(\alpha N)`$. If $`\mathrm{\Gamma }0`$, however, the asymptotic behavior of the reduced density-matrix, $`\sigma _{jj^{}}(t\mathrm{})`$, changes dramatically: all eigenfrequencies (except for the zero mode) obtain an imaginary part due to the second (damping) term in Eq. (23), so that only the stationary terms survive in the limit $`t\mathrm{}`$. This damping is illustrated in Fig. 2 which displays the numerical solution of Eqs. (Delocalization in the Anderson model due to a local measurement) for $`N=4`$, $`\mathrm{\Omega }_j=\overline{\mathrm{\Omega }}`$=const, and $`E_j/\overline{\mathrm{\Omega }}=\{0,2,4,1\}`$. The occupation of the first dot, $`P_1(t)=\sigma _{11}(t)`$, and the last dot, $`P_4(t)=\sigma _{44}(t)`$, is shown in Fig. 2 for $`\mathrm{\Gamma }=0`$ by the dashed lines, and for $`\mathrm{\Gamma }/\overline{\mathrm{\Omega }}=1`$ by the solid lines. One can clearly see that all oscillations decay for $`\mathrm{\Gamma }0`$, so that the density matrix reaches a stationary limit. Then we see the opposite of localization, as the probability of finding the electron in the last dot, $`P_4(t)`$, becomes the same as the probability of finding it in the first dot, $`P_1(t)`$. The delocalization phenomenon, illustrated by Fig. 2, can be proven analytically for any $`N`$. Indeed, let us consider Eqs. (Delocalization in the Anderson model due to a local measurement) in the asymptotic limit $`t\mathrm{}`$, where the electron density matrix reaches its stationary limit: $`\sigma _{jj^{}}(t\mathrm{})=u_{jj^{}}+iv_{jj^{}}`$. Since for the stationary solution $`_t\sigma _{jj^{}}0`$, Eqs. (Delocalization in the Anderson model due to a local measurement) become $`0`$ $`=`$ $`ϵ_{j^{}j}v_{jj^{}}+\mathrm{\Omega }_{j^{}1}v_{j,j^{}1}+\mathrm{\Omega }_j^{}v_{j,j^{}+1}\mathrm{\Omega }_{j1}v_{j1,j^{}}`$ (26) $`\mathrm{\Omega }_jv_{j+1,j^{}}+{\displaystyle \frac{\mathrm{\Gamma }}{2}}u_{jj^{}}(\delta _{1j}+\delta _{1j^{}})(1\delta _{jj^{}}),`$ $`0`$ $`=`$ $`ϵ_{j^{}j}u_{jj^{}}+\mathrm{\Omega }_{j^{}1}u_{j,j^{}1}+\mathrm{\Omega }_j^{}u_{j,j^{}+1}\mathrm{\Omega }_{j1}u_{j1,j^{}}`$ (28) $`\mathrm{\Omega }_ju_{j+1,j^{}}{\displaystyle \frac{\mathrm{\Gamma }}{2}}v_{jj^{}}(\delta _{1j}+\delta _{1j^{}})(1\delta _{jj^{}}).`$ Eqs. (Delocalization in the Anderson model due to a local measurement) have the unique solution $`v_{jj^{}}=0`$ and $`u_{jj^{}}=(1/N)\delta _{jj^{}}`$. This can be obtained by solving these equations sequentially, starting with $`j,j^{}=N`$, and then continuing for $`j,j^{}=N1,N2,\mathrm{}`$. Since $`u_{jj}\sigma _{jj}(t\mathrm{})`$, we finally obtain that $$\sigma _{jj^{}}(t\mathrm{})=\frac{1}{N}\delta _{jj^{}}.$$ (29) This corresponds to the totally delocalized electron state. Since Eq. (29) represents the unique solution of Eqs. (Delocalization in the Anderson model due to a local measurement), it implies that the asymptotic behavior of the electron density matrix is always given by Eq. (29) for any initial conditions. Note that this result is true only for $`\mathrm{\Gamma }0`$. Otherwise the solution of Eqs. (Delocalization in the Anderson model due to a local measurement) is not unique. Eq. (29) tells us that an arbitrarily weak interaction with the environment (detector) leads to delocalization in the Anderson model, even though this interaction affects only one of the sites. In other words, Anderson localization is unstable under infinitely small decoherence. One aspect of this instability is the importance of the order of limits $`t\mathrm{}`$ and $`N\mathrm{}`$. Taking $`t\mathrm{}`$ first, as above, gives delocalization, while taking $`N\mathrm{}`$ first would preserve localization. In the non-interacting model, $`\mathrm{\Gamma }=0`$, the order of limits is immaterial and the electron is localized. Even though a local interaction with the environment destroys the localization, the latter should affect the time-dependence of the observed system. We expect the delocalization time to increase exponentially with $`N`$ and to be dependent on both the decoherence rate and the localization length. This matter deserves further investigation. We would like to stress that our result is not an effect of finite temperature, as is so called the hopping conductivity. In the latter case, each site of the Anderson model interacts with the thermal bath; in our case, only one site is coupled to the detector (environment). If we were to let all the sites interact equally with the detector ($`I=I^{}`$, Fig. 1), we would obtain no delocalization in our model, since $`\mathrm{\Gamma }=0`$ in Eqs. (Delocalization in the Anderson model due to a local measurement), (23) (see also). Indeed, in this case Eq. (23) is equaivalent to the Scrödinger equation $`i_t\|\mathrm{\Psi }(t)=H_A\|\mathrm{\Psi }(t)`$ leading to Anderson localization. Note that there is no measurement when $`I=I^{}`$. The origin of delocalization in our case is therefore the break of coherence due to the measurement process. Delocalization of the Anderson model due to measurement has been studied previously. Yet the limit of a local and weak measurement has not been achieved. In the present work we include the detector in the quantum mechanical description, avoiding the use of the projection postulate in the course of measurement. This enables us to study delocalization due to local measurement and also in the limit of weak coupling with the measurement device. Another experimental setup for delocalization due to a local measurement is shown schematically in Fig. 3. It can be realized in atomic systems, for instance, in experiments with Rydberg atoms. For $`N=2`$ this setup is similar to a V-level system used for investigation of the quantum Zeno effect. The occupation of $`E_1`$ is monitored via spontaneous photon emission, where the $`01`$ Rabi transition is generated by a laser field. Using the same derivation as in the previous case, Fig. 1, we obtain Bloch equations for the reduced electron density matrix $`\sigma _{jj^{}}(t)`$ where $`j,j^{}=0,1,\mathrm{},N`$. The off-diagonal density-matrix elements are described by the same Eq. (22). Equation (20) for the diagonal density-matrix elements, however, is modified. Now it reads $`\dot{\sigma }_{jj}=i\mathrm{\Omega }_{j1}(\sigma _{j,j1}\sigma _{j1,j})`$ (30) $`+i\mathrm{\Omega }_j(\sigma _{j,j+1}\sigma _{j+1,j})\mathrm{\Gamma }(\delta _{j1}\delta _{j0})\sigma _{11}.`$ (31) The last term in Eq. (31) describes the rates due to spontaneous photon emission, Fig. 3. Here again the Bloch equations for the electron density matrix can be rewritten in Lindblad form, Eq. (23), with $`Q_{jj^{}}=\delta _{1j}\delta _{1j^{}}`$ and $`\stackrel{~}{Q}_{jj^{}}=\delta _{0j}\delta _{1j^{}}`$. (For $`N=2`$ Eq. (23) coincides with the optical Bloch equations used for analysis of a V-level system). Similar to the previous case, Fig. 1, Anderson localization is destroyed for any value of $`\mathrm{\Gamma }`$, and the asymptotic electron distribution, $`\sigma _{jj}(t\mathrm{})`$, does not depend on the initial electron state. Here, however, the electron density matrix in the asymptotic state is not a pure mixture, $`\sigma _{jj^{}}(t\mathrm{})0`$, and the probabilities $`\sigma _{jj}(t\mathrm{})`$ are not equally distributed between different wells (c.f. Eq. (29)). The delocalization of the Anderson model should also affect its transport properties. Indeed, by connecting the first and the last dot in Fig. 1 to leads (reservoirs) one can expect current to flow through the dot array whenever any of the dots is monitored. Indeed, the stationary current through coupled dots is proportional to the occupation probability of the last dot, attached to the collector. The current should appear with a delay after a voltage bias to the leads is switched on. This time delay is precisely the relaxation time needed for the electron to be delocalized. Anderson localization appears not only in quantum mechanics, but also in classical wave mechanics. Therefore the described delocalization due to local interaction with an environment should have a classical analogy. It can appear, for instance, in propagation of waves through coupled cavities with randomly distributed resonant frequencies. A wave cannot ordinarily penetrate through such a system due to the Anderson localization. Random vibration of one of the cavities, however, should destroy the localization, so that waves begin to penetrate through the system after some time delay, corresponding to the delocalization time. Such an experiment can also be done using the system of transparent plates with randomly varying thicknesses, described in. I am grateful to A. Buchleitner, B. Elattari, U. Smilansky and B. Svetitsky for very useful discussions and important suggestions.
warning/0002/math0002077.html	ar5iv	text	# Differential 3-knots in 5-space with and without self intersections ## 1. Introduction Immersions and embeddings form open subspaces of the space of $`C^r`$-maps $`S^k^{k+n}`$, $`r1`$. Smale showed that the path components of the space of immersions (or, which is the same, the regular homotopy classes of immersions) $`S^k^{k+n}`$ are in one to one correspondence with the elements of $`\pi _k(V_{k+n,k})`$, the $`k^{th}`$ homotopy group of the Stiefel manifold of $`k`$-frames in $`(k+n)`$-space. This far-reaching result translates problems in geometry to homotopy theory. Indicating the way back to geometry, Smale suggested the following problems (, p.329): “Find explicit representatives of regular homotopy classes…What regular homotopy classes have an embedding for representative?”. Explicit representatives of regular homotopy classes of immersions $`S^3^5`$ are given in Section 8.3. An answer to the second problem gives information about how the inclusion of the space of embeddings into the space of immersions is organized. In the case $`S^3^5`$, the group $`\pi _3(V_{5,3})`$, enumerating regular homotopy classes is infinite cyclic and the answer to the second problem was found by Hughes and Melvin . They proved that exactly every $`24^{\mathrm{th}}`$ regular homotopy class have an embedding for representative. In Theorem 1 below we describe the obstruction for a generic immersion $`S^3^5`$ to be regularly homotopic to an embedding in terms of geometric invariants of its self intersection. In dimensions where there are embeddings in different regular homotopy classes, it is impossible to express the regular homotopy class of a generic immersion in terms of its self intersection. This is in contrast to many other cases where this is possible: For example, in the cases $`S^k^{2k}`$, $`k2`$ the regular homotopy class is determined by the algebraic number of self intersection points (modulo 2 if $`k`$ is odd), see , and in the cases $`S^k^{2kr}`$, $`r=1,2`$ self intersection formulas for regular homotopy are given in and . Generic immersions have simple self intersections. For example, in the case $`S^1^3`$ a generic immersion has empty self intersection and thus, generic immersions are embeddings. Form this point of view, the analogue of classical knot theory in other dimensions is the study of path components of the space of generic immersions and we may think of generic immersions as knots with self intersections. The space of generic immersions is dense in the corresponding space of immersions and its complement is a stratified hypersurface. Using the stratification of this complement, Vassiliev introduced the notion of finite order invariants in classical knot theory. There are natural analogies of this notion for invariants of generic immersions in other dimensions (see Section 6.3). Arnold found first order invariants of generic plane curves ($`S^1^2`$). Theorem 2 below shows that, up to first order, the space of generic immersions $`S^3^5`$ is similar to the space of generic plane curves. The composition of an immersions $`S^3^4`$ and the inclusion $`i:^4^5`$ is an immersion into 5-space. Proposition 7.1.2 shows that two immersions of $`S^3`$ into 4-space are regularly homotopic in 5-space, after composing them with the inclusion, if and only if one of them is regularly homotopic in 4-space to the connected sum of the other one and a finite number of immersions regularly homotopic to the composition of the standard embedding and a reflection in a hyperplane in $`^4`$. Theorems 3, 4, and 5 give information about the self intersections of a generic immersions $`g:S^3^4`$ and its relation to the self intersection of a generic immersion $`f:S^3^5`$ regularly homotopic to $`ig`$. ## 2. Main results In this section, the main theorems of the paper are formulated. ### 2.1. Embeddings in the space of immersions Regular homotopy classes of immersions $`S^3^5`$ are known to form an infinite cyclic group $`\mathrm{𝐈𝐦𝐦}`$ under connected sum. The classes that contain embeddings form a subgroup $`\mathrm{𝐄𝐦𝐛}\mathrm{𝐈𝐦𝐦}`$. We now state Theorem 1 describing the extension $`\mathrm{𝐄𝐦𝐛}\mathrm{𝐈𝐦𝐦}`$ algebraically and devote the rest of this section to identify the homomorphisms involved there in topological terms. ###### Theorem 1. The following diagram of Abelian groups has exact rows, commutes and the vertical arrows are isomorphisms $$\begin{array}{ccccccccc}0& & \mathrm{𝐄𝐦𝐛}& & \mathrm{𝐈𝐦𝐦}& \stackrel{\lambda \beta }{}& _3_8& & 0\\ & & \sigma & & \mathrm{\Omega }& & & & & & \\ 0& & & \stackrel{\times 24}{}& & & _{24}& & 0\end{array}.$$ Theorem 1 is proved in Section 9.1. The homomorphism $`\sigma `$: Recall that if $`f:S^3^5`$ is an embedding then there exists a compact orientable 4-dimensional manifold $`V^4^5`$ with $`V^4=f(S^3)`$. We call such a manifold a Seifert-surface of $`f`$. Its signature $`\sigma (V^4)`$ is divisible by 16 and is known to depend only on $`f`$. For $`\xi \mathrm{𝐄𝐦𝐛}`$, define $`\sigma (\xi )`$ as $`\frac{\sigma (V)}{16}`$, where $`V`$ is a Seifert-surface of an embedding representing $`\xi `$. It is proved in that $`\sigma `$ induces an isomorphism $`\mathrm{𝐄𝐦𝐛}`$ and that $`\mathrm{𝐄𝐦𝐛}\mathrm{𝐈𝐦𝐦}`$ is a subgroup of index $`24`$. The homomorphism $`\mathrm{\Omega }`$: Any immersion $`f:S^3^5`$ is determined up to regular homotopy by its Smale invariant $`\mathrm{\Omega }(f)\pi _3(V_{5,3})`$ (see or Definition 3.2.1). For $`\xi \mathrm{𝐈𝐦𝐦}`$, define $`\mathrm{\Omega }(\xi )`$ as $`\mathrm{\Omega }(g)`$, where $`g`$ is an immersion representing $`\xi `$. The homomorphism $`\lambda `$: If $`f:S^3^5`$ is a generic immersion then its self intersection $`M_f^5`$ is a closed 1-dimensional manifold. Orientations of $`S^3`$ and $`^5`$ induce an orientation of $`M_f`$. We can push $`M_f`$ off the image of $`f`$ (see Section 6.2). Let $`\mathrm{lk}(f)`$ denote the linking number of the perturbed $`M_f`$ and $`f(S^3)`$ in $`^5`$ . For $`\xi \mathrm{𝐈𝐦𝐦}`$, define $`\lambda (\xi )_3`$ as $`\mathrm{lk}(g)`$ modulo 3, where $`g`$ is a generic immersion representing $`\xi `$. The homomorphism $`\beta `$: The normal bundle of an immersion $`f_0:S^3^5`$ is 2-dimensional orientable and therefore trivial. Hence, $`f_0`$ admits a normal vector field. This implies that $`f_0`$ is regularly homotopic to a (generic) immersion $`f_1:S^3^4`$, composed with the inclusion $`^4^5`$ (see ). Resolving the self intersection of a generic immersion $`f:S^3^4`$, we obtain a smooth surface $`F_f`$ (Lemma 5.1.5) and the immersion $`f`$ induces a pin (i.e. $`Pin^{}`$) structure on $`F_f`$ (Section 7.4). There is a one to one correspondence between pin structures on a surface $`F`$ and $`_4`$-quadratic functions $`q`$ on its first homology $`H_1(F;_2)`$. Pin structures on a surface are classified up to cobordism by the Brown invariant $`\beta (q)_8`$ of the corresponding function. Let $`\beta (f)_8`$ denote the Brown invariant of the quadratic function corresponding to the pin structure induced by $`f`$ on $`F_f`$. For $`\xi \mathrm{𝐈𝐦𝐦}`$, define $`\beta (\xi )_8`$ as $`\beta (g)`$, where $`g:S^3^4^5`$ is a generic immersion representing $`\xi `$. The author does not know how to calculate $`\beta (f)`$ for a generic immersion $`f:S^3^5`$ in terms of the geometry of its self intersection, without pushing it down to $`^4`$. However, $`\beta (f)`$ modulo 4, can be calculated in that way (see Section 9.2). ### 2.2. First order invariants of generic immersions In generic one-parameter families of immersions there are isolated instances of non-generic immersions. In generic one-parameter families of immersions $`S^3^5`$ such instances are immersions with one self tangency or one triple point. The same is true for generic one-parameter families of plane curves where one can distinguish two local types of self tangencies: direct (the tangent vectors point in the same direction) and reverse (the tangent vectors point in opposite directions). An invariant of generic immersions is a function which is constant on the path components of the space of generic immersions. Such a function may change when we pass through instances of non-generic immersions. Arnold found three independent first order invariants of generic plane curves: The invariant $`J^+`$ which changes at direct self tangency instances and does not change under other moves, the invariant $`J^{}`$ which changes under reverse self tangency moves and does not change under other moves, and the invariant Strangeness $`\mathrm{St}`$ which changes under triple point moves and does not change under other moves. For a generic immersion $`f:S^3^5`$, define $`J(f)`$ to be the number of self intersection components of $`f`$ and define $`L(f)=\frac{1}{3}(\mathrm{lk}(f)\stackrel{~}{\lambda }(f))`$, $`\stackrel{~}{\lambda }\{0,1,2\}`$ is a lifting of $`\lambda _3`$ (see Section 2.1). ###### Theorem 2. The invariant $`J`$ changes by $`\pm 1`$ under self tangency moves and does not change under triple point moves. The invariant $`L`$ changes by $`\pm 1`$ under triple point moves and does not change under self tangency moves. The invariants $`J`$ and $`L`$ are first order invariants. Moreover, if $`v`$ is any first order invariant then the restriction of $`v\|U`$, where $`U`$ is a path component of the space of immersions, is a linear combination of $`J\|U`$ and $`L\|U`$. Theorem 2 is proved in Section 6.4. Although there are two local types of selftangencies in the case $`S^3^5`$ (see Proposition 5.3.2), in contrast to the case of plane curves, $`J`$ can not be splitted into more refined first order invariants. Arnold defined the invariant $`\mathrm{St}`$ in such a way that it is additive under connected summation of plane curves and showed that this property together with the changes of $`\mathrm{St}`$ under local moves and its is orientation independence completely characterizes the invariant (up to a constant). The invariant $`L`$ is neither additive under connected summation nor symmetric with respect to orientation. To get the analogue of Arnold’s strangeness of plane curves for immersions $`f:S^3^5`$ we define $`\mathrm{St}(f)=\frac{1}{3}(\mathrm{lk}(f)+\mathrm{\Omega }(f))`$. Then $`\mathrm{St}`$ is additive under connected sum, changes exactly as $`L`$ under local moves, changes sign under composing immersions with an orientation reversing diffeomorphism of $`S^3`$, and is completely characterized by these properties up to a constant (see Proposition 6.5.1). The reason for stating Theorem 2 in terms of $`L`$ instead of $`\mathrm{St}`$ (which works equally well) is that $`L(f)`$ in contrast to $`\mathrm{St}(f)`$ can be calculated in terms of the self intersection of $`f`$. ### 2.3. Immersions into 4-space The self intersection of a generic immersion$`f:S^3^4`$ consists of 2-dimensional sheets of double points, 1-dimensional curves of triple points and isolated quadruple points (see Definition 5.1.1). Resolving the quadruple and triple points turns the self intersection into a smooth surface $`F_f`$ (Lemma 5.1.5). The following three theorems are proved in Section 9.2. ###### Theorem 3. A generic immersion $`g:S^3^4`$ has an odd number of quadruple points if and only if its (resolved) self intersection surface $`F_g`$ has odd Euler characteristic. As mentioned (Section 2.1), any immersion $`S^3^5`$ is regularly homotopic to a composition of an immersion $`S^3^4`$ and the inclusion $`^4^5`$. ###### Theorem 4. Let $`f:S^3^5`$ be a generic immersion and let $`g:S^3^4^5`$ be a generic immersion regularly homotopic to $`f`$. Then $`g`$ has an odd number of quadruple points if and only if $`f`$ has an odd number of self intersection components with connected preimage. We define a $`_4`$-valued invariant $`\tau `$ of generic immersions $`S^3^5`$ (see Section 7.5). If $`f`$ is a generic immersion then $`\tau (f)`$ is divisible by 2 if and only if $`f`$ has an odd number of self intersection components with connected preimage. ###### Theorem 5. Let $`f:S^3^5`$ be a generic immersion and let $`g:S^3^4^5`$ be a generic immersion regularly homotopic to $`f`$. If $`\tau (f)0`$ then the (resolved) self intersection surface $`F_g`$ of $`g`$ is nonorientable. ## 3. The Smale invariant and Stiefel manifolds In this section, some properties of the Smale invariant and the groups where it takes values are collected for later reference. Via the Smale invariant the set of regular homotopy classes of immersions is endowed with the structure of an Abelian group. We describe the group operations in topological terms. Most of the results presented here are known. Proofs are provided where there were hard to find references. ### 3.1. Homotopy groups of two Stiefel manifolds The Stiefel manifold $`V_{4,3}`$ is homotopy equivalent to $`SO(4)`$. Indeed, any orthonormal 3-frame in 4-space can be uniquely completed to a positively oriented $`4`$-frame. Let $`SO(4)\stackrel{𝑝}{}S^3`$ be the fibration with fiber $`SO(3)`$. Consider $`S^3`$ as the set of unit quaternions in $`^4`$, where we identify the vectors $`_1,\mathrm{},_4`$ in the standard base of $`^4`$ with $`1,i,j,k`$. For $`x,y^4`$ let $`xy^4`$ denote quaternionic product of $`x`$ and $`y`$. The map $`\sigma :S^3SO(4)`$, $`\sigma (x)y=xy`$, is a section of $`SO(4)\stackrel{𝑝}{}S^3`$. Thus, $`SO(4)`$ is diffeomorphic to $`S^3\times SO(3)`$. Let $`\varrho :S^3SO(3)`$ be the map $`\varrho (x)u=xux^1`$, where $`u`$ is a pure quaternion and $`^3`$ is identified with the set of pure quaternions. (That is, the span of the vectors $`i,j,k`$.) Let $`\rho :S^3SO(4)`$ be $`\varrho `$ composed with the inclusion of the fiber over $`(1,0,0,0)S^3`$. The following is immediate. ###### Lemma 3.1.1. $$\pi _3(SO(4))[\sigma ][\rho ].$$ Let $`SO(5)\stackrel{𝑝}{}V_{5,3}`$ be the fibration with fiber $`SO(2)`$ that maps an orthonormal $`5\times 5`$-matrix to the 3-frame consisting of its first three column vectors. Similarly, let $`SO(5)\stackrel{𝑟}{}S^4`$ be the fibration with fiber $`SO(4)`$. ###### Lemma 3.1.2. The homomorphism $`p_{}:\pi _3(SO(5))\pi _3(V_{5,3})`$ is an isomorphism. The homomorphism $`r_{}:\pi _3(SO(4))\pi _3(SO(5))`$ is an epimorphism with kernel $`N`$, where $`N`$ is the subgroup generated by $`[\sigma ]2[\rho ]`$. ###### Proof. The first statement follows by inspecting the homotopy sequence of the fibration. For the second, see , Chapter 8, Proposition 12.11. ∎ Using Lemma 3.1.2, we make the identifications $$\pi _3(V_{5,3})=\pi _3(SO(5))=\pi _3(SO(4))/N.$$ ### 3.2. The Smale invariant of an immersed 3-sphere We define the Smale invariant of an immersion $`f:S^3^n`$: Consider $`S^3`$ as the unit sphere in $`^4`$ and let $`s:S^3^4^n`$, denote the standard embedding. Fix a disk $`D^3S^3`$ containing the south pole and a framing $`X`$ of $`S^3D^3`$. Using regular homotopy, deform $`f`$ so that $`f\|D^3=s\|D^3`$. Choose a diffeomorphism $`r:H_+S^3D^3`$ of degree $`+1`$, where $`H_+`$ is the hemisphere $`\{x_00\}`$ in the unit sphere $`\{x_0^2+x_1^2+x_2^2+x_3^2=1\}`$ in $`^4`$. Let $`xx^{}`$ be the map $`(x_0,x_1,x_2,x_3)(x_0,x_1,x_2,x_3)`$ of the unit sphere in $`^4`$. Define $`\varphi _s^f:S^3V_{n,3}`$, $$\varphi _s^f(x)=\{\begin{array}{cc}df(X(r(x)))\hfill & \text{for }xH_+,\hfill \\ ds(X(r(x^{})))\hfill & \text{for }xH_{},\hfill \end{array}$$ where $`H_{}`$ is the hemisphere $`\{x_00\}`$. ###### Definition 3.2.1. The Smale invariant $`\mathrm{\Omega }(f)`$ of $`f`$ is $$\mathrm{\Omega }(f)=[\varphi _s^f]\pi _3(V_{n,3}).$$ Smale, showed that $`\mathrm{\Omega }`$ gives a bijection between the regular homotopy classes of immersions $`S^3^n`$ and the elements of $`\pi _3(V_{n,3})`$. ### 3.3. Calculating Smale invariants in 4- and 5-space In computations we will not use Definition 3.2.1 literally. We use a slightly different approach: Consider $$S^3=\{x^4=:x_0^2+x_1^2+x_2^2+x_3^2=1\},$$ where $``$ denotes the quaternions. The tangent space of $`S^3`$ at $`(1,0,0,0)`$ is the span of the vectors $`i,j,k`$. We trivialize $`TS^3`$ using the quaternion framing: $$Q(x)=(xi,xj,xk)T_xS^3,\text{ for }xS^3\text{.}$$ Let $`f:S^3^n`$ be an immersion. Then there is an induced map $`\mathrm{\Phi }_n^f:S^3V_{n,3}`$, $`\mathrm{\Phi }_n^f(x)=df(Q(x))`$. Assume that $`n=4`$. Then $`V_{n,3}=V_{4,3}=SO(4)`$. Thus, we get a map $`\mathrm{\Phi }_4^f:S^3SO(4)`$. ###### Lemma 3.3.1. Let $`f:S^3^4`$ be an immersion then $$\mathrm{\Omega }(f)=[\mathrm{\Phi }_4^f][\mathrm{\Phi }_4^s]\pi _3(SO(4)),$$ and $`[\mathrm{\Phi }_4^s]=[\sigma ]\pi _3(SO(4))`$.∎ Assume that $`n=5`$. Consider the fibration $`SO(5)V_{5,3}`$ as in Lemma 3.1.2. Since the normal bundle of $`f`$ is orientable 2-dimensional, it is trivial and we can lift $`\mathrm{\Phi }_5^f`$ to $`\mathrm{\Theta }^f:S^3SO(5)`$ and $`[\mathrm{\Theta }^f]\pi _3(SO(5))`$ is independent of this lifting. ###### Lemma 3.3.2. Let $`f:S^3^5`$ be an immersion then $$\mathrm{\Omega }(f)=[\mathrm{\Theta }^f][\mathrm{\Theta }^s]\pi _3(SO(5))=\pi _3(V_{5,3}),$$ and $`[\mathrm{\Theta }^s]=[\sigma ]+N\pi _3(SO(5))`$.∎ Let $`i:^4^5`$ denote the inclusion. ###### Lemma 3.3.3. If $`f:S^3^4`$ is an immersion then $$\mathrm{\Omega }(if)=\mathrm{\Omega }(f)+N\pi _3(SO(5))=\pi _3(V_{5,3}).$$ ###### Proof. Complete the framing of $`f`$ with a vector in the fifth direction and combine Lemma 3.3.1 and Lemma 3.3.2. ∎ ### 3.4. The immersion group Let $`\mathrm{𝐈𝐦𝐦}`$ denote the infinite cyclic group of regular homotopy classes of immersions $`S^3^5`$. The Smale invariant gives an isomorphism $`\mathrm{\Omega }:\mathrm{𝐈𝐦𝐦}\pi _3(V_{5,3})`$. First, we consider addition in $`\mathrm{𝐈𝐦𝐦}`$ and other groups of regular homotopy classes of immersions $`S^3^n`$: Given two immersions $`f,g:S^3^n`$ we define an immersion $`fg`$ as follows: Consider $`S^3^4`$ with coordinates $`x=(x_0,x_1,x_2,x_3)`$ as the subset characterized by $`x_i^2=1`$. Let $`a=(1,0,0,0)`$ and $`a^{}=(1,0,0,0)`$, respectively. Choose frames $`u_1,\mathrm{},u_n`$ at $`f(a)`$ and $`v_1,\mathrm{},v_n`$ at $`g(a^{})`$ that agree with the orientation of $`^n`$ and such that $`u_1,u_2,u_3`$ ($`v_1,v_2,v_3`$) are tangent to $`f(S^3)`$ (to $`g(S^3)`$) at $`f(a)`$ (at $`g(a^{})`$). We can assume, possibly after moving $`g(S^3)`$, that $`u_i=v_i`$, $`i=1,2,3`$ and that $`g(a^{})=f(a)+u_n`$. Moreover, we can deform the maps so that $`f(x)`$ $`=f(a)+{\displaystyle \underset{i=1}{\overset{3}{}}}x_iu_i\text{for}1ϵx_01,`$ $`g(x)`$ $`=g(a^{})+{\displaystyle \underset{i=1}{\overset{3}{}}}x_iv_i\text{for}1x_01+ϵ.`$ The immersion $`fg`$ is now obtained by running a tube from $`f(a)`$ to $`g(a^{})`$ with axis $`f(a)+tu_n`$. Details can be found in Kervaire , Section 2, where the following is proved. ###### Lemma 3.4.1. $$\mathrm{\Omega }(fg)=\mathrm{\Omega }(f)+\mathrm{\Omega }(g).$$ We call the immersion $`fg`$ the connected sum of $`f`$ and $`g`$. Secondly, we consider inversion in $`\mathrm{𝐈𝐦𝐦}`$: Given an immersion $`f:S^3^5`$ we define the immersion $`\widehat{f}`$ as follows: Let $`\widehat{f}=fr`$, where $`r:S^3S^3`$ is the restriction of a reflection through a hyperplane in $`^4`$. ###### Lemma 3.4.2. $$\mathrm{\Omega }(\widehat{f})=\mathrm{\Omega }(f)\pi _3(V_{5,3}).$$ ###### Proof. Let the disk $`US^3^4`$ in the definition of the Smale invariant be the hemisphere $`\{x_00\}`$ and let $`r:S^3S^3`$ be the map $$r(x_0,x_1,x_2,x_3)=(x_0,x_1,x_2,x_3).$$ Let $`R:^5^5`$ be the map $$R(y_1,y_2,y_3,y_4,y_5)=(y_1,y_2,y_3,y_4,y_5).$$ Then $`Rsr=s`$ and we can use the map $`\varphi _{Rsr}^{Rfr}:S^3V_{3,5}`$ to compute $`\mathrm{\Omega }(\widehat{f})`$. It is straightforward to check that this map is homotopic to $`\varphi _s^fr`$. Hence, $$\mathrm{\Omega }(\widehat{f})=(\varphi _{Rsr}^{Rfr})_{}[S^3]=(\varphi _s^fr)_{}[S^3]=(\varphi _s^f)_{}\left([S^3]\right)=\mathrm{\Omega }(f).$$ ## 4. Embeddings considered as immersions In this section we present the classification of embeddings $`S^3^5`$ up to regular homotopy. ### 4.1. Embeddings up to regular homotopy and signature Let $`f,g:S^3^5`$ be embeddings. Then $`fg`$ is regularly homotopic to an embedding and $`\widehat{f}`$ is an embedding. Thus, the regular homotopy classes that contain embeddings from a subgroup of $`\mathrm{𝐈𝐦𝐦}`$. We denote this subgroup $`\mathrm{𝐄𝐦𝐛}`$. We get a classification of embeddings up to regular homotopy as follows: Given an embedding $`f:S^3^5`$, we can find a compact connected orientable manifold $`V^4`$ embedded into $`^5`$ and such that $`V^4=f(S^3)`$. We call such a manifold a Seifert-surface of $`f`$. The orientation of $`f(S^3)`$ induces an orientation of $`V^4`$. Filling the $`S^3`$ on the boundary of $`V^4`$ with a 4-disk we get a closed connected oriented 4-manifold $`W^4`$. The cohomology sequence of the pair $`(W^4,V^4)`$ shows that the inclusion induces an isomorphism $`H^2(W^4;_2)H^2(V^4;_2)`$. The normal bundle of $`V^4`$ is a trivial 1-dimensional bundle. Hence, $`w_2(TV)=0`$ and therefore $`w_2(TW)=0`$, where $`w_2`$ is the second Stiefel-Whitney class. Thus, $`W`$ is a spin manifold. By Rokhlin’s theorem (see Milnor and Kervaire, ), the signature $`\sigma (W^4)`$ of $`W^4`$ is divisible by 16. Define $`\sigma (f)=\frac{\sigma (W)}{16}`$. (This definition agrees with that given in Section 2.1.) ###### Proposition 4.1.1. For $`\xi \mathrm{𝐄𝐦𝐛}`$, let $`\sigma (\xi )=\sigma (f)`$, where $`f`$ is an embedding representing $`\xi `$. Then $$\sigma :\mathrm{𝐄𝐦𝐛}$$ is an isomorphism. ###### Proof. See . ∎ ###### Proposition 4.1.2. The subgroup $`\mathrm{𝐄𝐦𝐛}\mathrm{𝐈𝐦𝐦}`$ has index $`24`$. ###### Proof. See . ∎ ## 5. Spaces of immersions In this section, generic immersions $`S^3^n`$, $`n=4,5`$ and their self intersections are studied. The space of immersions $`S^3^n`$ as described in the Introduction, will be denoted $`_n`$. It is an infinite dimensional manifold. The set of non-generic immersions in $`_n`$ is a stratified hypersurface $`\mathrm{\Sigma }_n`$. We describe its strata of codimension one, for $`n=4,5`$, and of codimension two, for $`n=5`$. ### 5.1. Generic immersions and their self intersections ###### Definition 5.1.1. An immersion $`f_4`$ is generic if it satisfies the following conditions: * For any $`w^5`$, $`f^1(w)`$ contains at most four points. * If $`f(x_1)=\mathrm{}=f(x_j)=w`$, $`2j4`$ for $`x_ix_jS^3`$ if $`ij`$, then $`df(T_{x_i}S^3)+_{ji}df(T_{x_j}S^3)=T_w^4`$. ###### Definition 5.1.2. An immersion $`f_5`$ is generic if it satisfies the following conditions: * For any $`w^5`$, $`f^1(w)`$ contains at most two points. * If $`f(x)=f(y)=w`$, for $`xyS^3`$, then $`df(T_xS^3)+df(T_yS^3)=T_w^5`$. If $`f:XY`$ is an immersion of manifolds then the self intersection of $`f`$ is the subset of points $`yY`$ such that $`f^1(y)`$ contains more than one point. We denote it $`M_f`$. We denote its preimage $`\stackrel{~}{M}_f`$. That is, $`\stackrel{~}{M}_f=f^1(M_f)X`$. ###### Lemma 5.1.3. Let $`f:S^3^5`$ be a generic immersion. Then $`M_f`$ and $`\stackrel{~}{M}_f`$ are closed 1-manifolds and $`f:\stackrel{~}{M}_fM_f`$ is a double cover. Moreover, there is an induced orientation on $`M_f`$. ###### Proof. The first statement follows from G2. If we order the oriented sheets coming together along $`M_f`$ then there is a standard way to assign an orientation to $`M_f`$. Since the codimension of $`\stackrel{~}{M}_f`$ is even this orientation is independent of the ordering. ∎ ###### Lemma 5.1.4. Let $`g:S^3^4`$ be a generic immersion. Then $`M_g`$ and $`\stackrel{~}{M}_g`$ are 2-dimensional stratified spaces, $$M_g=M_g^0M_g^1M_g^2\text{and}\stackrel{~}{M}_g=\stackrel{~}{M}_g^0\stackrel{~}{M}_g^1\stackrel{~}{M}_g^2,$$ where $`M_g^j`$ and $`\stackrel{~}{M}_g^j`$ are smooth manifolds of dimension $`j`$, $`j=0,1,2`$. The strata $`M_g^j`$ is the set of $`j`$-tuple points, $`\stackrel{~}{M}_g^j=g^1(M_g^j)`$, and $`g\|\stackrel{~}{M}_g^j`$ is a $`j`$-fold covering. ###### Proof. Immediate from g1 and g2. ∎ The next lemma shows how to resolve the self intersection of a generic immersion $`S^3^4`$. ###### Lemma 5.1.5. Let $`g:S^3^4`$ be a generic immersion. There exist closed surfaces $`\stackrel{~}{F}_g`$ and $`F_g`$ (and closed 1-manifolds $`\stackrel{~}{C}_g`$ and $`C_g`$), unique up to diffeomorphisms, and immersions $`s:\stackrel{~}{F}_gS^3`$ and $`t:F_g^4`$ ($`\sigma :\stackrel{~}{C}_gS^3`$ and $`\tau :C_g^4`$) such that $$\begin{array}{ccc}\stackrel{~}{F}_g& \stackrel{s}{}& \stackrel{~}{M}_gS^3\\ p& & g& & \\ F_g& \stackrel{t}{}& M_g^5\end{array},\begin{array}{ccc}\stackrel{~}{C}_g& \stackrel{\sigma }{}& \stackrel{~}{M}_g^0\stackrel{~}{M}_g^1S^3\\ \pi & & g& & \\ C_g& \stackrel{\tau }{}& M_g^0\stackrel{~}{M}_g^1^5\end{array},$$ commutes. The maps $`s`$ and $`t`$ are surjective, have multiple points only along $`\stackrel{~}{M}_g^0\stackrel{~}{M}_g^1`$ and $`M_g^0M_g^1`$, respectively and $`p`$ is the orientation double cover. The maps $`\sigma `$ and $`\tau `$ are surjective, have multiple points only along $`\stackrel{~}{M}_g^0`$ and $`M_g^0`$, respectively and $`\pi `$ is a $`3`$-fold cover. ###### Proof. This is immediate from the local pictures: Close to a $`j`$-tuple point $`M_g`$ is the intersection of $`j`$ 3-planes in general position in 4-space. ∎ We call $`F_g`$ the resolved self intersection surface of $`g`$. We use the notation $`F_g^j=t^1(M_g^j)`$ and $`\stackrel{~}{F}_g^j=s^1(\stackrel{~}{M}_g^j)`$, $`j=0,1`$. ### 5.2. The discriminant hypersurface and its stratum of codimension one The jet transversality theorem implies that the set of generic immersions is an open dense subset of $`_n`$, $`n=4,5`$. Its complement $`\mathrm{\Sigma }_n_n`$ will be called the discriminant hypersurface. The discriminant hypersurface is stratified, $`\mathrm{\Sigma }_n=\mathrm{\Sigma }_n^1\mathrm{\Sigma }_n^2\mathrm{}\mathrm{\Sigma }_n^{\mathrm{}}`$, where each stratum $`\mathrm{\Sigma }_n^k`$, $`k<\mathrm{}`$ is a smooth submanifold of codimension $`k`$ of $`_n`$ and $`\mathrm{\Sigma }_n^j`$ is contained in the closure of $`\mathrm{\Sigma }_n^1`$ for every $`j`$. The following two propositions follow by applying the jet transversality theorem to 1-parameter families of immersions. ###### Proposition 5.2.1. The codimension one stratum $`\mathrm{\Sigma }_4^1\mathrm{\Sigma }_4`$ is the set of all immersions $`f:S^3^4`$ such that * g1 and g2 holds except at one double point $`w=f(x)=f(y)^4`$, $`xyS^3`$ where $$dim\left(df(T_xS^3)+df(T_yS^3)\right)=3$$ or, * g1 and g2 holds except at one $`j`$-tuple point, $`3j5`$ $`w=f(x_1)=\mathrm{}=f(x_j)^4`$, $`x_ix_kS^3`$, if $`ik`$ where $$dim\left(df(T_{x_i}S^3)+\underset{rirk}{}df(T_{x_r}S^3)\right)=4,$$ for $`ki`$ but $$dim\left(df(T_{x_i}S^3)+\underset{ri}{}df(T_{x_r}S^3)\right)=3.$$ ###### Proposition 5.2.2. The codimension one stratum $`\mathrm{\Sigma }_5^1\mathrm{\Sigma }_5`$ is the set of all immersions $`f:S^3^5`$ such that either * G1 holds, $`f`$ has double points and G2 holds, except at one double point $`w=f(x)=f(y)^5`$, $`xyS^3`$, where $`dim\left(df(T_xS^3)+df(T_yS^3)\right)=4`$ or, * G2 holds and G1 holds except at one triple point $`w^5`$ with $`w=f(x_1)=f(x_2)=f(x_3)`$, where $`x_i`$, $`1i3`$ are three distinct points in $`S^3`$ and $$dim\left(df(T_{x_i}S^3)+\underset{ji}{}df(T_{x_j}S^3)\right)=4.$$ If (a) above holds for an immersion $`f`$, we say that the exceptional double point $`w`$ is a self tangency point of $`f`$. ### 5.3. Coordinate expressions and versal deformations Recall that a deformation of a map $`f_0`$ is a 1-parameter family of maps $`f_\lambda `$, parameterized by $`\lambda U`$, where $`U`$ is a neighborhood of $`0^m`$. A deformation $`f_\lambda `$ of a map $`f_0`$ is called versal if every deformation of $`f_0`$ is equivalent (up to left-right actions of diffeomorphisms) to one induced from $`f_\lambda `$. Below, $`f_0:S^3^n`$, $`n=4,5`$ will be an immersion in $`\mathrm{\Sigma }_n^1`$ with its exceptional multiple point at $`0^n`$, $`x,y,z,w,u`$ will denote coordinates in small 3-balls centered at the preimages of $`0`$ and $`f_t`$, $`tU`$ will be a versal deformation of $`f_0`$. Such a deformation can be assumed to be constant in $`t`$ outside of the coordinate balls. The proofs of the statements in this section and the next one are discussed in Section 5.5. ###### Proposition 5.3.1. Let $`f_0:S^3^4`$ be an immersion in $`\mathrm{\Sigma }_4^1`$ with exceptional multiple point $`0`$ and let $`f_t`$ be a versal deformation. Locally, at $`0`$, up to choice of coordinates around the preimages of $`0`$ and in $`^4`$, * if $`0`$ is double point then $`f_t`$ is of the form $`f_t(x)`$ $`=(x_1,x_2,x_3,0),`$ $`f_t(y)`$ $`=(y_1,y_2,y_3,y_1^2+y_2^2+ϵy_3^2+t),`$ where $`ϵ=\pm 1`$, or * if $`0`$ is a triple point then $`f_t`$ is of the form $`f_t(x)`$ $`=(x_1,x_2,x_3,0),`$ $`f_t(y)`$ $`=(y_1,y_2,0,y_3),`$ $`f_t(z)`$ $`=(z_1,z_2,z_3,z_1^2+ϵz_2^2+tz_3),`$ where $`ϵ=\pm 1`$, or * if $`0`$ is a quadruple point then $`f_t`$ is of the form $`f_t(x)`$ $`=(x_1,x_2,x_3,0),`$ $`f_t(y)`$ $`=(y_1,y_2,0,y_3),`$ $`f_t(z)`$ $`=(z_1,0,z_2,z_3),`$ $`f_t(w)`$ $`=(w_1,w_2,w_3,w_1^2+tw_2w_3),`$ or * if $`0`$ is a quintuple point then $`f_t`$ is of the form $`f_t(x)`$ $`=(x_1,x_2,x_3,0),`$ $`f_t(y)`$ $`=(y_1,y_2,0,y_3),`$ $`f_t(z)`$ $`=(z_1,0,z_2,z_3),`$ $`f_t(w)`$ $`=(0,w_1,w_2,w_3),`$ $`f_t(u)`$ $`=(u_1,u_2,u_3,tu_1u_2u_3).`$ ###### Proposition 5.3.2. Let $`f_0:S^3^5`$ be an immersion in $`\mathrm{\Sigma }_5^1`$ with exceptional multiple point $`0`$ and let $`f_t`$ be a versal deformation. Locally, at $`0`$, up to choice of coordinates around the preimages of $`0`$ and in $`^5`$, * if $`0`$ is double point then $`f_t`$ is of the form $`f_t(x)`$ $`=(x_1,x_2,0,x_3,0).`$ $`f_t(y)`$ $`=(y_1,y_2,y_1^2+ϵy_2^2+t,0,y_3).`$ where $`ϵ=\pm 1`$, or * if $`0`$ is a triple point then $`f_t`$ is of the form $`f_t(x)`$ $`=(x_1,x_2,x_3,0,0),`$ $`f_t(y)`$ $`=(y_1,t,0,y_2,y_3),`$ $`f_t(z)`$ $`=(0,z_2,z_1,z_2,z_3).`$ If $`ϵ=+1`$ ($`ϵ=1`$) in cases (a) in the above propositions we say that $`0`$ is an elliptic self tangency point (a hyperbolic self tangency point) of $`f_0`$. ### 5.4. The stratum of codimension two To study first order invariants of generic immersions (see Section 6.3) we need a description of the codimension two stratum of the discriminant hypersurface. We restrict attention to immersions into 5-space. ###### Proposition 5.4.1. The codimension two stratum $`\mathrm{\Sigma }_5^2\mathrm{\Sigma }`$ is the set of all immersions $`f_0:S^3^5`$ such that exactly one of the following holds * $`f_0`$ has two distinct self tangency points, * $`f_0`$ has two distinct triple points, * $`f_0`$ has one self tangency point and one triple point, * $`f_0`$ has a degenerate self tangency point at $`0=f(p)=f(q)`$, in which case its versal deformation $`f_{t,s}`$ is of the form $`f_{t,s}(x)`$ $`=(x_1,x_2,0,x_3,0),`$ $`f_{t,s}(y)`$ $`=(y_1,y_2,y_1^2+y_2(y_2^2+s)+t,0,y_3),`$ up to choice of coordinates $`x`$, $`y`$ around $`p`$, $`q`$, respectively, and coordinates in $`^5`$. The versal deformations of (a)-(d) are evident. They are just products of 1-dimensional versal deformations. Let $`E`$, $`Y`$ and $`T`$ denote the codimension one parts of $`\mathrm{\Sigma }_5`$ consisting of all immersions with one elliptic self tangency point, one hyperbolic self tangency point, and one triple point, respectively. Figure 1, which is a consequence of Proposition 5.4.1, shows the possible intersections of the discriminant hypersurface $`\mathrm{\Sigma }_5`$ and a small 2-disk in $`_5`$ which meets $`\mathrm{\Sigma }_5^2`$ and is transversal to $`\mathrm{\Sigma }_5`$. ### 5.5. Proofs The proofs of the propositions in Section 5.3 and Section 5.4 are all similar: First we need to find coordinates close to the preimages of the exceptional point such that the map is given by the expression stated. Then we show that the deformation given is infinitesimally versal. A standard theorem in singularity theory then implies that the deformation is versal. As an example we write out the proof of Proposition 5.3.1 (a): Since $`f_0`$ has a tangency at $`0`$ the tangent planes of the two sheets $`X`$ and $`Y`$ meeting there must agree. We may assume that this tangent plane is the plane of the first three coordinates. By the implicit function theorem we can choose coordinates so that the map of the first sheet is $$f_0(x)=(x_1,x_2,x_3,0).$$ We may now look upon the second sheet as the graph of a function $`\varphi :^3`$. The requirement that $`Y`$ is tangent to $`X`$ at $`0`$ implies that $`\varphi (0)=0`$ and $`\frac{\varphi }{y_i}(0)=0`$, $`i=1,2,3`$. Moreover, the second partials must be nondegenerate since the immersion is in the codimension 1 part of the discriminant hypersurface. Changing coordinates in $`^4`$, by adding a function of the first three coordinates that vanishes to order 3 to the fourth, we may assume that $`\varphi `$ is identical to the second order terms of its Taylor polynomial. Choosing the coordinates in $`Y`$ appropriately then gives $`f_0`$ the desired form. We must check infinitesimal versality. If $`(z_1,\mathrm{},z_4)`$ are coordinates on $`^4`$ then this amounts to showing that any smooth variations $$\alpha (x)=\underset{i_1}{\overset{4}{}}\alpha _i(x)\frac{}{z_i},\beta (x)=\underset{i_1}{\overset{4}{}}\beta _i(y)\frac{}{z_i},$$ can be written as $`\alpha (x)`$ $`={\displaystyle \underset{i=1}{\overset{3}{}}}a_i(x){\displaystyle \frac{f_0}{x_i}}+k(f_0(x))+c\dot{f}_t(x)\|_{t=0},`$ $`\beta (y)`$ $`={\displaystyle \underset{i=1}{\overset{3}{}}}b_i(y){\displaystyle \frac{f_0}{y_i}}+k(f_0(y))+c\dot{f}_t(y)\|_{t=0},`$ where $`k`$ is a vector field on $`^4`$ and $`c`$ is a constant. Writing these equations out we get the system $`\alpha _i(x)`$ $`=a_i(x)+k_i(x,0),i=1,2,3,`$ $`\alpha _4(x)`$ $`=k_4(x,0),`$ $`\beta _i(y)`$ $`=b_i(y)+k_i(y,y_1^2+y_2^2+ϵy_3^2),i=1,2,3,`$ $`\beta _4(y)`$ $`=2(y_1b_1(y)+y_2b_2(y)+y_3b_3(y))+k_4(y,y_1^2+y_2^2+ϵy_3^2)+c.`$ Let $`k_4(z)=\alpha _4(z_1,z_2,z_3)`$. Then the fourth equation holds. Let $`c=\beta _4(0)k_4(0)`$. Let $`\varphi (y)=\beta _4(y)k_4(f_0(y))c`$ and choose $`c`$ so that $`\varphi (0)=0`$. If $`y_i0`$ let $`b_1(y)`$ $`={\displaystyle \frac{1}{6y_1}}\left(\varphi (y_1,y_2,y_3)\varphi (0,y_2,y_3)+\varphi (y_1,0,y_3)\varphi (0,0,y_3)+\varphi (y_1,0,0)\right),`$ $`b_2(y)`$ $`={\displaystyle \frac{1}{6y_2}}\left(\varphi (y_1,y_2,y_3)\varphi (y_1,0,y_3)+\varphi (y_1,y_2,0)\varphi (y_1,0,0)+\varphi (0,y_2,0)\right),`$ $`b_3(y)`$ $`={\displaystyle \frac{1}{6y_3}}\left(\varphi (y_1,y_2,y_3)\varphi (y_1,y_2,0)+\varphi (0,y_2,y_3)\varphi (0,y_2,0)+\varphi (0,0,y_3)\right),`$ and if $`y_i=0`$ let $`b_1(y)`$ $`={\displaystyle \frac{1}{6}}\left({\displaystyle \frac{\varphi }{y_1}}(0,y_2,y_3)+{\displaystyle \frac{\varphi }{y_1}}(0,0,y_3)+{\displaystyle \frac{\varphi }{y_1}}(0,0,0)\right),`$ $`b_2(y)`$ $`={\displaystyle \frac{1}{6}}\left({\displaystyle \frac{\varphi }{y_2}}(y_1,0,y_3)+{\displaystyle \frac{\varphi }{y_2}}(y_1,0,0)+{\displaystyle \frac{\varphi }{y_2}}(0,0,0)\right),`$ $`b_3(y)`$ $`={\displaystyle \frac{1}{6}}\left({\displaystyle \frac{\varphi }{y_3}}(y_1,y_2,0)+{\displaystyle \frac{\varphi }{y_3}}(0,y_2,0)+{\displaystyle \frac{\varphi }{y_3}}(0,0,0)\right).`$ Then the last equation holds. Choosing first $`k_i(z)=k_i(z_1,z_2,z_3)`$, $`i=1,2,3`$ so that the three remaining $`\beta `$-equations hold and then choosing $`a_i`$, $`i=1,2,3`$ so that the remaining $`\alpha `$-equations hold, infinitesimal versality is proved. ## 6. Morse modifications, linking, and first order invariants In this section, we study how the self intersections of generic immersions and their preimages transform as the immersions cross the discriminant hypersurface. These transformations give rise to invariants of immersions: A function on $`_n`$, which is constant on path components will be called an invariant of regular homotopy. A function on $`_n\mathrm{\Sigma }_n`$ which is constant on path components of $`_n\mathrm{\Sigma }_n`$ is an invariant of generic immersions. In our study of invariants of generic immersion we are interested in how they change when we pass $`\mathrm{\Sigma }_n`$. This is described in terms of jumps: Let $`f_t:S^3^n`$ be a path in $`_n`$, $`n=4,5`$ intersecting $`\mathrm{\Sigma }_n^1`$ transversally at $`f_0`$. Let $`\delta >0`$ be small. Let $`v`$ be an invariant of generic immersions. Then $$v(f_0)=v(f_\delta )v(f_\delta ),$$ is a locally constant function on $`\mathrm{\Sigma }_n^1`$, defined up to sign. We call it the jump of $`v`$. In Section 6.3, we get rid of the sign ambiguity in the definition of $`v`$. We prove Theorem 2 and end the section with an axiomatic characterization of strangeness $`\mathrm{St}`$ for immersions $`S^3^5`$. ### 6.1. Morse modifications Throughout this section, let $`f_t`$ be a path in $`_n`$ intersecting $`\mathrm{\Sigma }_1`$ transversally at $`f_0`$ and let $`\delta 0`$ be small enough so that $`f_t`$ is generic for $`0<\|t\|\delta `$. Let $`n=5`$ and let $`f_0`$ have a self tangency point. By Proposition 5.3.2, the self intersection of $`f_\delta `$ is obtained from the self intersection of $`f_\delta `$ by a single Morse modification. Define $$J(f)=\text{the number of connected components of }M_f,$$ for generic immersions $`f`$. ###### Lemma 6.1.1. $`J`$ is an invariant of generic immersions $`S^3^5`$. It jumps by $`\pm 1`$ when crossing the self tangency part of $`\mathrm{\Sigma }_5^1`$ and remains constant when crossing the triple point part. ###### Proof. Immediate from Proposition 5.3.2. ∎ Let $`n=4`$ and let $`f_0`$ have an exceptional quadruple point. By Proposition 5.3.1 the number of quadruple points of $`f_\delta `$ differs from the number of quadruple points of $`f_\delta `$ by $`\pm 2`$. Define $$Q(f)=\text{the number of quadruple points of }f,\text{ and }Q_2(f)=Q(f)\mathrm{mod}2_2,$$ for generic immersions $`f`$. ###### Proposition 6.1.2. $`Q`$ is an invariant of generic immersions $`S^3^4`$. It jumps by $`\pm 2`$ when crossing the part of $`\mathrm{\Sigma }_4^1`$ that consists of immersions with an exceptional quadruple point and does not jump when crossing any other part of $`\mathrm{\Sigma }_4^1`$. Furthermore, $`Q_2`$ is an invariant of regular homotopy. ###### Proof. The first part is immediate from Proposition 5.3.1. Let $`f`$ and $`g`$ be regularly homotopic then $`f`$ and $`g`$ can be joined by a path in $`_4`$ which intersects $`\mathrm{\Sigma }_4^1`$ transversally. Since $`Q`$ jumps by $`\pm 2`$ or $`0`$ on such intersections, the lemma follows. ∎ Let $`n=4`$ and let $`f_0`$ have an exceptional triple point. By Proposition 5.3.1, $`C_{f_\delta }`$ is obtained from $`C_{f_\delta }`$ by a single Morse modification (for notation, see Section 5.1). Define $$T(f)=\text{the number of components of }C_f,$$ for generic immersions $`f`$. ###### Proposition 6.1.3. $`T`$ is an invariant of generic immersions $`S^3^4`$. It jumps by $`\pm 1`$ when crossing the part of $`\mathrm{\Sigma }_4^1`$ that consists of immersions with an exceptional triple point and does not jump when crossing any other part of $`\mathrm{\Sigma }_4^1`$. ###### Proof. Immediate from Proposition 5.3.1. ∎ Let $`n=4`$ and let $`f_0`$ have an exceptional double point (a self tangency point). By Proposition 5.3.1, $`F_{f_\delta }`$ is obtained from $`F_{f_\delta }`$ by a single Morse modification (for notation, see Section 5.1). Define $$D(f)=\chi (F_f)\text{ and }D_2(f)=D(f)\mathrm{mod}2_2,$$ for generic immersions $`f`$, where $`\chi `$ denotes the Euler characteristic. ###### Proposition 6.1.4. $`D`$ is an invariant of generic immersions $`S^3^4`$. It jumps by $`\pm 2`$ when crossing the part of $`\mathrm{\Sigma }_4^1`$ that consists of immersions with an exceptional double point and does not jump when crossing any other part of $`\mathrm{\Sigma }_4^1`$. Furthermore, $`D_2`$ is an invariant of regular homotopy. ###### Proof. Similar to the proof of Proposition 6.1.2. ∎ ### 6.2. Linking We construct an invariant of generic immersions $`S^3^5`$ which jumps under triple point moves and stays constant under self tangency moves: Let $`f:S^3^5`$ be a generic immersion with self intersection $`M_f`$. Consider the preimage $`\stackrel{~}{M}_f=f^1(M_f)`$. Choose a normal vector field $`w`$ along $`\stackrel{~}{M}_f`$ satisfying the following condition: If $`\stackrel{~}{M}_f^{}`$ denotes the result of pushing $`\stackrel{~}{M}_f`$ slightly along $`w`$ we require that $$[\stackrel{~}{M}_f^{}]=0H_1(S^3\stackrel{~}{M}_f).$$ ($``$) Any two vector fields satisfying this condition are homotopic. For existence, note that we can take $`w`$ as the normal vector field of the boundary in a Seifert-surface of the link $`\stackrel{~}{M}_fS^3`$. We define a vector field $`v`$ along $`M_f`$: For $`pM_f`$, let $$v(p)=dfw(p_1)+dfw(p_2),\text{ where }\{p_1,p_2\}=f^1(p).$$ Let $`M_f^{}`$ denote the result of pushing $`M_f`$ slightly along $`v`$. Then $`M_f^{}^5f(S^3)`$. By Alexander duality $`H_1(^5f(S^3))H^3(f(S^3))`$ and $`H^3(f(S^3))`$ since there is a triangulation of $`S^3`$ giving a triangulation of $`f(S^3)`$ after identifications in the $`0`$ and $`1`$ skeletons only. An orientation of $`S^3`$ gives a canonical generator of $`H_1(^5f(S^3))`$: The boundary of a small 2-disk intersecting $`f(S^3)`$ transversally in one point and oriented in such a way that the intersection number is positive. Note that this intersection number is independent of the ordering of the 2-disk and $`f(S^3)`$. Recall that there is an induced orientation on $`M_f`$ (Section 5.1) and hence on $`M_f^{}`$. Define $$\mathrm{lk}(f)=[M_f^{}]H_1(^5f(S^3))=\text{ and }\lambda (f)=\mathrm{lk}(f)\mathrm{mod}3_3.$$ Note that $`\mathrm{lk}(f)`$ is well defined since an homotopy between two vector fields $`w`$ and $`w^{}`$ satisfying ($``$) induces a homotopy between the shifted self intersections in $`^5f(S^3)`$. ###### Lemma 6.2.1. $`\mathrm{lk}`$ is an invariant of generic immersions. It jumps by $`\pm 3`$ when crossing the triple point part of $`\mathrm{\Sigma }_5^1`$ and remains constant when crossing the self tangency part. Furthermore, $`\lambda `$ is an invariant of regular homotopy. ###### Proof. The second part follows from the first exactly as in Proposition 6.1.2. We prove the first part: Suppose that $`f_t`$ is a path in $`_5`$ intersecting $`\mathrm{\Sigma }_5^1`$ transversally at $`f_0`$ and let $`\delta >0`$ be small. We restrict attention to a small neighborhood of the exceptional multiple point of $`f_0`$ since all the immersions $`f_t`$, $`\|t\|\delta `$ can be assumed to agree outside of this neighborhood. Assume that $`f_0`$ has an elliptic self tangency point $`f_0(p)=f_0(q)`$. Using Proposition 5.3.2 we can write $`f_t(x)`$ $`=(x_1,x_2,t,x_3,0).`$ $`f_t(y)`$ $`=(y_1,y_2,y_1^2+y_2^2,0,y_3),`$ where $`x`$ are coordinates around $`p`$ and $`y`$ are coordinates around $`q`$. The preimages of the newborn self intersection circle $`c`$ of $`f_\delta `$ are $`\{x_1^2+x_2^2=\delta ,x_3=0\}`$ and $`\{y_1^2+y_2^2=\delta ,y_3=0\}`$. We may assume that the field $`w`$ is given by $`w(x)=\frac{}{x_3}`$ and $`w(y)=\frac{}{y_3}`$. Now shift $`c`$ a small distance $`ϵ`$ along $`df_\delta w(x)+df_\delta w(y)`$ to get $`c^{}`$. Then $`c^{}`$ is the boundary of the 2-disk $$(r\delta \mathrm{cos}(v),r\delta \mathrm{sin}(v),\delta ,ϵ,ϵ),0v2\pi ,\mathrm{\hspace{0.17em}0}r1,$$ which does not intersect $`f_\delta (S^3)`$. Hence, $`\mathrm{lk}(f_\delta )=\mathrm{lk}(f_\delta )`$. Assume that $`f_0`$ has a hyperbolic self tangency point. As above we can write $`f_t(x)`$ $`=(x_1,x_2,t,x_3,0).`$ $`f_t(y)`$ $`=(y_1,y_2,y_1^2y_2^2,0,y_3).`$ The preimages of the self intersection are $`\{x_1^2x_2^2=\delta ,x_3=0\}`$ and $`\{y_1^2y_2^2=\delta ,y_3=0\}`$. We can choose the field $`w`$ so that $`w(x)=\frac{}{x_3}`$ and $`w(y)=\frac{}{y_3}`$ close to $`p`$ and $`q`$, respectively. We note that with vector fields as above and equal outside neighborhoods of $`p`$ and $`q`$ the condition ($``$) holds for $`f_\delta `$ if and only if it holds for $`f_\delta `$. The rest of the argument is similar to the elliptic case. Assume that $`f_0`$ has a triple point. According to Proposition 5.3.2 we can find coordinates $`x`$, $`y`$ and $`z`$ centered at $`p_1`$, $`p_2`$ and $`p_3`$, respectively such that close to the triple point $`f_0(p_1)=f_0(p_2)=f_0(p_3)`$ we have $`f_t(x)`$ $`=(x_1,x_2,x_3,0,0),`$ $`f_t(y)`$ $`=(y_1,t,0,y_2,y_3),`$ $`f_t(z)`$ $`=(0,z_2,z_1,z_2,z_3).`$ Denote the neighborhoods of $`p_1`$, $`p_2`$ and $`p_3`$ by $`X`$, $`Y`$ and $`Z`$, respectively. If we orient these by declaring the frames $`(_1,_2,_3)`$ to have the positive orientation then the oriented self intersections of $`f_0`$ are the lines $`f_t(X)f_t(Y)`$ $`=(s,t,0,0,0),`$ $`s,`$ $`f_t(X)f_t(Z)`$ $`=(0,0,s,0,0),`$ $`s,`$ $`f_t(Z)f_t(Y)`$ $`=(0,t,0,t,s),`$ $`s.`$ The inverse image of the self intersection of $`f_t`$, $`t<0`$ is shown in Figure 2. We calculate $`\mathrm{lk}(f_\delta )`$ in terms of $`\mathrm{lk}(f_\delta )`$ and have to take two things into consideration: First the motion of $`Y`$ itself. Second, the changes in the field $`w`$ that this motion causes. Assume that the shifting distance is very small in comparison to $`\delta `$. Let $`l_t(X,Z)`$ denote the intersection $`f_t(X)f_t(Z)`$ shifted along $`v=dfw(x)+dfw(z)`$. Let $`D_t(X,Z)`$ denote a part of a disc in $`^5`$ bounded by $`l_t(X,Z)`$. We may assume that $`D_t(X,Z)`$ is a shift along an extension of $`v`$ of a disk in $`f_t(X)`$. First, the intersection $`D_t(X,Z)f_t(Y)`$ does not depend on the field $`w`$ since the shift is assumed to be very small in comparison to $`t`$. We calculate the change in $`D_t(X,Z)f_t(Y)`$ from the preimage in $`X`$. From Figure 3 it follows that the algebraic number of intersection points in $`D_\delta (X,Z)f_\delta (Y)`$ differs from that corresponding to $`D_\delta (X,Z)f_\delta (Y)`$ by $`+1`$. We now take the field $`w`$ into consideration. In Figure 4 we see a Seifert-surface in $`S^3`$ of $`\stackrel{~}{M}_f`$ close to a crossing point. In Figure 5 we see $`w`$ chosen as the inward normal in a Seifert-surface of $`\stackrel{~}{M}_{f_t}`$. As we move through the triple point the direction of rotation of the vector field $`w`$ is changed. This change in rotation gives rise to a new positive intersection point in $`D_t(X,Z)f_t(Z)`$ for $`t>0`$. As seen in Figure 6. Similarly, we get a new positive intersection point in $`D_t(X,Z)f_t(X)`$. Thus, in total the algebraic number of intersection points in $`D_t(X,Z)(f_t(X)f_t(Y)f_t(Z))`$ increases by $`1`$ when $`t`$ is changed from $`\delta `$ to $`\delta `$. Similarly, the intersection numbers corresponding to $`D_t(X,Y)(f_t(X)f_t(Y)f_t(Z))`$ and $`D_t(Y,Z)(f_t(X)f_t(Y)f_t(Z))`$ increase by $`1`$ when $`t`$ is changed from $`\delta `$ to $`\delta `$. It follows that $`\mathrm{lk}(f_\delta )\mathrm{lk}(f_\delta )=3`$. This proves the lemma. ∎ ###### Corollary 6.2.2. For $`\xi `$ in $`\mathrm{𝐈𝐦𝐦}`$, let $`\lambda (\xi )=\lambda (f)`$, where $`f`$ is a generic immersion representing $`\xi `$. Then $$\lambda :\mathrm{𝐈𝐦𝐦}_3,$$ is a homomorphism. ###### Proof. As a function, $`\lambda `$ is well defined by Lemma 6.2.1. Clearly, it is additive under connected sum and thus a homomorphism. ∎ ###### Remark 6.2.3. In knot theory one assigns a sign to crossings as in Figure 2 by comparing the orientation given by $`(v_1,v_2,v_{12})`$ to that of the ambient space, where $`v_1`$ is the tangent vector to the first branch, $`v_2`$ the tangent vector of the second, and $`v_{12}`$ a vector from the second to the first branch. We note that all crossings appearing in Figure 2 are positive. Changing the orientation of one of the sheets changes the sign of all crossings. Hence, all three crossings appearing close to a triple point always have the same sign. ### 6.3. Finite order invariants In this section we summarize the properties of finite order invariants of generic immersions that are needed to prove Theorem 2. Recall that we have defined the jump $`v`$ of an invariant $`v`$ of generic immersions. It was defined only up to sign. To get rid of this sign we need a coorientation of $`\mathrm{\Sigma }_n^1`$: If $`f_0\mathrm{\Sigma }_n^1`$ then there is a neighborhood $`U`$ of $`f_0`$ in $`_n`$ which is cut into two parts by $`\mathrm{\Sigma }_n^1`$. We make a choice of a positive part and a negative part of $`U`$. A coherent choice like that for all $`f_0\mathrm{\Sigma }_n^1`$ is a coorientation of $`\mathrm{\Sigma }_n`$. A coorientation enables us to make $`v`$ well defined: In the definition $$v(f_0)=v(f_\delta )v(f_\delta ),$$ we require that $`f_\delta `$ is on the positive and $`f_\delta `$ on the negative side of $`\mathrm{\Sigma }_n^1`$ at $`f_0`$. As mentioned, $`v`$ is then locally constant on $`\mathrm{\Sigma }_n^1`$ and we may consider $`v`$ as an element of $`H^0(\mathrm{\Sigma }_n^1)`$. Let $`f_t`$ be a path in $`_5`$ intersecting $`\mathrm{\Sigma }_5^1`$ transversally at $`f_0`$. We coorient $`\mathrm{\Sigma }_5`$ as follows: If $`f_0EY`$ then $`f_\delta `$ is on the positive side of $`EY`$ at $`f_0`$ if $`M_{f_\delta }`$ has more components than $`M_{f_\delta }`$. If $`f_0T`$ then $`f_\delta `$ is on the positive side of $`T`$ if $`\mathrm{lk}(f_\delta )\mathrm{lk}(f_\delta )`$. It is easy to check that this coorientation is continuous (see , Section 7.6). This enables us to iterate the above construction and define inductively $$^{k+1}v(f_{0,0,\mathrm{},0})=^k(f_{\delta ,0,\mathrm{},0})^k(f_{\delta ,0,\mathrm{},0}).$$ Here $`f_{0,\mathrm{},0}`$ is an immersion with $`k+1`$ distinct degenerate points (self tangencies or triple points) and $`f_{t,0,\mathrm{},0}`$ is a path in $`\mathrm{\Sigma }_5^k\mathrm{\Sigma }_5^{k+1}`$, where $`\mathrm{\Sigma }_5^j`$ is the space of immersions with $`j`$ distinct degenerate points, intersecting $`\mathrm{\Sigma }_5^{k+1}`$ transversally. Then $`^kv`$ is an element in $`H^0(\mathrm{\Sigma }_5^k)`$ (see , Remark 9.1.2). We say that an invariant $`v`$ is of order $`k`$ if $`^{k+1}v0`$. Finally, we remark that the space of invariants of generic immersions, i.e.$`H^0(_n\mathrm{\Sigma }_n)`$, splits as a direct sum over the path components of $`_n`$. That is $`H^0(_n\mathrm{\Sigma }_n)=_UH^0(U(\mathrm{\Sigma }_nU))`$, where $`U`$ runs over the path components of $`_n`$. ### 6.4. Proof of Theorem 2 By Lemmas 6.1.1 and 6.2.1, $`J`$ and $`L`$ are invariants which changes as claimed. Clearly, $`L`$ and $`J`$ are independent. Since $`J`$ ($`L`$) is $`1`$ ($`0`$) on all self tangency parts of $`\mathrm{\Sigma }_5^1`$ and is $`0`$ ($`1`$) on all triple point parts of $`\mathrm{\Sigma }_5^1`$, it follows from Proposition 5.4.1 (see Figure 1) that $`^2J0`$ and $`^2L0`$. Thus, they both have order one. Let $`U`$ be a path component of $`_5`$ and $`v`$ an invariant, the jump of which is constant on $`TU`$ and $`(EY)U`$ (for notation see Section 5.3). Then the jump of $`v\|U`$ is a linear combination of the jumps of $`J\|U`$ and $`L\|U`$ and it follows that $`v\|U`$ is (up to a constant term) a linear combination of $`J\|U`$ and $`L\|U`$. Thus, the theorem follows once we show that any first order invariant has this property. Let $`f`$ and $`g`$ be regularly homotopic immersions with one self tangency point of the same kind each or with one triple point each. Using a diffeotopy of $`^5`$ we can move the exceptional point of $`g`$ to that of $`f`$ and using a diffeotopy of $`S^3`$ we can make the maps agree in small disks around the preimages of this point. Choose unknotted paths in $`S^3`$ connecting these disks. Since there are no 1-knots in $`^5`$, we can after diffeotopy of $`^5`$ assume that the maps agree also in a neighborhood of these arcs. After this is done we have two immersions that agree on a 3-disk in $`S^3`$ and are regularly homotopic. We must show that they are regularly homotopic keeping this 3-disk fixed. However, by the Smale-Hirsch h-principle the obstruction to this regular homotopy is the difference of the Smale invariants which is zero. Choosing a generic regular homotopy we see that we can join $`f`$ to $`g`$ with a path in $`\mathrm{\Sigma }_5^1\mathrm{\Sigma }_5^2`$ intersecting $`\mathrm{\Sigma }_5^2`$ transversally only at immersions with two distinct exceptional points. Let $`v`$ be a first order invariant. Then $`^2v0`$. This and the fact that $`v`$ is locally constant on $`\mathrm{\Sigma }_5^1`$ implies that $`v`$ remains constant along a path as above. Hence, if $`f`$ and $`g`$ both belong to either one of $`EU`$, $`YU`$, or $`TU`$ then $`v(f)=v(g)`$. It remains to show that $`v\|EU`$ equals $`v\|YU`$: Consider the last part of Figure 1 in Section 5.4. Let $`f`$ be an immersion in the $`E`$-branch and $`g`$ an immersion in the $`Y`$-branch. Pick a small loop $`h_t`$, $`0t1`$, $`h_1=h_0`$ around the cusp such that it intersects the $`E`$-branch in $`f`$ and the $`Y`$-branch in $`g`$. Then $`0=v(h_1)v(h_1)=v(f)v(g)`$. Hence, $$v\|EUv(f)=v(g)v\|YU.$$ The theorem follows. ∎ ### 6.5. $`L`$ and $`\mathrm{St}`$ Arnold’s invariant $`\mathrm{St}`$ for plane curves (see ) can be characterized axiomatically as an invariant which jumps by $`1`$ on triple points, does not jump under self tangencies, and is additive under connected summation. The invariant $`L=\frac{1}{3}(\mathrm{lk}+\stackrel{~}{\lambda })`$ jumps by one on triple points but is not additive under connected sum. The reason for this is that the invariant $`\stackrel{~}{\lambda }`$ of order zero is not additive under connected sum. The invariant $`\mathrm{lk}`$ is additive under connected sum. It jumps by $`\pm 3`$ on $`T`$ and not at all on $`EY`$. Thus, $`\mathrm{lk}`$ and $`J`$ generates the space of $``$-valued invariants. We want an invariant which is additive under connected sum and together with $`J`$ generates the space of $``$-valued invariants. To accomplish this we note that, either $`\mathrm{lk}+\mathrm{\Omega }`$ or $`\mathrm{lk}\mathrm{\Omega }`$ is divisible by $`3`$, where $`\mathrm{\Omega }`$ is the Smale invariant, which we may consider integer valued after the choice of generator of $`\pi _3(V_{5,3})`$ made in Lemma 3.1.2. In Lemma 8.3.3 we calculate $`\mathrm{lk}(f)`$ for an immersion with $`\mathrm{\Omega }(f)=1`$ and the result is $`\mathrm{lk}(f)=2`$. Since $`\lambda =\mathrm{lk}\mathrm{mod}3`$ is an invariant of regular homotopy it follows that $`\frac{1}{3}(\mathrm{lk}(f)+\mathrm{\Omega }(f))`$ is an integer valued invariant which jumps by $`\pm 1`$ on triple points, does not jump on self tangencies, and is additive under connected sum (since both $`\mathrm{lk}`$ and $`\mathrm{\Omega }`$ are). We therefore define $$\mathrm{St}(f)=\frac{1}{3}((\mathrm{lk}(f)+\mathrm{\Omega }(f))).$$ To complete the analogy with Arnold’s $`\mathrm{St}`$, we need to define the coorientation of the triple point stratum $`T`$ of $`\mathrm{\Sigma }_5^1`$ in terms of the local picture near a triple point. Recall that all crossings in the preimage close to a triple point have the same sign (Remark 6.2.3). By Lemma 6.2.1 our coorientation of $`T`$ agrees with the following one defined in local terms: We say that we pass $`T`$ in the positive direction if the crossing signs in the preimage close to the triple point changes from positive to negative. ###### Proposition 6.5.1. There is a unique invariant $`\mathrm{St}`$ such that it jumps by $`+1`$ ($`1`$) under positive (negative) triple point moves, it does not jump under self tangency moves, it is additive under connected sum, its value on an immersion changes sign if the immersion is composed with an orientation reversing diffeomorphism of $`S^3`$, and takes the values $`1`$ on $`f_{\frac{1}{2}}`$ (see Proposition 8.3.1). ###### Proof. A generic immersion $`g`$ with $`\mathrm{\Omega }(g)0`$ is regularly homotopic to a connected sum of $`\|\mathrm{\Omega }(g)\|_+`$ copies of $`f_{\frac{1}{2}}`$ or $`f_{\frac{1}{2}}r`$, where $`r`$ is an orientation reversing diffeomorphism. A generic immersion $`g`$ with $`\mathrm{\Omega }(g)=0`$ is regularly homotopic to $`f_{\frac{1}{2}}f_{\frac{1}{2}}`$. We know the values of our invariants on connected sums of $`f_{\pm \frac{1}{2}}`$. Thus, we know the value on any other generic immersion by adding the jumps along a path connecting it to such a connected sum. This proves uniqueness. For existence, we need to show that $`\mathrm{St}`$ changes sign if the orientation is reversed. For $`\mathrm{\Omega }`$ we already know this (Lemma 3.4.2). To see that the same is true for $`\mathrm{lk}`$ note that the induced orientation of $`M_f`$ does not change if we reverse the orientation on $`S^3`$. The orientation of $`f(S^3)`$ does change and hence $`\mathrm{lk}(f)=[M_f^{}]H_1(^5f(S^3))`$ does change sign. ∎ ## 7. Pin structures and twist framings In this section we show that there is an induced pin structure on the self intersection surface of a generic immersion $`S^3^4`$. The Brown invariant of this pin structure is unchanged under regular homotopy. Actually an even stronger result holds, if two generic immersions into $`^4`$ are regularly homotopic in $`^5`$ after composing them with the inclusion $`^4^5`$ then the corresponding Brown invariants are equal. Using the geometry of self intersections of generic immersions $`S^3^5`$, we define a $`_4`$-valued invariant of regular homotopy. ### 7.1. Immersions into $`^4^5`$ Let $`g:S^3^4`$ be an immersion. Its Smale invariant $`\mathrm{\Omega }(g)`$ is an element of $`\pi _3(SO(4))=[\sigma ][\rho ]`$ (see Lemma 3.1.1). The normal bundle of $`g`$ is 1-dimensional and orientable. Hence, it is trivial and $`g`$ admits a normal field $`n:S^3S^3`$, such that the frame $`(n(x),dfQ(x))`$ (see Section 3.3) gives the positive orientation of $`^4`$. If $`\mathrm{\Omega }(g)=m[\sigma ]+n[\rho ]`$ then, from Lemma 3.1.1, it follows that $`m+1=\mathrm{deg}(n)`$, the normal degree of $`g`$. Let $`s:S^3^4`$ be the standard embedding and $`\widehat{s}=sr`$, where $`r`$ is a reflection in a hyperplane as in Lemma 3.4.2. ###### Lemma 7.1.1. $$\mathrm{\Omega }(\widehat{s})=2[\sigma ]+[\rho ]\pi _3(V_{4,3})=\pi _3(SO(4)).$$ ###### Proof. Let $`i:^4^5`$ be the inclusion. Then $`is`$ and $`i\widehat{s}`$ are regularly homotopic in $`^5`$. Hence, by Lemma 3.3.3, $`\mathrm{\Omega }(\widehat{s})NSO(4)`$. The normal degree of $`\widehat{s}`$ equals the normal degree of $`s`$ with opposite sign. Thus, $`\mathrm{\Omega }_4(\widehat{s})=2[\sigma ]+n[\rho ]`$. But then, $`\mathrm{\Omega }_4(\widehat{s})NSO(4)`$ implies that $`n=1`$. ∎ Let $`g,h:S^3^4`$ be immersions. ###### Proposition 7.1.2. If $`ig`$ and $`ih`$ are regularly homotopic in $`^5`$ then either $`g`$ and $`h\widehat{s}\mathrm{}\widehat{s}`$ or $`h`$ and $`g\widehat{s}\mathrm{}\widehat{s}`$ are regularly homotopic in $`^4`$. ###### Proof. If $`ig`$ and $`ih`$ are regularly homotopic in $`^5`$ then $`\mathrm{\Omega }(g)\mathrm{\Omega }(h)N`$. Now, $`\mathrm{\Omega }(\widehat{s})`$ is a generator of $`N`$ and the Smale invariant is additive under connected sum. ∎ ### 7.2. Functionals on curves in self intersections The invariants we are about to define originates from functionals on curves in self intersections. Before we can construct these functionals we need some preliminaries: Let $`𝐈`$ denote the identity matrix. ###### Definition 7.2.1. Define $`𝐀(4)SO(4)`$ and $`𝐀(5)SO(5)`$ by $$𝐀(4)=\left[\begin{array}{cccc}1& 0& 0& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right],𝐀(5)=\left[\begin{array}{ccccc}1& 0& 0& 0& 0\\ 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 1\\ 0& 1& 0& 0& 0\\ 0& 0& 1& 0& 0\end{array}\right].$$ The group $`\{𝐈,𝐀(n)\}`$ acts on $`SO(n)`$ from the right. Denote the corresponding quotient space $`SO(n)/𝐀(n)`$ and let $`p`$ be the projection. ###### Proposition 7.2.2. For $`n=4`$ and $`n=5`$, $$\pi _1(SO(n)/𝐀(n))=_4.$$ ###### Proof. See , Proposition 5.2. ∎ Let $`u\pi _1(SO(n))`$ denote the nontrivial element. The functional on curves in a self intersection surface of a generic immersion into 4-space is constructed as follows (for notation, see Proposition 5.1.5): Let $`g:S^3^4`$ be a generic immersion. Let $`c`$ be a closed curve in $`F_g`$ meeting $`F_g^0F_g^1`$ transversally and let $`c_S=p^1(c)`$. Note that $`c_S`$ is either a union of two disjoint circles or one circle. Choose a parameterization $`r:Ic`$ of $`c`$ and a normal vector $`n:ITF_g`$ of $`c`$ in $`TF_g`$. Then $`r`$ lifts to two parameterizations $`r_1`$ and $`r_2`$. If $`c_S`$ is connected then the path product $`r_1r_2`$ is a parameterization of $`c_S`$. Choose normal fields $`\nu _1`$ and $`\nu _2`$ of $`\stackrel{~}{M}_g`$ along $`sr_1`$ and $`sr_2`$, respectively. Let $`n_1`$ and $`n_2`$ be such that $`dp(n_i)=n`$. Assume that $`\nu _i`$ is chosen so that $`(ds\dot{r}_i,dsn_i,\nu _i)`$ is a positively oriented frame, $`i=1,2`$. Then $`(dt\dot{r},dtn,dg\nu _1,dg\nu _2)`$ is a framing or twist framing $`X`$ of $`T^4`$ along $`tc`$. Let $`\stackrel{}{c}_S`$ denote the loop (or loops) in $`SO(TS^3)`$ (the principal $`SO(3)`$-bundle of the tangent bundle of $`S^3`$) represented by $`sc_S`$ with the framing $`(ds\dot{r}_i,dsn_i,\nu _i)`$ along $`sr_i`$, $`i=1,2`$ and let $`[\stackrel{}{c}_S]`$ denote the corresponding element in $`H_1(SO(TS^3);_2)`$. Let $`[c,r,X]\pi _1(SO(4)/𝐀(4))`$ denote the homotopy class of the loop induced by $`tc`$ with (twist) framing $`X`$ as above (which we may have to orthonormalize). Then define ###### Definition 7.2.3. $$\omega _g(c)=\xi _S,[\stackrel{}{c}_S]p_{}(u)+[c,r,X]+p_{}(u)\pi _1(SO(4)/𝐀(4)),$$ where $`\xi _SH^1(SO(TS^3);_2)`$ denotes the unique spin structure on $`S^3`$. This definition is independent of all choices (see , Lemma 7.6). We extend the functional $`\omega _g`$ to collections $`L=\{c_1,\mathrm{},c_m\}`$ of oriented closed curves in $`F_g`$ transversal to $`F_g^0F_g^1`$: $$\omega _g(L)=\underset{i=1}{\overset{m}{}}\omega _g(c_i).$$ Let $`L`$ be a collection of oriented closed curves in $`F_g`$ with $`m`$ transverse intersections and $`L^{}`$ be the collection obtained from smoothing each intersection, respecting orientation. Then $$\omega _g(L^{})=\omega _g(L)+mp_{}(u),$$ ($``$) and if $`K`$ and $`L`$ are isotopic collections then $$\omega _g(K)=\omega _g(L),$$ ($``$) see , Lemma 7.15. The functional on self intersection circles of generic immersions into $`5`$-space is constructed as follows: Let $`f:S^3^5`$ be a generic immersion. Let $`cM_f`$ be a component and $`c_S=f^1(c)`$. Note that $`c_S`$ is either a union of two disjoint circles or one circle. Choose a parameterization $`r`$ of $`c`$, agreeing with the orientation induced on $`c`$ if $`k`$ is odd. Then $`r`$ lifts to two parameterizations $`r_1`$ and $`r_2`$. If $`c_S`$ is connected then the path product $`r_1r_2`$ is a parameterization of $`c_S`$. Choose orthonormal framings $`Y_1`$ and $`Y_2`$ of $`N(c_SS^3)`$ (the normal bundle of $`c_S`$ in $`S^3`$) along $`r_1`$ and $`r_2`$ respectively. Then $`(\dot{r},dfY_1,dfY_2)`$ is a framing or twist framing $`X`$ of $`T^5`$ along $`c`$. Let $`\stackrel{}{c}_S`$ denote the loop (or loops) in $`SO(TS^5)`$ (the principal $`SO(5)`$-bundle of the tangent bundle of $`S^5`$) represented by $`c_S`$ with the framing $`(\dot{r}_i,Y_i)`$ and $`[\stackrel{}{c}_S]`$ the corresponding element in $`H_1(SO(TS^5);_2)`$. Let $`[c,r,X]\pi _1(SO(5)/𝐀(5))`$ denote the homotopy class of the loop induced by $`c`$ with the (twist) framing $`X`$ (which we may have to orthonormalize). Then define ###### Definition 7.2.4. $$\omega _f(c)=\xi ,[\stackrel{}{c}_S]p_{}(u)+[c,r,X]+p_{}(u)\pi _1(SO(5)/𝐀(5)),$$ where $`\xi H^1(SO(TS^5);_2)`$ denotes the unique spin structure on $`S^5`$. This definition is independent of all choices (see , Lemma 6.5). ### 7.3. Pin structures on surfaces Let $`V`$ be a vector space over $`_2`$ with a nonsingular symmetric bilinear form $`(x,y)xy`$. A $`_4`$-quadratic function on $`V`$ is a function $`q:V_4`$ such that $`q(x+y)=q(x)+q(y)+2(xy)`$, for all $`x,yV`$. There are four indecomposable $`_4`$-quadratic spaces: $`𝒫_+=(_2(a),,q),`$ $`aa=1,`$ $`q(a)=1,`$ $`𝒫_{}=(_2(a),,q),`$ $`aa=1,`$ $`q(a)=1,`$ $`𝒯_0=(_2(b)_2(c),,q),`$ $`bb=cc=0,bc=1,`$ $`q(b)=q(c)=0,`$ $`𝒯_4=(_2(b)_2(c),,q),`$ $`bb=cc=0,bc=1,`$ $`q(b)=q(c)=2.`$ A $`_4`$-quadratic space $`V`$ is called split if it contains a subspace $`H`$ such that $`dimH=\frac{1}{2}dimV`$, $`q\|H=0`$ and $`HH=\{0\}`$. Two $`_4`$-quadratic spaces $`V`$ and $`W`$ belong to the same Witt class if there exists spilt spaces $`S_1`$ and $`S_2`$ such that $`VS_1WS_2`$. The Witt classes forms the Witt group which is generated by $`[𝒫_+]`$ (the Witt class of $`𝒫_+`$) with relations $`8[𝒫_+]=0`$, $`4[𝒫_+]=[𝒯_4]`$, and $`[𝒫_+]+[𝒫_{}]=0`$. Given a $`_4`$-quadratic space $`V`$ with quadratic function $`q`$ we define $$\lambda (V,q)=\underset{xV}{}e^{\frac{\pi i}{2}q(x)}.$$ Then $$\lambda (V,q)=\sqrt{2}^{dimV}\left(\frac{1+i}{\sqrt{2}}\right)^m.$$ Since $`\frac{1+i}{\sqrt{2}}`$ is an $`8^{th}`$ root of unity, $`m`$ modulo 8 is well defined. This is Brown’s invariant. It is denoted $`\beta (V,q)`$ and gives an isomorphism between the Witt group and $`_8`$. We shall sometimes write $`\beta (q)`$, dropping $`V`$. More details about $`_4`$-quadratic spaces can be found in . Let $`F`$ be a surface. By a pin structure on $`F`$ we shall mean a $`Pin^{}`$-structure on the tangent bundle $`TF`$ of $`F`$. There is a 1-1 correspondence between pin structures on $`F`$ and $`_4`$-quadratic functions on $`H_1(F;_2)`$, see , Theorem 3.2. ### 7.4. Invariance of the Witt class and a homomorphism $`\mathrm{𝐈𝐦𝐦}_8`$ Let $`g:S^3^4`$ be a generic immersion. Choose, once and for all an isomorphism $$\varphi :\pi _1(SO(4)/𝐀(4))_4.$$ (This choice is discussed in Section 9.2.) By equations ($``$) and ($``$) in Section 7.2 and Lemma 3.4 in , $`\varphi \omega _g`$ induces a $`_4`$-quadratic function $`q_g:H_1(F_g;_2)_4`$ (and hence a pin structure on $`TF_g`$). We define $$\beta (g)=\beta (q_g),$$ the Brown invariant of the $`_4`$-quadratic function $`q_g`$. ###### Lemma 7.4.1. Let $`g_0,g_1:S^3^4`$ be regularly homotopic generic immersions. Then $`\beta (g_0)=\beta (g_1)`$ ###### Proof. Connect $`g_0`$ to $`g_1`$ by a path $`g_t`$ in $`_4`$ which is transversal to $`\mathrm{\Sigma }_4`$. This means that $`g_t`$ intersects $`\mathrm{\Sigma }_4`$ only transversally at finitely many points in $`\mathrm{\Sigma }_4^1`$. Clearly, $`\beta `$ remains constant as long as we do not cross $`\mathrm{\Sigma }_4`$ and we must show that $`\beta `$ remains unchanged when we cross $`\mathrm{\Sigma }_4^1`$: Passing $`\mathrm{\Sigma }_4^1`$ at $`g_{t_0}`$, where $`g_{t_0}`$ has a $`j`$-tuple point $`j2`$ does not change $`F_{g_t}`$. It changes the preimage in $`F_g`$ of $`M_g^{4j}`$. Since the curves representing a basis of $`H_1(F_{g_t};_2)`$, used to calculate $`\beta `$ can be chosen so that they do not meet the discs where these changes occur, $`\beta `$ does not change at such crossings. Passing $`\mathrm{\Sigma }_4^1`$ at $`g_{t_0}`$, where $`g_{t_0}`$ has a degenerate double point does change $`F_{g_t}`$ by a Morse modification. If it has index $`2`$ or $`0`$ then $`F_{g_t}`$ is changed by addition or subtraction of an $`S^2`$-component. This does not affect $`\beta `$. If it has index $`1`$ then $`F_{g_t}`$ is changed by addition or subtraction of a handle. In this handle there is a newborn circle $`c`$, $`q_f(c)=0`$ and $`\beta `$ remains unchanged, see , Lemmas 5.2.1-8. ∎ Let $`f:S^3^5`$ be an immersion. The normal bundle of $`f`$ is trivial (Section 3.3) which implies (, Theorem 6.4) that $`f`$ is regularly homotopic to an immersion $`f_1:S^3^4`$ composed with the inclusion $`^4^5`$. Moreover, we may assume that $`f_1`$ is generic. Define $$\beta (f)=\beta (f_1).$$ ###### Proposition 7.4.2. For $`\xi \mathrm{𝐈𝐦𝐦}`$, let $`\beta (\xi )=\beta (f)`$, where $`f`$ is an immersion representing $`\xi `$. Then $$\beta :\mathrm{𝐈𝐦𝐦}_8,$$ is a homomorphism. ###### Proof. We show first that $`\beta `$ is well defined: If $`f_1`$ and $`g_1`$ are immersions into $`^4`$ which are regularly homotopic to $`f`$ in $`^5`$ then $`f_1`$ is regularly homotopic to $`g_1\widehat{s}\mathrm{}\widehat{s}`$ (Proposition 7.1.2). We can perform the connected sum of two generic immersions into $`^4`$ in such a way that the self intersection of the immersion obtained is the union of the self intersections of the summands and at most two new $`S^2`$ self intersection components. Hence, performing connected sum with the embedding $`\widehat{s}`$ does not change $`\beta `$. Thus, $`\beta :\mathrm{𝐈𝐦𝐦}_8`$ is well defined by Lemma 7.4.1. By the argument above, $`\beta `$ is additive under connected summation. Hence, it is a homomorphism. ∎ ### 7.5. A homomorphism $`\mathrm{𝐈𝐦𝐦}_4`$ Let $`f:S^3^5`$ be a generic immersion. Choose, once and for all, an isomorphism $$\psi :\pi _1(SO(5)/𝐀(5))_4.$$ (This choice is discussed in Section 9.2.) ###### Definition 7.5.1. A component $`c`$ of the self intersection $`M_f`$ is right twist framed if $`\psi (\omega _f(c))=1`$. It is left twist framed if $`\psi (\omega _f(c))=1`$. If $`c_S`$ is disconnected then there is an induced spin structure on $`c`$, see , Proposition 2.10. This spin structure on $`c`$ is trivial (spin-cobordant to zero) if $`\omega _f(c)=0`$ and nontrivial if $`\omega _f(c)=p_{}(u)`$ (see , proof of Theorem 7.30). We define $`r(f)`$, $`l(f)`$ and $`n(f)`$ as the number of self intersection components of a generic immersion $`f`$ that are right twist framed, left twist framed and has the nontrivial spin structure, respectively. Define $$\tau (f)=r(f)+2n(f)l(f)\mathrm{mod}4_4$$ ###### Proposition 7.5.2. For $`\xi \mathrm{𝐈𝐦𝐦}`$, let $`\tau (\xi )=\tau (f)`$, where $`f`$ is a generic immersion representing $`\xi `$. Then $$\tau :\mathrm{𝐈𝐦𝐦}_4,$$ is a homomorphism. ###### Proof. That $`\tau `$ is well defined follows from , Propositions 5.1.1-4. Clearly, $`\tau `$ is additive under connected sum. Hence, it is a homomorphism. ∎ ## 8. Special immersions In this section, we construct immersions $`S^3^5`$ with arbitrary Smale invariant by rotating a 2-disk with a kink i 5-space and use an immersion $`P^3^4`$ to show that the homomorphism $`\beta `$ (see Proposition 7.4.2) is nontrivial. ### 8.1. The Whitney kink. As in Whitney we consider the map $`g:^2^4`$ given by the equations $`g_1(x,y)`$ $`=x{\displaystyle \frac{2x}{u}}`$ $`g_2(x,y)`$ $`=y`$ $`g_3(x,y)`$ $`={\displaystyle \frac{1}{u}}`$ $`g_4(x,y)`$ $`={\displaystyle \frac{xy}{u}},`$ where $`u=(1+x^2)(1+y^2)`$. This is an immersion with one transversal double point: $$(0,0,1,0)=g(1,0)=g(1,0).$$ Moreover, if $`\|x\|`$ or $`\|y\|`$ is large then $`g`$ is close to the standard embedding $`^2^2\times 0^4`$. Thus, we can change $`g`$ slightly so that it agrees with the standard embedding outside some large disk. The Whitney kink enjoys the following symmetry property: If $`L:^4^4`$ and $`R:^2^2`$ are the linear maps $$L(z_1,z_2,z_3,z_4)=(z_1,z_2,z_3,z_4)\text{and}R(x,y)=(x,y),$$ then $`LgR=g`$. The differential of $`g`$ is $$dg_{(x,y)}=\left[\begin{array}{cc}1\frac{2(1x^2)}{(1+x^2)u}& \frac{4xy}{(1+y^2)u}\\ 0& 1\\ \frac{2x}{(1+x^2)u}& \frac{2y}{(1+y^2)u}\\ \frac{y(1x^2)}{(1+x^2)u}& \frac{x(1y^2)}{(1+y^2)u}\end{array}\right]$$ Let $`_x`$ and $`_y`$ be unit vector fields on $`^2`$ in the $`x`$ and $`y`$ directions, respectively. Applying $`dg`$ to these we get a map of $`^2`$ into $`V_{4,2}`$, which outside some large disk is constantly equal to $`(_1,_2)`$, where $`_i`$ is the unit vector in the $`z_i`$-direction in $`^4`$. Clearly, this map is homotopic in $`V_{4,2}`$ to a map of the form $`(v_1(x,y),_2)`$, where $`v_1(x,y)`$ is orthonormal to $`_2`$. Let $`D`$ be a large disk in $`^2`$. A straightforward calculation shows that the degree of $`v_1:(D^2,D^2)S^2`$ is one. Let $`B`$ be a 2-dimensional hemisphere in $`^4`$ which is flat close to its north pole. To make this precise, fix a small $`a>0`$ and let $`B`$ be a 2-disk such that $`B\{z^4:z_3\sqrt{116a^2}\}=`$ $`=\{(\xi ,\eta ,\zeta ,0):\xi ^2+\eta ^2+\zeta ^2=1,0\zeta \sqrt{116a^2}\},`$ $`B\{z^4:z_3\sqrt{19a^2}\}=`$ $`=\{(\xi ,\eta ,\sqrt{19a^2},0):\xi ^2+\eta ^24a^2\}.`$ Let $`(\xi ,\eta )`$, $`\xi ^2+\eta ^21`$, as above be coordinates on $`B`$. Then we can define an immersion $`K:B^4`$ with the following properties: The map $`K`$ equals the inclusion for $`\xi ^2+\eta ^2a^2`$. On the remaining part $`D_a`$ of $`B`$ the map $`K`$ is a suitably scaled Whitney kink and $`L(K(R(\xi ,\eta )))=K(\xi ,\eta )`$. Let $`(r,\vartheta )`$, $`0r\pi `$ and $`0\vartheta 2\pi `$ be polar coordinates on the disk $`D_a`$, then the map $`D_aV_{4,2}`$ induced by the differential of $`K`$ is homotopic to $`dK_{(r,\vartheta )}(_\xi )`$ $`\mathrm{cos}(r)_1+\mathrm{sin}(r)\left(\mathrm{cos}(\vartheta )_3\mathrm{sin}(\vartheta )_4\right),`$ $`dK_{(r,\vartheta )}(_\eta )`$ $`_2,`$ ($``$) since the first of these equations defines a map $`(D_a,D_a)S^2`$ of degree one. ### 8.2. A modified standard embedding. Let $`(y_1,\mathrm{},y_5)`$ be coordinates on $`^5`$. Let $`_i`$ be the unit vector in the $`y_i`$-direction. Let $`0\theta 2\pi `$ and let $`_+^4(\theta )^5`$ be the subset $$\{(x_0,x_1,x_2\mathrm{cos}\theta ,x_2\mathrm{sin}\theta ,x_3):x_20\}.$$ Let $`D(\theta )_+^4(\theta )`$ be the disk $$\{(x_0,x_1,x_2\mathrm{cos}\theta ,x_2\mathrm{sin}\theta ,0):x_20,x_0^2+x_1^2+x_2^2=1\}.$$ As $`\theta `$ varies from $`0`$ to $`2\pi `$, $`D(\theta )`$ sweeps the standard $`S^3`$ in $`^5`$. We note that if $`\theta _1\theta _2`$, $`\theta _10`$ then $$D(\theta _1)D(\theta _2)=\{(x_0,x_1,0,0,0):x_0^2+x_1^2=1\},$$ and $`D(0)=D(2\pi )`$. A straightforward calculation shows that the column vectors $`S^i(\theta )`$ of the matrix $`S(0,0,\mathrm{cos}\theta ,\mathrm{sin}\theta ,0,0)=\mathrm{\Theta }^s(0,0,\mathrm{cos}\theta ,\mathrm{sin}\theta ,0,0)`$ (see Lemma 3.3.2 for notation) are $`S^1(\theta )`$ $`=\mathrm{sin}(\theta )_3\mathrm{cos}(\theta )_4,`$ $`S^2(\theta )`$ $`=\mathrm{cos}(\theta )_1\mathrm{sin}(\theta )_2,`$ $`S^3(\theta )`$ $`=\mathrm{sin}(\theta )_1+\mathrm{cos}(\theta )_3,`$ $`S^4(\theta )`$ $`=_5,`$ $`S^5(\theta )`$ $`=\mathrm{cos}(\theta )_3+\mathrm{sin}(\theta )_4.`$ We now modify the standard embedding: Replace the disk $`D(\theta )`$ by $`B(\theta )`$, where $`B(\theta )`$ is obtained by rotating $`B_+^4(0)`$ to $`_+^4(\theta )`$. The embedding of $`S^3`$ so obtained is clearly diffeotopic to the standard embedding. We may assume that the framing it induces on $`D_a(\theta )`$ depends only on $`\theta `$ and equals $`S(\theta )`$. ### 8.3. Constructing immersions with arbitrary Smale invariant Let the linear map $`L_\alpha :^5^5`$ be given by the matrix $$\left[\begin{array}{cc}\begin{array}{cc}\mathrm{cos}\alpha & \mathrm{sin}\alpha \\ \mathrm{sin}\alpha & \mathrm{cos}\alpha \end{array}& \\ & 𝐈_3\end{array}\right].$$ Let $`(\xi ,\eta )`$ be coordinates on $`B(\theta )`$, as in Section 8.1. Let $`R_\alpha :B(\theta )B(\theta )`$ be the map that is given by the matrix $$\left[\begin{array}{cc}\mathrm{cos}\alpha & \mathrm{sin}\alpha \\ \mathrm{sin}\alpha & \mathrm{cos}\alpha \end{array}\right],$$ in the coordinates $`(\xi ,\eta )`$. Consider $`S^3=_{0\theta 2\pi }B(\theta )`$. For every half integer $`m\frac{1}{2}`$ we define $`f_m:S^3^5`$, $$f_m\|B(\theta )=L_{m\theta }KR_{m\theta },$$ where $`K`$ is as in Section 8.1. Then $`f_m`$ is well defined since it agrees with the inclusion on $`B(\theta )`$ for each $`\theta `$ and since $`K`$ is symmetric, i.e. $`L_{2m\pi }KR_{2m\pi }=K`$. The self intersection of $`f_m`$ is one circle. It is traced out by the double point of $`K`$ under rotation. The preimage of this circle is a torus knot of type $`(2m,2)`$ if $`m`$ is not an integer and a two component link with unknotted components that are linked with linking number $`m`$ if $`m`$ is an integer. ###### Proposition 8.3.1. The Smale invariant $`\mathrm{\Omega }(f_m)`$ satisfies, $$\mathrm{\Omega }(f_m)=2m[\sigma ]+N\pi _3(V_{5,3}).$$ (See Lemma 3.1.2 for notation.) Thus, as $`m`$ runs through $`\frac{1}{2}`$, $`f_m`$ runs through the regular homotopy classes of immersions $`S^3^5`$. The proof of Proposition 8.3.1 constitutes Section 9.3. ###### Corollary 8.3.2. The homomorphism (see Proposition 7.5.2) $$\tau :\mathrm{𝐈𝐦𝐦}_4,$$ is surjective. ###### Proof. The immersion $`f_{\frac{1}{2}}`$ represents a generator of $`\mathrm{𝐈𝐦𝐦}`$. It has one self intersection circle with connected preimage. Thus, $`\tau (f_{\frac{1}{2}})`$ is a generator of $`_4`$. ∎ ###### Lemma 8.3.3. The invariant $`\mathrm{lk}`$ (see Lemma 6.2.1) takes the values $$\mathrm{lk}(f_m)=4m.$$ We prove Lemma 8.3.3 in Section 9.4. ###### Corollary 8.3.4. The homomorphism (see Corollary 6.2.2) $$\lambda :\mathrm{𝐈𝐦𝐦}_3,$$ is surjective. ###### Proof. By Lemma 8.3.3, $`\lambda (f_{\frac{1}{2}})`$ is a generator of $`_3`$. ∎ ### 8.4. An immersion into 4-space As in Lashof and Smale , Theorem 3.4, we consider an immersion $`k:S^2^4`$ with one transversal double point. The normal bundle of this immersion has Euler number 2 and the boundary of a small tubular neighborhood of it is therefore an immersion $`h:US^2P^3^4`$, where $`US^2`$ denotes the unit tangent bundle of $`S^2`$. Let $`\pi :S^3P^3`$ be the universal double covering. Let $`f=h\pi :S^3^4`$. Then $`f`$ is an immersion. Let $`g:S^3^4`$ be a generic immersion regularly homotopic to $`f`$. ###### Proposition 8.4.1. The invariant $`\beta (g)`$ (see Proposition 7.4.1) is a generator of $`_8`$, the immersion $`g`$ has an odd number of quadruple points and the Euler characteristic of $`F_g`$ is odd. Proposition 8.4.1 is proved in Section 9.5 ###### Corollary 8.4.2. The $`_2`$-valued regular homotopy invariants $`Q_2`$ and $`D_2`$ defined in Section 6.1 are non-trivial.∎ ###### Corollary 8.4.3. The homomorphism (see Proposition 7.4.2) $$\beta :\mathrm{𝐈𝐦𝐦}_8$$ is an epimorphism.∎ ## 9. Proofs In this section we prove Theorems 1, 34, and 5, Proposition 8.3.1, Lemma 8.3.3, and Proposition 8.4.1 ### 9.1. Proof of Theorem 1 Proposition 4.1.1 shows that $`\sigma `$ is an isomorphism. Proposition 4.1.2 shows that $`\mathrm{𝐄𝐦𝐛}`$ is a subgroup of the infinite cyclic group $`\mathrm{𝐈𝐦𝐦}`$ of index 24. The Smale invariant $`\mathrm{\Omega }`$ takes values in $`\pi _3(V_{5,3})`$. Choosing a generator, we identify this group with $``$. For one of the two possible choices, the first part of the diagram is correct. (Making the other choice, the second homomorphism in the second row would be $`\times (24)`$ instead of $`\times 24`$). To prove the theorem it is enough to show that $`\tau \beta `$ is surjective. Lemma 8.3.3 shows that $`\tau (f_{\frac{1}{2}})`$ is a generator of $`_3`$ and if $`g:S^3^4^5`$ is as in Proposition 8.4.1 then $`\beta (g)`$ is a generator of $`_8`$. Hence, if either $`\beta (f_{\frac{1}{2}})`$ is a generator of $`_8`$ or $`\tau (g)`$ is a generator of $`_3`$ then $`\tau \beta `$ is surjective. Assume that this is not the case. Then $`\beta (f_{\frac{1}{2}})`$ is even and $`\tau (g)=0`$ and hence $$(\tau (f_{\frac{1}{2}}g),\beta (f_{\frac{1}{2}}g))=(\tau (f_{\frac{1}{2}}),\beta (f_{\frac{1}{2}}))+(\tau (g),\beta (g)),$$ is a generator of $`_3_8`$. Thus, $`\tau \beta `$ is surjective. (Since $`f_{\frac{1}{2}}`$ represents a generator of $`\mathrm{𝐈𝐦𝐦}`$, it actually follows that $`\beta (f_{\frac{1}{2}})`$ is odd.) ∎ ### 9.2. Proofs of Theorems 34, and 5 Consider the invariants $`\beta `$, $`\tau `$, $`D_2`$ and $`Q_2`$, defined in Section 6.1. Since $`D_2`$ and $`Q_2`$ are both additive under connected summation (see the proof of Lemma 7.4.2), it follows exactly as for $`\beta `$ that they are invariant under regular homotopy in $`^5`$. They both define homomorphisms $`\mathrm{𝐈𝐦𝐦}_2`$ and Lemma 8.4.1 implies that they are both nontrivial and therefore equal. Proof of Theorem 3: Let $`f:S^3^4`$ be a generic immersion. Consider $`f`$ as an immersion into $`^5`$ and evaluate the $`_2`$-valued invariants induced by $`D_2`$ and $`Q_2`$. By the above they are equal. The theorem follows.∎ Let $`\mu :\mathrm{𝐈𝐦𝐦}_2`$ denote the homomorphism obtained from $`Q_2`$ and $`D_2`$. We have the following sequence: $$\mathrm{𝐈𝐦𝐦}\stackrel{𝛽}{}_8\stackrel{r_4}{}_4\stackrel{r_2}{}_2,$$ where $`r_2`$ and $`r_4`$ denote reduction modulo 2 and reduction modulo 4, respectively. The Brown invariant of a pin structure on a surface reduced modulo 2 equals the Euler characteristic of the surface reduced modulo 2. Therefore, $`r_2r_4\beta =\mu `$. Moreover, if $`f:S^3^5`$ represents a generator of $`\mathrm{𝐈𝐦𝐦}`$ then $`r_4(\beta (f))`$ is a generator of $`_4`$ and so is $`\tau (f)`$. As a consequence, $$\beta (h)\mathrm{mod}4=\pm \tau (h)_4$$ for any immersion $`h:S^3^5`$. The sign in this formula depends on the choices of the isomorphisms $`\varphi `$ (Section 7.4) and $`\psi `$ (Section 7.5). For any one of the two possible choices of $`\varphi `$ there is a unique choice (out of two possible) of $`\psi `$ such that the sign in the above formula is $`+`$. Proof of Theorem 4: Let $`f:S^3^5`$ be a generic immersion and $`g:S^3^4`$ a generic immersion regularly homotopic to $`f`$. Then $`\tau (f)`$ is odd if and only if $`r_4(\beta (g))`$ is odd. But $`r_4(\beta (g))`$ is odd if and only if $`\beta (g)`$ is odd which is equivalent to $`F_g`$ having odd Euler characteristic. Apply Theorem 3.∎ Proof of Theorem 5: Let $`f`$ and $`g`$ be as above. If $`F_g`$ is orientable then $`\beta (g)`$ is divisible by $`4`$. Hence, $`r_4(\beta )=0`$ and therefore $`\tau (f)=0`$.∎ ### 9.3. Proof of Proposition 8.3.1 For simplicity of notation we let $`s:S^3^5`$ denote the modified standard embedding and $`S:S^3SO(5)`$ denote the map $`\mathrm{\Theta }^s`$ associated to it. Fix $`m`$ and let $`F:S^3SO(5)`$ denote the map $`\mathrm{\Theta }^{f_m}`$ associated to the immersion $`f_m:S^3^5`$. Let $`\mathrm{\Lambda }_m(x)=S^1(x)F(x)`$, where $`S^1(x)`$ is the inverse of the matrix $`S(x)`$ and $`S^1(x)F(x)`$ is the matrix product. To compute the Smale invariant of $`f_m`$ we must compute the homotopy class (see Lemma 3.3.2) $$[\mathrm{\Theta }^{f_m}][\mathrm{\Theta }^s]=[F][S]=[\mathrm{\Lambda }_m]\pi _3(SO(5))=\pi _3(V_{5,3}).$$ Since $`f_m=s`$ outside the solid torus $`T=_{0\theta 2\pi }D_a(\theta )`$, we have $$\mathrm{\Lambda }_m(x)=S^1(x)F(x)=\left[\begin{array}{cc}𝐈_3& \\ & 𝐉(x)\end{array}\right],\text{ for }xS^3T,$$ where $`𝐉(x)SO(2)`$. Let $`(\theta ,r,\phi )`$, $`\theta ,\phi [0,2\pi ]`$, $`r[0,\pi ]`$ be coordinates on $`T`$. Here $`\theta `$ indicates which $`D_a(\theta )`$ we are in and on $`D_a(\theta )`$ we have $$\xi =\frac{ra}{\pi }\mathrm{cos}(\phi ),\eta =\frac{ra}{\pi }\mathrm{sin}(\phi ).$$ Using equations ($``$) in Section 8.1 we can compute the first columns of $`F\|T`$. The result is the following: $`F^1(\theta ,r,\phi )=`$ $`\mathrm{sin}(\theta )_3\mathrm{cos}(\theta )_4,`$ $`F^2(\theta ,r,\phi )=`$ $`\mathrm{cos}(\theta m\theta ))v_1(\theta ,r,\phi )\mathrm{sin}(\theta m\theta )v_2(\theta ,r,\phi )`$ $`F^3(\theta ,r,\phi )=`$ $`\mathrm{sin}(\theta m\theta )v_1(\theta ,r,\phi )+\mathrm{cos}(\theta m\theta )v_2(\theta ,r,\phi ),`$ where $`v_1(\theta ,r,\varphi )=`$ $`\mathrm{cos}(r)\left[\mathrm{cos}(m\theta )_1+\mathrm{sin}(m\theta )_2\right]`$ $`+\mathrm{sin}(r)\left[\mathrm{cos}(\phi m\theta )\left(\mathrm{cos}(\theta )_3+\mathrm{sin}(\theta )_4\right)\mathrm{sin}(\phi m\theta )_5\right],`$ $`v_2(\theta ,r,\phi )=`$ $`\mathrm{sin}(m\theta )_1+\mathrm{cos}(m\theta )_2.`$ We note that $`F^1(x)=S^1(x)`$, for all $`xS^3`$. Hence, $$\mathrm{\Lambda }_m(x)=\left[\begin{array}{cc}1& \\ & 𝐌(x)\end{array}\right],\text{for all }xS^3,$$ where $`𝐌(x)SO(4)`$. Thus, to compute $`[\mathrm{\Lambda }_m]\pi _3(SO(5))`$ it is enough to compute the homotopy class of $`𝐌:S^3SO(4)`$ which, as we shall see, can be calculated if we know the first two columns of $`𝐌(x)`$. If $`𝐌=(m_i^j)`$ then $`m_j^i=S^{i1},F^{j1}`$, $`i,j=2,3,4,5`$, where $`,`$ is the usual inner product on $`^5`$. Hence, we have all the information we need to determine $`[\mathrm{\Lambda }_m]`$. In what follows, we first study how to calculate the homotopy class of $`𝐌`$ and then carry out the necessary calculations. In terms of our standard generators of $`\pi _3(SO(4))`$ we have, for $`𝐌:S^3SO(4)`$, $`[𝐌]=u[\sigma ]+v[\rho ]\pi _3(SO(4))`$, $`u,v`$. The integer $`u`$ is simply the degree of the map $`p𝐌=𝐌^1:S^3S^3`$, where $`SO(4)\stackrel{𝑝}{}S^3`$ is the fibration described in Section 3.1 and $`𝐌^1(x)`$ is the first column vector of $`𝐌(x)`$. To compute $`v[\rho ]`$, recall from Section 3.1 that $`SO(4)S^3\times SO(3)`$ and thus, we must compare the homotopy class of the map $`p_2𝐌:S^3SO(3)`$, where $`p_2:SO(4)SO(3)`$ is the projection onto the second factor in the product space $`SO(4)S^3\times SO(3)`$, to that of $`p_2\rho :S^3SO(3)`$. The product structure on $`SO(4)`$ is obtained by using the section $`\sigma :S^3SO(4)`$. The projection $`p_2:SO(4)SO(3)`$ is thus determined by the equation $$\sigma (p(x))^1x=\left[\begin{array}{cc}1& \\ & p_2(x)\end{array}\right],xSO(4).$$ We note that $`p_2\rho =\varrho :S^3SO(3)`$. Let $`𝐍=p_2𝐌`$. To compare the homotopy classes of these maps we can use the fibration $`SO(3)\stackrel{p}{}S^2`$: Since $`p_{}:\pi _3(SO(3))\pi _3(S^2)`$ is an isomorphism we might as well compare the homotopy classes of the maps $`p\varrho `$ and $`p𝐍`$ both mapping $`S^3`$ to $`S^2`$. By the Pontryagin construction the homotopy class of any map $`S^3S^2`$ is determined by its Hopf invariant (the linking number of two regular fibers). The map $`\varrho :S^3SO(3)`$, has Hopf invariant one if endow $`S^3`$ with the orientation that is coherent with the quaternion framing. Thus, to determine $`[𝐍]\pi _3(SO(3))`$ we need only calculate the Hopf invariant of $`p_2𝐍=𝐍^1:S^3S^2`$, where $`𝐍^1(x)`$ is the first column vector in the matrix $`𝐍(x)`$. This, in turn, means that the integer $`v`$ can be computed as the Hopf invariant of the map $`𝐍^1:S^3S^2`$, where $`S^3`$ has the same orientation as above. We now turn to the actual calculations. We have $`𝐌^k=_{i=1}^4S^{i+1},F^{k+1}_i`$, $`k=1,2`$ where $`_i`$, $`i=1,\mathrm{},4`$ is the standard basis in $`^4`$ and $`,`$ is the standard inner product on $`^5`$. Computing the necessary inner products we find that if $`xTS^3`$ then $$𝐌^1(x)=𝐌^1(\theta ,r,\phi )=\left[\begin{array}{c}\mathrm{cos}^2(m\theta \theta )\mathrm{cos}(r)+\mathrm{sin}^2(m\theta \theta )\\ \frac{1}{2}\left(\mathrm{cos}(r)+1\right)\mathrm{sin}(2m\theta 2\theta )\\ \mathrm{cos}(m\theta \theta )\mathrm{sin}(r)\mathrm{sin}(\phi m\theta )\\ \mathrm{cos}(m\theta \theta )\mathrm{sin}(r)\mathrm{cos}(\phi m\theta )\end{array}\right],$$ $$𝐌^2(x)=𝐌^2(\theta ,r,\phi )=\left[\begin{array}{c}\frac{1}{2}\left(\mathrm{cos}(r)+1\right)\mathrm{sin}(2m\theta 2\theta )\\ \mathrm{sin}^2(m\theta \theta )\mathrm{cos}(r)+\mathrm{cos}^2(m\theta \theta )\\ \mathrm{sin}(m\theta \theta )\mathrm{sin}(r)\mathrm{sin}(\phi m\theta )\\ \mathrm{sin}(m\theta \theta )\mathrm{sin}(r)\mathrm{cos}(\phi m\theta )\end{array}\right]$$ and if $`xS^3T`$ then $$𝐌^1(x)=\left[\begin{array}{c}1\\ 0\\ 0\\ 0\end{array}\right],𝐌^2(x)=\left[\begin{array}{c}0\\ 1\\ 0\\ 0\end{array}\right].$$ To find the first column of $`𝐍`$ we must calculate $`\sigma (p(𝐌))^1𝐌^2`$. For $`y^4`$, $`\sigma (p(𝐌))y=𝐌^1y`$, (for notation see Section 3.1). For $`q`$, let $`\overline{q}`$ denote its conjugate then $`\sigma (p(𝐌))^1y=\overline{𝐌^1}y`$. Multiplying gives $$\overline{𝐌^1(x)}𝐌^2(x)=\left[\begin{array}{c}0\\ \mathrm{cos}(r)\\ \mathrm{sin}(r)\mathrm{cos}(\phi \theta )\\ \mathrm{sin}(r)\mathrm{sin}(\phi \theta )\end{array}\right]\text{ if }xT\text{.}$$ Hence, the the first column of $`𝐍`$ is given by $$𝐍^1(x)=\left[\begin{array}{c}\mathrm{cos}(r)\\ \mathrm{sin}(r)\mathrm{cos}(\phi \theta )\\ \mathrm{sin}(r)\mathrm{sin}(\phi \theta )\end{array}\right]\text{ if }xT\text{ and }𝐍^1(x)=\left[\begin{array}{c}1\\ 0\\ 0\end{array}\right]\text{ if }xS^3T\text{.}$$ To determine the degree of $`𝐌^1`$ we consider preimages of the regular value $`(0,0,1,0)S^3`$. For $`xS^3T`$, $`𝐌^1(x)(0,0,1,0)`$ and for $`xT`$ we get the equations $`r`$ $`={\displaystyle \frac{\pi }{2}},`$ $`\theta `$ $`={\displaystyle \frac{l\pi }{m1}},`$ $`\phi `$ $`={\displaystyle \frac{ml\pi }{m1}}+(1)^l{\displaystyle \frac{\pi }{2}},`$ where $`l=0,1,\mathrm{},\|2m3\|`$. Thus, there are $`\|2m2\|`$ points in the preimage of $`(0,0,1,0)`$. The orientation of $`T`$ induced from the quaternion framing on $`S^3`$ is the same as that given by the basis $`(_\theta ,_r,_\phi )`$. Thus, the orientation of $`S^3`$ at $`(0,0,1,0)`$ induced by $`𝐌^1`$ is $`o((m1)_2,_1,_4)`$. The orientation of $`S^3`$ at $`(0,0,1,0)`$ induced from the standard orientation on $`^4`$ is $`o(_1,_2,_4)`$. So the sign of each point in the preimage is the same. It is positive if $`m<1`$ and negative if $`m>1`$. Thus, $`\mathrm{deg}(𝐌^1)=u=(2m2)=2m+2`$. To compute the Hopf invariant of $`𝐍^1:S^3S^2`$ we consider the preimage of the regular value $`(0,0,1)`$ and see that it is the curve $`c`$ defined by the equations $$r=\frac{\pi }{2},\varphi \theta =\frac{\pi }{2}.$$ Moreover, the vector field $`\frac{}{r}`$ along $`c`$, is mapped by $`d𝐍`$ to the vector $`(1,0,0)`$. Shifting the curve $`c`$ along this vector field we get a curve $`c^{}`$ and $`\mathrm{lk}(c,c^{})=1`$ using the above orientation on $`T`$. So, the Hopf invariant of $`𝐍^1`$ is $`1`$ and thus $`v=1`$. Collecting these results, we finally get $`\mathrm{\Omega }(f_m)`$ $`=[\mathrm{\Lambda }_m]=[𝐌]+N=((2m+2)[\sigma ][\rho ])+N`$ $`=2m[\sigma ]+N\pi _3(SO(5))=\pi _3(V_{5,3}),`$ as claimed.∎ ### 9.4. Proof of Lemma 8.3.3 Let $`(r,\theta ,\phi )`$ be coordinates on the solid torus where the Whitney kinks are placed. The preimage of the self intersection $`c_m`$ of $`f_m`$ is $`r=ϵ`$, $`\phi =m\theta `$ or $`\phi =m\theta +\pi `$, where $`ϵ`$ is some small positive number. Consider the vector field $`_r`$ along this preimage. If we shift $`c_m`$ along $$v(p)=(df_m(_r(p_1))+df_m(_r(p_2)),$$ where $`f_m^1(p)=\{p_1,p_2\}`$ then we obtain a curve $`c_m^{\prime \prime }`$, which bounds a 2-disk in the $`x_3x_4`$-plane without intersections with $`f_m(S^3)`$. Thus, $`\mathrm{lk}(c_m^{\prime \prime },f_m(S^3))=0`$. Assume that $`m`$ is not an integer. Then $`c_S=f_m^1(c_m)`$ is connected. Shifting $`c_S`$ along $`_r`$ gives a curve $`c_S^{\prime \prime }`$ with $`\mathrm{lk}(c_S,c_S^{\prime \prime })=4m`$. Thus, the framing of $`c_S`$ which satisfies condition $`()`$ in Section 6.2 is the one that differs from $`_r`$ by $`4m`$ rotations. Assume that $`m`$ is an integer. Then $`c_S=c_1c_2`$ have two components and $`\mathrm{lk}(c_1,c_2)=m`$. Shifting $`c_i`$ along $`_r`$ gives a curve $`c_i^{\prime \prime }`$ with $`\mathrm{lk}(c_i,c_i^{\prime \prime })=m`$. Thus, the framings of the $`c_i`$ that satisfy condition $`()`$ in Section 6.2 differ from $`_r`$ by $`2m`$ rotations each. The rotations that we need to add to $`_r`$ can be made locally supported. Using a local model it is easy to see how they affect the linking number: Let $`X`$ and $`Y`$ be the 3-planes $`X=(x_1,x_2,x_3,0,0)`$ and $`Y=(y_1,0,0,y_2,y_3)`$. Then $`XY`$ is the line $`l=(t,0,0,0,0)`$. Shifting $`l`$ along $`ϵ(_2+_4)`$ we get $`l^{}=(t,ϵ,0,ϵ,0)`$. Shifting $`l`$ along $`ϵ(\mathrm{cos}(t)_2+\mathrm{sin}(t)_3+_4)`$, $`\pi t\pi `$ we get the curve $`l^{\prime \prime }=(t,ϵ\mathrm{cos}(t),ϵ\mathrm{sin}(t),0,ϵ,0)`$. Let $`D=(1s)l^{}+sl^{\prime \prime }`$, $`0s1`$. Then $`D`$ is a 2-chain with boundary $`l^{\prime \prime }l^{}`$, $`DX=\mathrm{}`$ and $`DY=(0,0,0,ϵ,0)`$ with sign $`1`$. Thus, for each rotation in the negative direction we decrease the linking number by 1. Adding the appropriate number of rotations to the framing $`_r`$ and shifting the curve $`c_m`$ along the corresponding vector field we obtain a curve $`c_m^{}`$ with $$\mathrm{lk}(c_m^{},f_m(S^3))=\mathrm{lk}(c_m^{\prime \prime },f_m(S^3))4m=4m.$$ Thus, $`\mathrm{lk}(f_m)=4m`$.∎ ### 9.5. Proof of Proposition 8.4.1 The self intersection of $`h`$ is located close to the transverse double point $`h(p)=h(q)=0`$ of $`S^2`$. Let $`x`$ and $`y`$ be coordinates on small disks $`D_1`$ and $`D_2`$ around $`p`$ and $`q`$ respectively. We may assume that $$k(x)=(x_1,x_2,0,0)\text{and}k(y)=(0,0,y_1,y_2).$$ Close to $`p`$ and $`q`$, $`US^2`$ is of the form $`D_i^2\times S^1`$, $`i=1,2`$ and the immersion $`h`$ is: $$h(x,\theta )=(x_1,x_2,ϵ\mathrm{cos}(\theta ),ϵ\mathrm{sin}(\theta )),h(y,\phi )=(ϵ\mathrm{cos}(\phi ),ϵ\mathrm{sin}(\phi ),y_1,y_2),$$ where $`0\theta 2\pi `$ and $`0\phi 2\pi `$ are coordinates on $`S^1`$ in the fiber close to $`p`$ and $`q`$, respectively. Thus, $`M_h=T^2`$ and the preimage $`\stackrel{~}{M}_h`$ is $`c_i\times S^1D_i\times S^1`$, where $`c_i`$ is a circle of radius $`ϵ`$ around $`0D_i`$. We can find $`P^2US^2`$ as follows: Fix a point $`s`$ in $`S^2`$ and a unit tangent vector $`v`$ at that point. The space of 2-planes $`\mathrm{\Pi }`$ through $`0`$ in $`^3`$ is $`P^2`$. Let $`r_\mathrm{\Pi }`$ be the reflection in $`\mathrm{\Pi }`$. The map $`\mathrm{\Pi }(r_\mathrm{\Pi }(s),r_\mathrm{\Pi }(v))`$ is an embedding of $`P^2`$ into $`US^2`$. Consider the double cover $`p:S^3US^2`$. Observe that $`K=p^1(P^2)S^2`$ subdivides $`S^3`$ into two hemispheres $`D_+^3`$ and $`D_{}^3`$. The intersection $`\stackrel{~}{M}_hP^2`$ consists of two circles, isotopic to $`c_i\times 1D_i\times S^1`$, $`i=1,2`$. Altering the embedding of $`P^2`$ slightly we may assume that $`\stackrel{~}{M}_hP^2=c_1\times 1c_2\times 1`$. The covering $`p:S^3US^2`$ is nontrivial on fibers so $`p^1(\stackrel{~}{M}_h)`$ consists of two tori $`T_iS^3`$, $`i=1,2`$ and $`p\|T_i:T_ic_i\times S^1`$ is a double cover. Thus, $`KT_i`$ consists of two meridians in $`T_i`$, $`i=1,2`$. Choose a normal vector field $`\nu `$ to $`h`$. Let $`\delta >0`$ be small and let $`\varphi :S^3(\delta ,\delta )`$ be positive on $`D_+`$, negative on $`D_{}`$, $`0`$ on $`K`$, and have nonvanishing derivative in the normal direction of $`K`$. (Map $`(S^3,K)`$ to $`(S^3,\{x_4=0\})^4`$ and use a suitably scaled height function.) We use $`\varphi `$ and $`\nu `$ to perturb $`f`$. For $`xS^3`$, let $$g(x)=f(x)+\varphi (x)\nu (x).$$ Then $`g`$ is an immersion. Clearly, $`g`$ is an embedding outside a small neighborhood of $`T_1T_2K`$ in $`S^3`$. Moreover, inside this neighborhood it has only transverse double points if we stay away from $`T_1T_2`$. To determine the self intersection of $`g`$ we can therefore restrict attention to neighborhoods of $`T_1`$ and $`T_2`$. Let $`E_i\times S^1=p^1(D_i\times S^1)`$. If $`(x_1,x_2,e^{i\alpha })`$, $`(y_1,y_2,e^{i\omega })`$ are coordinates on $`E_1\times S^1`$ and $`E_2\times S^1`$, respectively then $`p:E_i\times S^1D_i\times S^1`$ is given by $`p(x_1,x_2,e^{i\alpha })`$ $`=(x_1,x_2,e^{i2\alpha }),`$ $`p(y_1,y_2,e^{i\omega })`$ $`=(y_1,y_2,e^{i2\omega }).`$ Normal fields along $`h(D_1\times S^1)`$ and $`h(D_2\times S^1)`$ are $`\nu _1(x_1,x_2,e^{i\theta })`$ $`=\mathrm{cos}(\theta )_3+\mathrm{sin}(\theta )_4,`$ $`\nu _1(y_1,y_2,e^{i\phi })`$ $`=\mathrm{cos}(\theta )_1+\mathrm{sin}(\theta )_2.`$ We may assume that the function $`\varphi \|E_i\times S^1`$ is given by $$\varphi (x_1,x_2,e^{i\alpha })=\delta \mathrm{sin}(\alpha )\text{and}\varphi (y_1,y_2,e^{i\omega })=\delta \mathrm{sin}(\omega ).$$ Then, in $`E_i\times S^1`$, $`g=f+\varphi \nu `$ is given by $`g(x_1,x_2,e^{i\alpha })`$ $`=(x_1,x_2,\left(ϵ+\delta \mathrm{sin}(\alpha )\right)\mathrm{cos}(2\alpha ),\left(ϵ+\delta \mathrm{sin}(\alpha )\right)\mathrm{sin}(2\alpha )),`$ $`g(y_1,y_2,e^{i\omega })`$ $`=(\left(ϵ+\delta \mathrm{sin}(\omega )\right)\mathrm{cos}(2\omega ),\left(ϵ+\delta \mathrm{sin}(\omega )\right)\mathrm{sin}(2\omega ),y_1,y_2).`$ Thus, $`g`$ has one quadruple point at $`(ϵ,0,ϵ,0)`$, with preimages $$(x_1,x_2,e^{i\alpha })=(ϵ,0,\pm 1)=(y_1,y_2,e^{i\omega }),$$ and $`g`$ has four triple lines: $`\alpha `$ $`(ϵ,0,\left(ϵ+\delta \mathrm{sin}(\alpha )\right)\mathrm{cos}(2\alpha ),\left(ϵ+\delta \mathrm{sin}(\alpha )\right)\mathrm{sin}(2\alpha )),`$ $`\omega `$ $`(\left(ϵ+\delta \mathrm{sin}(\omega )\right)\mathrm{cos}(2\omega ),\left(ϵ+\delta \mathrm{sin}(\omega )\right)\mathrm{sin}(2\omega ),ϵ,0),`$ where $`\alpha ,\omega (0,\pi )`$ or $`\alpha ,\omega (\pi ,2\pi )`$. Constructing $`F_g`$ by gluing the nine double point sheets of $`g`$ together gives $`F_g=T^2P^2`$ as shown in Figure 7, where the six intersection points of the solid lines represents $`F_g^0`$ and the remaining twelve open arcs of the solid lines represents $`F_g^1`$. It follows that $`\beta (g)`$ as well as $`\chi (F_g)`$ are odd. Since $`\beta `$ and $`D_2`$ are invariants of regular homotopy in $`^5`$, the proposition follows. ∎ ## Acknowledgments I want to thank Oleg Viro for valuable ideas and Ryszard Rubinsztein for fruitful discussions. I also want to thank the referee for pointing out that the problem on representing regular homotopy classes by embeddings was solved in . This shortened the paper considerably.
warning/0002/astro-ph0002313.html	ar5iv	text	# Resolving the extragalactic hard X-ray background NASA Goddard Space Flight Center, Code 662, Greenbelt, MD 20771 Institute for Astronomy, University of Hawaii, 2680 Woodlawn Drive, Honolulu, HI 96822 Astronomy Department, University of Maryland, College Park, MD 20742 To be published in Nature The origin of the hard ($`\mathrm{𝟐}\mathrm{𝟏𝟎}`$ keV) X-ray background has remained mysterious for over 35 years. Most of the soft ($`\mathbf{0.5}\mathrm{𝟐}`$ keV) X-ray background has been resolved into discrete sources, which are primarily quasars; however, these sources do not have the flat spectral shape required to match the X-ray background spectrum. Here we report the results of an X-ray survey 30 times more sensitive than previous studies in the hard band and four times more sensitive in the soft band. The sources detected in our survey account for at least 75 per cent of the hard X-ray background. The mean X-ray spectrum of these sources is in good agreement with that of the background. The X-ray emission from the majority of the detected sources is unambiguously associated with either the nuclei of otherwise normal bright galaxies or optically faint sources, which could either be active nuclei of dust enshrouded galaxies or the first quasars at very high redshifts. For some time after the discovery of the cosmic X-Ray background (XRB)giacconi , there was considerable controversy over whether the background arose from a superposition of discrete sources or from thermal bremsstrahlung emission from a hot intergalactic gas. We now know that the bulk of the XRB cannot originate in a uniform hot intergalactic medium since a strong Compton distortion on the cosmic microwave background spectrum was not observed by the FIRAS instrument on COBE mather2 ; wright . At soft X-ray energies ($`0.52`$ keV) the XRB has been extensively studied with the ROSAT satellite. The deepest ROSAT source counts reach $`1000`$ per square degree at a limiting flux of $`10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, and at this level $`7080`$ per cent of the XRB is resolved into discrete sources hasinger . The great majority of the optical identifications of a complete sample of 50 ROSAT sources, at a limiting flux of $`5\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup>, are unobscured active galactic nuclei (AGN) schmidt1 . However, because the objects detected in the soft band do not have the spectrum of the XRB, a new population of absorbed or flat spectrum objects are needed to make up the background at higher energies. Detailed models developed to resolve this “spectral paradox” assumed that most of the flux in the XRB is produced by active galaxies that are obscured by dust. When deep imaging sky surveys with the ASCAueda ; ueda1 ; ueda2 ; cagnoni and BeppoSAXfiore satellites became possible in the hard ($`>2`$ keV) X-ray band, $`30`$ per cent of the hard XRB was resolved, but only indirect identifications of the optical counterparts could be made. The Chandra satelliteweisskopf , with its great sensitivity over a wide energy range, excellent image quality, superb positional accuracy, and reasonable field-of-view, can directly image the sources that make up the hard XRB. We have therefore carried out a deep imaging survey of the Hawaii Deep Survey Field SSA13 with the ACIS-S instrument on Chandra to resolve the hard XRB and to identify the nature of the sources that produce it. We chose to centre on the SSA13 fieldlilly , which has existing multiwavelength observationswindhorst ; songaila ; barger , to maximize the immediate identification of optical/near-infrared (NIR) counterparts and redshifts for the X-ray source detections. We find that above a flux threshold of $`2.5\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`210`$ keV), we can account for at least $`75`$ per cent of the sky flux, with the main uncertainty being the sky flux itself. Our deep optical observations show a rich assortment of hard X-ray sources which could not have been discovered by previous satellites. Chandra X-ray Survey of SSA13 The SSA13 observation was performed on 1999 December 3–4 for an elapsed time of 100.9 ks. The optical axis of the telescope at RA(2000)$`=13^h12^m21.40^s`$, Dec(2000)$`=42^{}41^{^{}}20.96^{^{\prime \prime }}`$ was positioned on the back illuminated CCD (S3) of ACIS since this detector has a much better soft X-ray sensitivity than the front illuminated chips. Furthermore, since the back illuminated detectors did not suffer the radiation damage which affected the front illuminated chips in orbit, they are well characterised by extensive ground-based calibrations. The overall sensitivity of the instrument spans a wide energy range from 0.2 to 10 keV. Two energy-dependent images of the S3 chip were generated in the hard ($`210`$ keV) and soft ($`0.52`$ keV) bands, as was a $`210`$ keV image of the front illuminated S2 chip that covered a neighboring region. We extracted sources independently for the hard and soft band images. Sources brighter than $`3.2\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`210`$ keV) or $`3\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`0.52`$ keV; S3 chip only) which lie within 6 arcminutes of the optical axis are given in Table 1, ordered by right ascension; the table contains 22 sources selected in the hard band and a further 15 sources selected solely in the soft band. Details of the extraction and calibration of the X-ray data and of the optical photometry may be found in the table footnote. Number Counts and the Resolution of the X-ray Background The cumulative counts per square degree, $`N(>S)`$, are the sum of the inverse areas of all sources brighter than flux $`S`$. Sources at the faintest fluxes can be detected only at smaller off-axis angles where the PSF and vignetting corrections are smaller; thus, the area diminishes with flux. In Fig. 1a, b we present our cumulative counts per square degree (filled squares) in the soft and hard bands, respectively, with $`1\sigma `$ uncertainties from the Poisson error in the number of detected sources (jagged solid lines). To the limiting flux levels of $`2.3\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`0.52`$ keV) and $`2.5\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`210`$ keV), simulations show that the counts are nearly complete and that Eddington bias is unimportant; thus, the raw counts accurately represent the true counts. Our soft band counts are in excellent agreement with the deep ROSAT counts in the Lockman Hole from Hasinger et al.hasinger in the region of overlap. At fainter fluxes our new counts fall at the lower limit of their fluctuation analysis, which suggests an ongoing flattening. An area-weighted maximum likelihood fitmurdoch of a single power-law to the $`0.52`$ keV counts over the flux range $`2.370\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> is given by the relation $$N(>S)=185\times (S/7\times 10^{15})^{0.7\pm 0.2}$$ (1) where the errors on the power-law index are 68% confidence. Likewise, a power-law fit to the $`210`$ keV counts over the flux range $`2.520\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> is given by $$N(>S)=170\times (S/2\times 10^{14})^{1.05\pm 0.35}$$ (2) where the counts intercept the ASCA extrapolation at the upper end of the flux range. Though the range in indices is consistent with the power-law index of 1.5 seen at brighter fluxes, the counts are significantly lower than an extrapolation of the ASCA counts. The source contributions to the XRB can be obtained by summing the individual fluxes divided by area or, more indirectly, by integrating $`SdN`$ using the power-law fits. We list the directly summed source contributions to the XRB in the two bands in Table 2, along with previous determinations by ROSAT and ASCA. With the additional 10 per cent contribution from our data to the soft band, a maximum flux of $`1.1\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> remains to be accounted for. In the hard band, the combination of the present results with the ASCA measurements at higher fluxes means that at least $`75`$ per cent of the background (using the highest published normalisation) is resolved to the currently observed flux limits. Optical Properties of the X-ray Sources We have compared our X-ray images with existingcowie96 ; wilson deep $`HK^{}`$, $`I`$, $`B`$, and $`U^{}`$ images obtained with the Keck 10 m and UH 2.2 m telescopes. Because of the excellent $`1`$ arcsec X-ray positional accuracy, we can, in most cases, securely identify the optical counterparts to the X-ray sources. In Fig. 2a, b we show thumbnail $`I`$-band images of all of the X-ray sources in Table 1. Only one source, CXO J131159.3+423928 (significant in both the hard and soft X-ray images), is significantly extended in the X-ray images; it is probably a high redshift cluster. The optical image (thumbnail 35 of Fig. 2a) is centred on a faint ($`I=23`$) galaxy which lies at the centre of a region of enhanced galaxy density. In addition to the probable cluster, the hard sample contains two quasars, eight bright galaxies, and eleven optically faint ($`I>23`$) objects, while the soft sample contains five quasars, five bright galaxies, and fifteen unidentified optically faint objects. Morphologically the X-ray selected bright galaxy population consists of a mixture of early spirals and elliptical galaxies. Three of the bright galaxies show possible signs of interaction with nearby bright neighbors while the remainder are clearly isolated. Our Keck spectra for the quasars show broad MgII or CIV and Lyman alpha lines. In some of the bright galaxy spectra clear AGN signatures are present (e.g., a broad absorption line galaxy with P-Cygni profiles at $`z=1.320`$). However, although subtle AGN signatures may be present in their optical spectra, most of the bright galaxies would not have been identified in an optical survey as AGN. We illustrate this in Fig. 3 with spectra for three of the bright galaxies. The X-ray Spectrum The photon intensity of the XRB, $`P(E)`$, where $`E`$ is the photon energy in keV and $`P(E)`$ has units of \[photons cm<sup>-2</sup> s<sup>-1</sup> keV<sup>-1</sup> sr<sup>-1</sup>\], can be approximated by a power-law, $`P(E)=AE^\mathrm{\Gamma }`$. The HEAO1 A-2 experimentmarshall found that the XRB spectrum from $`315`$ keV was well described by a photon index $`\mathrm{\Gamma }1.4`$, and this result has been confirmed and extended to lower energies by recent analyses of ASCAgendreau ; chen ; miyaji ; ishisaki and BeppoSAXvecchi data. The photon indices of the individual sources given in Table 1 were computed from the ratios of the counts in the $`0.52`$ keV band to those in the $`210`$ keV band, assuming each source could be described by a single power-law. There is an extremely wide range of hardness in both samples, ranging from negative indices to $`\mathrm{\Gamma }=1.8`$ in the hard-selected sample and from $`\mathrm{\Gamma }=0.1`$ to values above 2 in the soft sample. The composite (counts-weighted) photon index is $`1.22\pm 0.03`$ in the hard sample and $`1.42\pm 0.04`$ in the soft sample. The progressive hardening of the soft sample as we move to fainter fluxes is a continuation of a trend seen in the ROSAT sampleshasinger93 ; almaini . The combined spectrum of all the soft X-ray sources of Table 1 is well fit by a single power-law over the $`0.310`$ keV range with an index of $`1.42\pm 0.07`$ and an extinction corresponding to the galactic $`N(H)=1.4\times 10^{20}`$ cm<sup>-2</sup>. If we assume that 75 per cent of the $`210`$ keV background has an index of 1.22 and that the remaining 25 per cent of the background comes from sources that have fluxes greater than $`1\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`210`$ keV) and an average photon indexueda1 of 1.63, then the index of the combined sample is 1.38, which agrees extremely well with the spectrum of the hard X-ray background marshall ; gendreau ; chen ; miyaji ; ishisaki ; vecchi . Inspection of Table 1 suggests that the hardest sources tend to correspond to the bright galaxies, with the optically faint objects having intermediate hardness, and the quasars being the softest of the sources observed. We illustrate this more clearly in Fig. 4 where we have plotted $`I`$-band magnitude versus photon index. Of the $`I22`$ objects, more than half are galaxies, and the majority of these have $`\mathrm{\Gamma }<1`$. The five known quasars all lie in the $`\mathrm{\Gamma }>1.7`$ range, consistent with that of most brighter AGN. The faint sources spread over a wide range of indices that overlap both of the other populations. We can quantify this by generating the counts-weighted averages for each population separately. For the hard-selected sample, we find that the bright galaxies (#9, 12, 26, 29) have an average photon index of $`0.59\pm 0.06`$, the faint objects (#1, 6, 7, 8, 9, and 22) have $`1.33\pm 0.06`$, and the two quasars have $`1.76\pm 0.07`$. For the soft-selected sample, the 15 unidentified objects with $`I>23`$ have a composite index of $`1.35\pm 0.06`$, which is almost identical to that of the optically faint objects in the hard sample, and the quasars have a composite index of $`1.80\pm 0.12`$. The Source of the Background Our data conclusively show that AGN are the major contributors to the hard X-ray background. Many of our sources agree with the predictions of XRB synthesis modelssetti ; madau ; matt ; comastri ; zdziarski ; smith ; gilli ; schmidt2 ; miyaji2 constructed within the framework of AGN unification schemes to account for the spectral intensity of the hard XRB and to explain the X-ray source counts in the hard and soft energy bands. In the unified scheme, the orientation of a molecular torus surrounding the nucleus determines the classification of the source. The models invoke, along with a population of unobscured AGN, whose nuclear emission we see directly, a substantial population of intrinsically obscured AGN whose hydrogen column densities of $`N_H10^{21}10^{25}`$ cm<sup>-2</sup> around the nucleus block our line-of-sight. The AGN that make up the hard XRB come in two main flavors: roughly 40 per cent are luminous early-type galaxies (both ellipticals and early spirals) in the redshift range from $`z=0`$ to just beyond $`z=1`$, and roughly 50 per cent have faint or, in some cases, undetectable optical counterparts. Most of these objects would not have been found even in sensitive optical surveys for AGN. The bright galaxy population is extremely hard with an average photon index of $`\mathrm{\Gamma }=0.59`$. The X-ray sources are point-like and centred on the galaxy nuclei, which suggests that they are produced by accretion onto the central black holes that are known to be present in such systems. The hardness of the X-ray spectra indicates that these X-ray sources are highly obscured. Such sources were described by Moran et al.moran based on Einstein data. After hard X-ray components were discovered by Allen et al.allen in ASCA spectra of six nearby giant elliptical galaxies, a modeldm1 was constructed which was able to account for a large fraction of the XRB with objects of this type. The model predicted that a significant fraction of the hard number counts at fluxes $`<10^{14}`$ erg cm<sup>-2</sup> s<sup>-1</sup> could arise from sources at low redshift, as is indeed now observed to be the case. The absolute $`K`$ magnitudes of these sources lie between $`24`$ and $`26`$ ($`\mathrm{H}_\mathrm{o}=65`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_o=0.5`$), or from just below to several times the $`L_{}`$ luminosity, and their rest-frame $`210`$ keV luminosities range from $`5\times 10^{41}`$ to $`3\times 10^{43}`$ erg s<sup>-1</sup>. These sources are at too low redshifts to be likely submillimeter candidates; however, they should be far-infrared sources, which SIRTF and other upcoming airborne and space missions should be able to detect. The optically faint sources have an average photon index of $`\mathrm{\Gamma }=1.3`$. These sources could either be a smooth continuation to $`z>1`$ of the bright early-type galaxies with obscured luminous X-ray nuclei, more distant obscured AGN, or something more exotic, such as extremely high redshift ($`z5`$) quasars. For this final possibility, the objects would be invisible in the $`B`$-band because of scattering by the foreground intergalactic neutral hydrogen. In the soft sample, eleven of the optically faint sources have $`B>26`$ and could lie in this category. This places an upper limit on the surface density of this type of source of 0.26 per square arcminute to the $`3\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> limit of the $`0.52`$ keV sample, which is slightly lower than the predictions of the toy model of Haiman and Loebhaiman for X-ray selected high-redshift quasars. The handful of objects which are detected in the NIR but absent in $`B`$ are the most promising candidates for this type of object and, in some cases (e.g., object 19 in the soft sample) may be bright enough for follow-up with NIR spectroscopy to test the hypothesis. Comparison with submillimeterbarger1 ; hughes ; gunn , far-infrared, and radio samples should allow us to determine what fraction of objects in these surveys are X-ray emitting AGN. As near-infrared spectra and photometric redshift estimates of the optically faint sources are established, we will be able to refine the obscured AGN models and determine whether any of the faint sources are indeed very high redshift quasars. With the X-ray, optical, and submillimeter samples all now approaching the full resolution of their respective backgrounds, we are close to achieving a complete cosmic census of the population of galaxies and AGN. Acknowledgements We thank J. Halpern and G. Hasinger for comments which greatly improved the first draft of this paper. We acknowledge E. Boldt for his many years of pioneering work concerning the X-ray background and R. Giacconi, whose insight and enthusiasm have inspired this subject. We thank the CXC, L. Van Speybroeck, M. Weisskopf and the MSFC team, M. Bautz, G. Garmire, and the ACIS team for building and operating such an excellent observatory. We acknowledge the use of HEASARC software. A.J. Barger acknowledges support from the Hubble and Chandra fellowship programs. Notes to Table 1 X-ray sources in the SSA13 field selected in the hard ($`210`$ keV; S2 and S3 chips) or soft ($`0.52`$ keV; S3 chip) bands. The X-ray images were prepared using xselect and associated ftools at GSFC. ACIS grades 0, 2, 3, 4, and 6 were used, and columns at the boundaries of the readout nodes were rejected. Counts lying within a $`5^{\prime \prime }`$ diameter aperture were measured, together with the background in a $`5^{\prime \prime }7.5^{\prime \prime }`$ radius annulus, at $`2^{\prime \prime }`$ intervals along the field. The distribution of detected counts is Poisson. A cut of 17 counts in the hard S3 image and 10 counts in the other two images represents a $`<10^7`$ probability threshold against background fluctuations and ensures a $`<20`$% probability of a single spurious source detection in the entire sample. The source counts were corrected for the enclosed energy fraction within the aperture. For the S3 chip the flux calibrations were made using an array of effective areas versus energy at 12 positions and an assumed power-law spectrum having counts-weighted mean photon indices $`\mathrm{\Gamma }=1.2`$ ($`210`$ keV) and $`\mathrm{\Gamma }=1.4`$ ($`0.52`$ keV). The galactic $`N(H)=1.4\times 10^{20}`$ cm<sup>-2</sup> is too low to affect the flux conversions. For the S2 chip a single conversion factor of $`2.6\times 10^{11}`$ erg cm<sup>-2</sup> ct<sup>-1</sup> was used. Using the on-axis flux calibrations of $`2.5\times 10^{11}`$ ($`210`$ keV) and $`2.9\times 10^{12}`$ ($`0.52`$ keV) to convert the S3 counts per second to flux, we determine limiting minimum fluxes of $`3.2\times 10^{15}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`210`$ keV) and $`3.0\times 10^{16}`$ erg cm<sup>-2</sup> s<sup>-1</sup> ($`0.52`$ keV). The table is restricted to sources with off-axis angles $`<6^{}`$, where $`>50`$% of the energy is enclosed within a $`2.5^{\prime \prime }`$ radius, and to sources where the noise, computed from the variance of the background and signal, is less than one third the signal. The $`15^{\prime \prime }`$ borders of each chip have incomplete exposure times due to the spacecraft dither so objects detected in these borders were not included in our counts analysis. The NIR and optical magnitudes are computed in $`1.5^{\prime \prime }`$ radii intervals. $`I`$ is Kron-Cousins, $`B`$ is Johnson, $`HK^{}`$ is a broad filter centred at $`1.9`$ microns, and $`U^{}`$ is a $`300`$ Å filter centred at $`3400`$ Å. Lower limits are $`1\sigma `$. Notes to Table 2 The statistical errors on our observed sky brightnesses dominate the systematic errors, which are expected to be less than 10 per cent. To be consistent with Hasinger et al.hasinger , we converted our soft band sky brightness of $`6.0\pm 1.5\times 10^{13}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> to the $`12`$ keV range using the measured mean photon index. This result was then compared with the $`12`$ keV background (where galactic contamination is less than at lower energies) of $`3.74.4\times 10^{12}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> from Gendreau et al.gendreau , using a fit to ASCA data, and Chen et al.chen , using a fit to joint ASCA/ROSAT data. In the hard band the summed counts are compared with the $`210`$ keV background of $`1.62.3\times 10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup> deg<sup>-2</sup> from Marshall et al.marshall , using a fit to HEAO1 A2 data, and Vecchi et al.vecchi , using a fit to BeppoSAX data.
warning/0002/hep-ph0002001.html	ar5iv	text	# Compact Hyperbolic Extra Dimensions: Branes, Kaluza-Klein Modes and Cosmology \[ ## Abstract We reconsider theories with low gravitational (or string) scale $`M_{}`$ where Newton’s constant is generated via new large-volume spatial dimensions, while Standard Model states are localized to a 3-brane. Utilizing compact hyperbolic manifolds (CHM’s) we show that the spectrum of Kaluza-Klein (KK) modes is radically altered. This allows an early universe cosmology with normal evolution up to substantial temperatures , and completely negates the constraints on $`M_{}`$ arising from astrophysics. Furthermore, an exponential hierarchy between the usual Planck scale and the true fundamental scale of physics can emerge with only $`𝒪(1)`$ coefficients. The linear size of the internal space remains small. The proposal has striking testable signatures. PACS:12.10.-g, 11.10.Kk, hep-ph/0002001, CERN-TH-2000-038, 11.25.M,04.50.+h CWRU-P1-00, SU-ITP-00/05 \] Recent work has heralded a renewed interest in higher-dimensional space-times, a key new concept being the localization of matter, and even gravity, to branes embedded in the extra dimensions . In the canonical example of , space-time is a direct product of ordinary 4D space-time and a (flat) spatial $`d`$-torus of common linear size $`R`$ and volume $`V_{\mathrm{new}}=R^d`$, while Standard Model particles are localized on a 3-brane of thickness $`M_{}^1`$, where $`M_{}`$ is the new fundamental higher-dimensional gravitational (or string) scale. The low energy effective 4D Planck scale $`M_P`$ is then given by the Gauss’s Law relation, $`M_P^2=M_{}^{2+d}R^d`$. The hierarchy between $`M_P`$ and $`M_{}`$ can be very large if $`RM_{}1`$. For example, if $`d=2`$ and $`R\mathrm{mm}`$, then $`M_{}\mathrm{TeV}`$. The hierarchy $`M_P/\mathrm{TeV}`$ thus becomes a problem of understanding the size of the extra dimensions in such a model . Remarkably, models with $`R`$ approaching the sub-millimeter range are not excluded , but astrophysics and cosmology do place significant bounds. In particular, the evolution of the early universe at temperatures just above those at the epoch of Big Bang Nucleosynthesis (BBN) is inevitably, and dramatically altered. This narrow range of normal evolution prior to BBN makes it difficult to implement baryogenesis, moduli dilution etc. The most important model-independent constraints on such models arise from the production of light KK modes of the graviton. These KK modes are the eigenmodes of the appropriate Laplace operator $`\mathrm{\Delta }`$ on the internal space, and it is of central importance in the following that all the constraints depend on the form of the spectral density of this operator, which in turn depends completely on the topology and geometry of the internal space. In this letter we argue that attractive alternate choices of compactification imply significantly weaker constraints, admitting in particular a standard 4D Friedmann-Robertson-Walker (FRW) evolution up to high temperatures. These compactifications employ a topologically non-trivial internal space— a $`d`$-dimensional compact hyperbolic manifold (CHM). They also throw into a new light the problem of explaining the large hierarchy $`M_P/\mathrm{TeV}`$, since even though the volume of these manifolds is large, their linear size $`L`$ is only slightly larger than the new fundamental length scale ($`L30M_{}^1`$ for example), thus only requiring numbers of $`𝒪(10)`$. CHM’s are obtained from $`H^d`$, the universal covering space of hyperbolic manifolds (those admitting constant negative curvature), by modding out by an appropriate freely acting discrete subgroup $`\mathrm{\Gamma }`$ of the isometry group of $`H^d`$ . (If $`\mathrm{\Gamma }`$ is not freely-acting, then the resulting quotient is a non-flat non-smooth orbifold. We will not discuss this interesting case here.) These manifolds have been much discussed recently as the possible structure of ordinary 3-space , and play an important role in the theory of classical and quantum “chaotic” systems, where the spectra of Laplacian operators are also vital . Here we will consider space-times of the form $`M^4\times (H^d/\mathrm{\Gamma }\|_{\mathrm{free}})`$ ($`M^4`$ is a FRW 4-manifold) with metric $$G_{IJ}dz^Idz^J=g_{\mu \nu }^{(4)}(x)dx^\mu dx^\nu +R_c^2g_{ij}^{(d)}(y)dy^idy^j.$$ (1) Here $`R_c`$ is the physical curvature radius of the CHM, so that $`g_{ij}(y)`$ is the metric on the CHM normalized so that its Ricci scalar is $`=1`$, and $`\mu =0,\mathrm{},3`$, $`i=1,\mathrm{},d`$. Because they are locally negatively curved, CHM’s exist only for $`d2`$. Their properties are well understood only for $`d3`$; however, it is known that CHM’s in dimensions $`d3`$ possess the important property of rigidity . As a result, these manifolds have no massless shape moduli. Moreover, the volume of the manifold, in units of the curvature radius $`R_c`$, cannot be changed while maintaining the homogeneity of the geometry. Hence, the stabilization of such internal spaces reduces to the problem of stabilizing a single modulus, the curvature length or the “radion”. Of course, in a complete high-energy theory, (e.g. string theory), there will be massive $`𝒪(M_{})`$ excitations of the would-be shape moduli, and more important for the constraints, the massive KK modes. To uncover the physics of these models one must consider the spectrum of small fluctuations $`h`$ in the metric around the background eq. (1), $`G_{IJ}G_{IJ}+\mathrm{e}^{ip.x}h_{IJ}(y)`$. There are 3 different types of KK fluctuations that so arise: $`h_{\mu \nu }`$, the spin-2 piece; $`h_{ij}`$, with indices only in the internal directions, giving spin-0 fields for the 4D observer; and the mixed case $`h_{i\mu }`$, giving spin-1 4D fields. The 4D KK masses of these states are the eigenvalues of the appropriate internal-space Laplacians acting on $`h(y)`$, the correct Laplacian differing between these 3 cases. In the most important spin-2 case the operator is the Laplace-Beltrami operator $`\mathrm{\Delta }_{LB}`$ (the Laplacian on scalar functions in the internal space), defined by $$\mathrm{\Delta }_{LB}\varphi (y)=\|g(y)\|^{1/2}_i\left(\|g(y)\|^{1/2}g^{ij}_j\varphi (y)\right).$$ (2) There are no known analytic expressions for the individual eigenvalues of $`\mathrm{\Delta }_{LB}`$ on a CHM of any dimension. However, despite the extremely complicated topology and geometry of CHM’s with arbitrarily large volume, a number of simple facts are generally true. First, by a variational argument, the spectrum of $`\mathrm{\Delta }_{LB}`$ is bounded from below, and the lowest eigenmode is just the constant function on the CHM. This zero mode is the internal space wave-function of the massless spin-2 4D graviton. As it is a constant, the effective 4D Planck mass depends only on the volume of the (highly curved) internal space. For example, suppose that the internal space was a 3-sphere of radius $`r`$, cut out of an $`H^3`$ of curvature radius $`R_c`$. Its volume $`\mathrm{Vol}(r)`$ grows exponentially for $`rR_c`$, $$\mathrm{Vol}(r)=\pi R_c^3[\mathrm{sinh}(2r/R_c)2r/R_c].$$ (3) In general, the total volume of a smooth compact hyperbolic space in any number of dimensions is $$\mathrm{Vol}_{\mathrm{new}}=R_c^de^\alpha ,$$ (4) where $`\alpha `$ is a constant, determined by topology. (For $`d=3`$ it is known that there is a countable infinity of orientable CHM’s, with dimensionless volumes, $`e^\alpha `$, bounded from below, but unbounded from above. Moreover, the $`e^\alpha `$ do not become sparsely distributed with large volume.) In addition, because the topological invariant $`e^\alpha `$ characterizes the volume of the CHM, it is also a measure of the largest distance $`L`$ around the manifold. CHM’s are globally anisotropic; however, since the largest linear dimension gives the most significant contribution to the volume, one can employ eq. (3), or its generalizations to $`d3`$, to find an approximate relationship between $`L`$ and $`\mathrm{Vol}_{\mathrm{new}}`$. For $`LR_c/2`$ the appropriate asymptotic relation, dropping irrelevant angular factors, is $$e^\alpha \mathrm{exp}\left((d1)L/R_c\right).$$ (5) Thus, in strong contrast to the flat case, the expression for $`M_P`$ depends exponentially on the linear size, $$M_P^2=M_{}^{2+d}R_c^de^\alpha =M_{}^{2+d}R_c^d\mathrm{exp}\left((d1)L/R_c\right).$$ (6) The most interesting case (and as we will see later, most reasonable) is the smallest possible curvature radius, $`R_cM_{}^1`$. Taking $`M_{}`$ TeV then yields $$L35M_{}^1=10^{15}\mathrm{mm}.$$ (7) Therefore, one of the most attractive features of a CHM internal space is that to generate an exponential hierarchy between $`M_{}`$ TeV, and $`M_P`$ requires only that the linear size $`L`$ be very mildly tuned. We now return to the important topic of the non-zero eigenmodes of $`\mathrm{\Delta }_{LB}`$ on CHM’s, and to the astrophysical and cosmological implications of these KK modes. Recall that in flat models, the KK modes are extremely light, $`m_{KK}R^110^4\mathrm{eV}`$, and very numerous, $`N_{KK}M_P^2/M_{}^210^{32}`$ . As a result, even though these modes are individually only weakly coupled, with strength $`1/M_P`$, they can be copiously produced by energetic processes on our brane, and observational limits then constrain the fundamental scale. The tightest astrophysical constraint comes from supernova physics, leading to a lower bound of $`M_{}50\mathrm{T}\mathrm{e}\mathrm{V}`$ if $`d=2`$, and of $`M_{}3\mathrm{T}\mathrm{e}\mathrm{V}`$ for $`d=3`$ . There are also severe limits on the maximum temperature (the “normalcy temperature” $`T_{}`$) above which the evolution of the universe must be non-standard . This temperature is found by equating the rates for cooling by the usual process of adiabatic expansion, and by the new process of evaporation of KK gravitons into the bulk. This gives $`T_{}10`$ MeV for $`d=2`$, up to $`T_{}10`$ GeV when $`d=6`$. As we will now see, for us the situation is much improved. First, by the compactness of the internal space, the spectrum of $`\mathrm{\Delta }_{LB}`$ on a CHM is discrete and has a gap between the zero mode and the first excited KK state. The size of this gap is all important. Second, most of the eigenmodes of $`\mathrm{\Delta }_{LB}`$ on a CHM have wavelengths less than $`R_c`$, and the number density of these modes is well approximated by the usual Weyl asymptotic formula $$n(k)=(2\pi )^d\mathrm{\Omega }_{(d1)}V_dk^{d1},$$ (8) where $`\mathrm{\Omega }_{(d1)}=\mathrm{Area}(S^{d1})`$. There can also be a few lighter supercurvature modes, with wavelengths as large as the longest linear distance in the manifold, and masses thus bounded below by $`L^1`$. There is no simple expression for the spectral density of supercurvature modes, although the Selberg trace formula provides some information on the full spectrum of $`\mathrm{\Delta }_{LB}`$. Nevertheless bounds on the first non-zero eigenvalue are known. In the best-studied CHM case of $`d=2`$ we have the following theorem : Consider a compact (oriented) Riemann surface $`S_g`$ of arbitrary genus $`g2`$, with metric of constant negative curvature -1. Then for every $`ϵ`$, there exists $`NZ^+`$ such that for $`g>N`$ there exists an $`S_g`$ with first eigenvalue $$\lambda _1(S_g)(Cϵ),$$ (9) where $`C171/784`$ by earlier work . Restoring units, a large enough volume (and thus genus) $`d=2`$ CHM will have first eigenvalue $`171/(784R_c^2)`$. Moreover, Brooks has conjectured that for $`d=2`$ a typical CHM chosen at random will have first eigenvalue $`1/4R_c^2`$ with positive probability $`P`$, perhaps even with $`P1`$ as the genus $`g\mathrm{}`$ . The analogous conjecture in $`d=3`$ is more problematic, but has also been made . Numerical studies of the spectra of even small volume $`d=3`$ CHM’s show that they have very few modes with $`\lambda <R_c`$ . The crucial result is that the first KK modes are exponentially more massive than the very light $`m_{KK}1/V^{1/d}`$ found in the flat case. This drastically raises the threshold for their production. Even making the pessimistic assumption that the spectral density of the supercurvature modes satisfies eq. (8) for $`k>1/L`$, the astrophysical bounds of and completely disappear since the lightest KK mode has a mass (at least 30 GeV), much greater than the temperature of even the hottest astrophysical object. Similarly the large KK masses imply a much higher normalcy temperature $`T_{}`$, up to which the evolution of our brane-localized 4D universe can be normal radiation-dominated FRW. Approximate numerical evaluation shows that $`T_{}`$ is understandably sensitive to the gap to the first non-zero KK mass, ranging from 2 GeV to 10 GeV (for $`d=2`$ to $`d=6`$) if $`m_{KK,1}1/L\mathrm{TeV}/35`$, and from 20 GeV to 40 GeV if $`m_{KK,1}\mathrm{TeV}/2`$ as suggested by the Brooks conjecture. (In all cases taking $`M_{}=1`$ TeV. Raising $`M_{}`$ raises $`T_{}`$.) So far we have concentrated on the spectrum of $`\mathrm{\Delta }_{LB}`$ appropriate for the spin-2 KK excitations. What about the spin-0(1) excitations? In both cases the detailed form of the Laplacian is modified. For example, in the spin-0 case the correct operator is the Liechnerowicz Laplacian, $$(\mathrm{\Delta }_{\mathrm{LL}}h)_{ij}=(D^kD_kh_{ij}+R_{ikjl}h^{kl}),$$ (10) where $`D_i`$ is the covariant derivative. The Mostow-Prasad rigidity theorem for CHM’s of dimension $`d3`$ tells us that $`\mathrm{\Delta }_{\mathrm{LL}}`$ has no zero modes. Although we know of no rigorous bounds for the first eigenvalue of this operator, an inspection of the generalized Selberg trace formulae supports the conjecture that the gap is of similar size to the Laplace-Beltrami case, a result that is physically reasonable. Finally for the spin-1 fluctuations $`h_{i\mu }`$ recall that these zero modes would correspond to KK gauge-bosons (the original motivation of Kaluza and Klein!), and are directly related to the continuous isometries of the compact space. But, as a result of the quotient by $`\mathrm{\Gamma }`$, CHM’s have no such isometries, and thus there are no massless KK gauge bosons. The non-zero KK modes once again have a mass gap that is at least as large as $`1/L`$ and is more likely close to $`1/R_c`$, as in the previous cases. Thus these additional types of fluctuation do not disturb our estimates. We have not yet addressed why it is almost automatic that there exist solutions of the form of eq. (1). Since CHM’s are just quotients of $`H^d`$ by a discrete identification under $`\mathrm{\Gamma }\mathrm{Isom}(H^d)`$, it is possible to find solutions of our form whenever there exists a uniform negative bulk cosmological constant (CC), given one constraint: $`R_cM_{}^1`$ and $`e^\alpha \mathrm{exp}\left((d1)L/R_c\right)1`$ must be realized consistently with our ansatz of a factorizable geometry with a static internal space, together with the vanishing of the 4D long-distance ($`L`$) CC. To ensure a static internal space, the small curvature radius of the internal space must be balanced in the field equations by the bulk CC, $`\mathrm{\Lambda }_{4+d}M_{}^{4+d}`$. Both these quantities contribute to the effective long-distance 4D CC, $`\mathrm{\Lambda }_4`$, on our brane, and typically do not cancel. Furthermore, one cannot just set $`\mathrm{\Lambda }_4`$ to zero by adjusting the tension or energy density $`f^4`$ of our 3-brane, because this requires $`f^4M_{}^4`$, violating our basic assumption that a low-energy effective theory is valid on the brane (and perturbing the geometry, possibly destroying our assumption that it is factorizable). To address this problem we must examine the form of the total 4D potential energy density $`V`$, which in the effective theory depends only on $`R_c`$ ($`e^\alpha `$ is an invariant), and which arises from the dimensional reduction of the full bulk and brane actions . For a 3-brane embedded in $`(4+d)`$ dimensions, the bulk and brane actions are respectively: $`S_{\mathrm{bulk}}`$ $`={\displaystyle d^{4+d}x\sqrt{\|g_{(4+d)}\|}\left(M_{}^{d+2}+\mathrm{\Lambda }_m\right)}`$ (11) $`S_{\mathrm{brane}}`$ $`={\displaystyle d^4x\sqrt{\|g_{(4)}^{\mathrm{induced}}\|}\left(f^4+\mathrm{}\right)},`$ (12) where $`_m`$ is the bulk matter field Lagrangian. Reduction of these actions gives a 4D potential energy density of the form $$V(R_c)=\mathrm{\Lambda }R_c^de^\alpha M_{}^4e^\alpha (M_{}R_c)^{d2}+W(R_c),$$ (13) to which we must add the brane tension $`f^4`$. The first two terms arise from the $`(4+d)`$ bulk CC term, and the curvature of the internal space. Now consider expanding $`W(R_c)`$, which comes from $`_m`$, as a Laurent series in $`R_c`$ $$W(R_c)=\underset{p}{}a_p\frac{M_{}^4}{(R_cM_{})^p},$$ (14) with dimensionless coefficients $`a_p`$. (Gauge or scalar field kinetic energies can give such terms with $`p>0`$ .) If the determination of the minimum is dominated by a competition between any two terms in $`V`$, then at this minimum $`VV_{\mathrm{min}}0`$. Moreover, $`V_{\mathrm{min}}`$ is enhanced by $`e^\alpha `$ over the “natural” value $`M_{}^4`$. However, the vanishing of the 4D CC demands $`V_{\mathrm{min}}\|_{\mathrm{tot}}=0`$. This cannot be achieved by adjusting the brane tension such that $`\|f^4\|M_{}^4`$. Fortunately there is an attractive alternative. If three or more $`R_c`$-dependent terms in $`V(R_c)`$ are all important at the minimum (for example the CC and curvature terms, and one of the matter terms from $`W`$) then we can tune the coefficients $`a_p`$ such that $`V_{\mathrm{min}}=0`$, without needing $`f^4M_{}^4`$. Thus, our basic assumptions remain consistent. Moreover, this tuning is particularly natural in our case precisely because we want the minimum to occur for a curvature radius close to the fundamental scale $`R_cM_{}^1`$, at which we expect the high-scale theory to produce many different terms that contribute roughly in an equal way. (This is exactly the opposite situation from the large flat extra dimension case where the minimum has to occur at a length scale much greater than $`M_{}^1`$.) This one fine-tuning is just the usual 4d CC problem, about which we have nothing to add. Having shown that there do exist solutions of our form, another significant result follows from this analysis. The most severe problem bedeviling the usual large extra dimension scenario is the radion moduli problem in the early universe . In our case this problem is much weakened. The radion, which is the light mode corresponding to dilations of the internal space, gets its mass from the stabilizing potential $`V(R_c)`$. Generally, in the flat extra dimension scenario, the radion mass $`m_r`$ is of size $`M_{}^2/M_P10^3`$ eV, so that it is very easily excited during the exit from inflation. Furthermore, since its couplings are $`1/M_P`$ suppressed, its life-time is longer than the age of the universe, so that it would unacceptably dominate our current expansion. In our case, however, the radion mass is greatly increased because the second derivative of the potential at its minimum is enhanced by a factor of $`e^\alpha `$, $`V_{\mathrm{min}}^{\prime \prime }=𝒪(e^\alpha M_{}^6)`$. Thus $$m_r^2=\frac{1}{2}\frac{R_c^2V^{\prime \prime }(R_c)}{e^\alpha M_{}^{d+2}R_c^d}\frac{1}{R_c^2},$$ (15) which is close to $`M_{}^2\mathrm{TeV}^2`$. Therefore, the radion lifetime is $`TM_P^2/M_{}^3`$, much shorter than in the case of flat extra dimensions, and only slightly longer than the age of the universe at nucleosynthesis, even if $`M_{}\mathrm{TeV}`$. Moreover, it is (comparatively) easy to dilute away any unwanted radion oscillations by a period of late inflation. While cosmologically and astrophysically much safer, models with internal compact hyperbolic spaces do have testable signatures accessible to collider experiments. Since KK modes abound close to the fundamental scale, Standard Model particle collisions with center-of-mass energies near this scale will result in the production of KK particles, detectable by a distinctive missing energy signature . Although this is qualitatively similar to the scenario of , the spectrum of KK modes one will see is quite distinctive. While the scale of KK masses is set by $`R_c^1`$, their ratios and multiplicities are in almost one-to-one correspondence with the topology of the internal manifold . A full exploration of these experimental signatures will require a more complete investigation of the spectrum of large CHM’s, in particular the issues of isospectrality and homophonicity of such manifolds. It is quite likely that such CHM’s have other implications for higher-dimensional physics. Besides a more detailed study of the question of radion stabilization, effects such as wavefunction scarring and brane-manifold dynamics are currently under investigation. We thank L. Alvarez-Gaume, R. Brooks, N. Cornish, S. Dimopoulos, C. Gordon, A. Gamburd, H. Mathur, J. Ratcliffe and J. Weeks for discussions, and the Stanford (JMR, GDS) and LBNL (JMR) theory groups for hospitality. Support was provided by the A.P. Sloan Foundation (JMR), the NSF (NK: NSF-PHY-9870115, GDS: NSF-CAREER), and the DOE (GDS, MT).
warning/0002/cond-mat0002435.html	ar5iv	text	# References ## Figure captions Normalized ’excitation profile’ $`A_n(\mathrm{\Omega },b)=A(\mathrm{\Omega },b)/A_{max}`$ versus $`\mathrm{\Omega }/\mathrm{\Omega }_{max}`$ for $`b=1,\mathrm{\Omega }_{max}\tau =0.736`$ (full line) and $`b=0.7,\mathrm{\Omega }_{max}\tau =0.185`$ (dashed line). Also shown is the imaginary part of the susceptibility, $`ϵ^{\prime \prime }(\mathrm{\Omega }\tau )`$, where $`\tau `$ is the relaxation time, (dotted line) for comparison. The lines for $`A_n(\mathrm{\Omega },0.7)`$ and $`ϵ^{\prime \prime }(\mathrm{\Omega }\tau )`$ have been shifted by $`0.5`$ and $`1.0`$ units, respectively. $`10^3\mathrm{\Delta }\mathrm{\Phi }(t)`$ versus rescaled time $`t/t_{max}`$ for (a) a heterogeneous and (b) a homogeneous scenario for various burn frequencies $`\mathrm{\Omega }`$ and $`\beta ^2E_P^2/\mathrm{\Delta }\mu ^2=0.05`$. In both cases the equilibrium response decays as $`\mathrm{\Phi }(t)=e^{(t/1s)^{0.7}}`$. The time $`t_{max}=1.0`$s in case (a) and $`t_{max}=0.095`$s in case (b). The used frequencies are: heterogeneous scenario: $`\mathrm{\Omega }\tau =5.0`$ (1), $`1.0`$ (2), $`0.2`$ (3) $`0.1`$ (4); homogeneous scenario: $`\mathrm{\Omega }\tau =5.0`$ (1), $`1.0`$ (2), $`0.1`$ (3), $`0.05`$ (4); Experimental data in (a) are adapted from ref.. Here, $`\mathrm{\Omega }\tau `$=1.02 (2) and $`\mathrm{\Omega }\tau `$=0.203 (3).
warning/0002/hep-ph0002147.html	ar5iv	text	# Updated Global Analysis of the Atmospheric Neutrino Data in terms of neutrino oscillations ## I Introduction Together with the solar neutrino problem the atmospheric neutrino anomaly constitutes the second evidence for physics beyond the Standard Model. Indeed a large number of observations have been performed concerning electron– and muon–neutrino fluxes produced by hadronic showers initiated by cosmic–ray interactions in the upper atmosphere. Apart from the first iron–calorimeter detectors , all experiments, which also entail different detection techniques, have steadily reported a deficit of the collected number of events with respect to the theoretical expectations . Although the knowledge of the fluxes of atmospheric neutrinos is affected by uncertainties which range from about 20% to 30%, the expected ratio $`R(\mu /e)`$ of the muon neutrino ($`\nu _\mu +\overline{\nu }_\mu `$) over the electron neutrino flux ($`\nu _e+\overline{\nu }_e`$) is known with much better confidence, since the uncertainties associated with the absolute fluxes largely cancel out. It is fair to ascribe an overall uncertainty less than about 5% to the calculated $`R(\mu /e)`$ ratio. Since this ratio is measured to be substantially smaller than the expectations, one faces an anomaly which only seems possible to account for in terms of non–standard neutrino properties. Super–Kamiokande high statistics observations indicate that the deficit in the total ratio $`R(\mu /e)`$ is due to the neutrinos arriving in the detector at large zenith angles. They also show that the $`e`$-like events do not present any compelling evidence of a zenith-angle dependent suppression, while the $`\mu `$-like event rates are substantially suppressed at large angles. Different explanations to these features have been proposed and discussed in the literature . The simplest and most direct possibility is represented by the oscillation of muon neutrinos $`\nu _\mu `$ into either a $`\nu _\tau `$ or a sterile neutrino $`\nu _s`$ . The oscillation hypothesis provides a very good explanation for this smaller-than-expected ratio, which is also simple and well-motivated theoretically. In this paper we present our updated analysis of all the available data on atmospheric neutrinos in terms of neutrino oscillation. We include in the analysis all the existing experimental results obtained so far. In addition to the previous data samples of Frejus, Nusex, IMB and Kamioka experiments, we also consider the recent Soudan2 data which refers to an exposure of 4.6 kton-yr, the full data set of the 52 kton-yr of Super–Kamiokande, including both contained events and upgoing muon data, and finally the results on upgoing muons reported from the MACRO and Baksan experiments. We critically discuss the analysis of the various individual data sets. Moreover we consider the combined information that can be derived from all the experimental evidences so–far obtained. We hope this will contribute to a more complete understanding of the atmospheric neutrino anomaly in the framework of a global analysis based on a self–consistent theoretical calculation of the event rates. This has the advantage of calibrating all relevant elements in the analysis, such as theoretical atmospheric neutrino fluxes, in a consistent way. This allows different experiments to be compared in a meaningful way. A new element which we now add to our previous investigations is the inclusion of the upgoing muon data and the corresponding theoretical flux determinations. We anticipate that from the statistical analysis the $`\nu _\mu \nu _\tau `$ emerges as the oscillation channel which provides the best agreement with the combined data. The $`\nu _\mu \nu _s`$ channels, although slightly disfavoured, cannot be statistically ruled out on the basis of the global fit to the full set of observables. The outline of the paper is the following: in Sect. II we briefly recall the theoretical calculation of the event rates for contained events and upgoing muon fluxes, as well as the calculation of the oscillation probabilities. Sect. III discusses the statistical approach of our analysis and reports on the results of the fits. Sect. IV presents a discussion of the results and the conclusions. We include more details on the statistical analysis in an Appendix. ## II Atmospheric Neutrino Induced Events in Underground Detectors Atmospheric neutrinos can be detected in underground experiments by direct observation of their charged current interaction inside the detector. These are the so-called contained events, which can be further classified into fully contained events when the charged lepton (either electron or muon) produced by the neutrino interaction does not escape the detector, and partially contained muons when the muon, produced inside, leaves the detector. In the case of Kamiokande and Super–Kamiokande, the contained data sample is further divided into sub-GeV events with visible energy below 1.2 GeV and multi-GeV events when the lepton energy is above such cutoff. On the average, sub-GeV events arise from neutrinos of several hundreds of MeV while multi-GeV events are originated by neutrinos with energies of the order of several GeV. Higher energy muon neutrinos and anti-neutrinos can also be detected indirectly by observing the muons produced by charged current interactions in the vicinity of the detector. These are the so called upgoing muons. Should the muon stop inside the detector, it will be classified as a “stopping” muon, while if the muon track crosses the full detector the event is classified as a “through-going” muon. Typically stopping muons arise from neutrinos of energies around ten GeV while through-going muons are originated by neutrinos with energies of the order of hundred GeV. ### A Contained events At present, six underground experiments have collected data on contained events. Three of them, Kamiokande IMB and Super–Kamiokande use water-Cerenkov detectors, while the other three, Fréjus , NUSEX and Soudan2 are iron calorimeter detectors. For a given neutrino conversion mechanism, the expected number of $`\mu `$-like and $`e`$-like contained events, $`N_\alpha `$, $`\alpha =\mu ,e`$ can be computed as: $$N_\mu =N_{\mu \mu }+N_{e\mu },N_e=N_{ee}+N_{\mu e},$$ (1) where $`N_{\alpha \beta }`$ $`=`$ $`n_tT{\displaystyle \frac{d^2\mathrm{\Phi }_\alpha }{dE_\nu d(\mathrm{cos}\theta _\nu )}\kappa _\alpha (h,\mathrm{cos}\theta _\nu ,E_\nu )P_{\alpha \beta }\frac{d\sigma }{dE_\beta }\epsilon (E_\beta )𝑑E_\nu 𝑑E_\beta d(\mathrm{cos}\theta _\nu )𝑑h}`$ (2) and $`P_{\alpha \beta }`$ is the conversion probability of $`\nu _\alpha \nu _\beta `$ for given values of $`E_\nu ,\mathrm{cos}\theta _\nu `$ and $`h`$, i.e., $`P_{\alpha \beta }P(\nu _\alpha \nu _\beta ;E_\nu ,\mathrm{cos}\theta _\nu ,h)`$. In the Standard Model (SM), the only non-zero elements are the diagonal ones, i.e. $`P_{\alpha \alpha }=1`$ for all $`\alpha `$. In Eq.(2) $`n_t`$ denotes the number of targets, $`T`$ is the experiment running time, $`E_\nu `$ is the neutrino energy and $`\mathrm{\Phi }_\alpha `$ is the flux of atmospheric neutrinos of type $`\alpha =\mu ,e`$ for which we use the Bartol flux; $`E_\beta `$ is the final charged lepton energy and $`\epsilon (E_\beta )`$ is the detection efficiency for such charged lepton; $`\sigma `$ is the neutrino-nucleon interaction cross section, and $`\theta _\nu `$ is the angle between the vertical direction and the incoming neutrinos ($`\mathrm{cos}\theta _\nu `$=1 corresponds to the down-coming neutrinos). In Eq. (2), $`h`$ is the slant distance from the production point to the sea level for $`\alpha `$-type neutrinos with energy $`E_\nu `$ and zenith angle $`\theta _\nu `$. Finally, $`\kappa _\alpha `$ is the slant distance distribution, normalized to one . For the angular distribution of events we integrate in the corresponding bins for $`\mathrm{cos}\theta _\beta `$ where $`\theta _\beta `$ is the angle of the detected lepton, taking into account the opening angle between the neutrino and the charged lepton directions as determined by the kinematics of the neutrino interaction. In average the angle between the directions of the final-state lepton and the incoming neutrino ranges from $`70^{}`$ at 200 MeV to $`20^{}`$ at 1.5 GeV. The neutrino fluxes, in particular in the sub-GeV range, depend on the solar activity. In order to take this fact into account, we use in Eq. (2) a linear combination of atmospheric neutrino fluxes $`\mathrm{\Phi }_\alpha ^{max}`$ and $`\mathrm{\Phi }_\alpha ^{min}`$ which correspond to the most active Sun (solar maximum) and quiet Sun (solar minimum), respectively, with different weights which depend on the running period of each experiment . Following Ref. we explicitly verify in our present reanalysis the agreement of our predictions with the experimental Monte Carlo predictions, leading to a good confidence in the reliability of our results for contained events. ### B Upward Going Muons Several experiments have obtained data on these neutrino induced muons. In our analysis we consider the results from three experiments: Super–Kamiokande , MACRO and Baksan. They present their upgoing muon data in the form of measured muon fluxes. We obtain the effective muon fluxes for both stopping and through-going muons by convoluting the survival $`\nu _\mu `$ probabilities (calculated as in Sect. II.C) with the corresponding muon fluxes produced by the neutrino interactions with the Earth. We include the muon energy loss during propagation both in the rock and in the detector according to Refs. and we take into account also the effective detector area for both types of events, stopping and through-going. Schematically $$\mathrm{\Phi }_\mu (\theta )_{S,T}=\frac{1}{A(L_{min},\theta )}_{E_{\mu ,min}}^{\mathrm{}}\frac{d\mathrm{\Phi }_\mu (E_\mu ,\mathrm{cos}\theta )}{dE_\mu d\mathrm{cos}\theta }A_{S,T}(E_\mu ,\theta )𝑑E_\mu ,$$ (3) where $`{\displaystyle \frac{d\mathrm{\Phi }_\mu }{dE_\mu d\mathrm{cos}\theta }}`$ $`=`$ $`N_A{\displaystyle _{E_\mu }^{\mathrm{}}}𝑑E_{\mu 0}{\displaystyle _{E_{\mu 0}}^{\mathrm{}}}𝑑E_\nu {\displaystyle _0^{\mathrm{}}}𝑑X{\displaystyle _0^{\mathrm{}}}𝑑h\kappa _{\nu _\mu }(h,\mathrm{cos}\theta ,E_\nu )`$ (5) $`{\displaystyle \frac{d\mathrm{\Phi }_{\nu _\mu }(E_\nu ,\theta )}{dE_\nu d\mathrm{cos}\theta }}P_{\mu \mu }{\displaystyle \frac{d\sigma (E_\nu ,E_{\mu 0})}{dE_{\mu 0}}}F_{rock}(E_{\mu 0},E_\mu ,X)`$ $`N_A`$ is the Avogadro number, $`E_{\mu 0}`$ is the energy of the muon produced in the neutrino interaction and $`E_\mu `$ is the muon energy when entering the detector after traveling a distance $`X`$ in the rock. $`\mathrm{cos}\theta `$ labels both the neutrino and the muon directions which to a very good approximation at the relevant energies are collinear. We denote by $`F_{rock}(E_{\mu 0},E_\mu ,X)`$ the function which characterizes the energy spectrum of the muons arriving the detector. Following standard practice in our calculations we use the analytical approximation obtained by neglecting the fluctuations during muon propagation in the Earth. In this case the three quantities $`E_{\mu 0}`$, $`E_\mu `$, and $`X`$ are not independent: $$_0^{\mathrm{}}F_{rock}(E_{\mu 0},E_\mu ,X)𝑑X=\frac{1}{d_\mu (E_\mu )/dX},$$ (6) where $`d_\mu (E_\mu )/dX`$ is the average muon energy loss due to ionization, bremsstrahlung, $`e^+e^{}`$ pair production and nuclear interactions in the rock according to Refs. . Equivalently, the pathlength traveled by the muon inside the Super–Kamiokande detector is given by the muon range function in water $$L(E_\mu )=_0^{E_\mu }\frac{1}{d_\mu (E_\mu ^{})/dX}𝑑E_\mu ^{}.$$ (7) In Eq.(3) $`A(L_{min},\theta )=A_S(E_\mu ,\theta )+A_T(E_\mu ,\theta )`$ is the projected detector area for internal path-lengths longer than $`L_{min}`$ which for Super–Kamiokande is $`L_{min}`$ = 7 m. Here $`A_S`$ and $`A_T`$ are the corresponding effective areas for stopping and through-going muon trajectories. For Super–Kamiokande we compute these effective areas using the simple geometrical picture given in Ref. . For a given angle, the threshold energy cut for Super–Kamiokande muons is obtained by equating Eq.(7) to $`L_{min}`$, i.e. $`L(E_\mu ^{\mathrm{t}h})=L_{min}`$. In contrast with Super–Kamiokande, Baksan and MACRO present their results as muon fluxes for $`E_\mu >1`$ GeV, after correcting for detector acceptances. Therefore in this case we compute the expected fluxes as Eqs. (3) and (5) but without the inclusion of the effective areas. We have explicitly verified the agreement of our predictions for upgoing muons with the experimental Monte Carlo predictions from Super–Kamiokande, Baksan and MACRO. The agreement can be observed by comparing our Standard Model predictions for the angular distributions in Fig. 6 with the corresponding distributions in Refs. and . We find an agreement to the 5% and 1% level, respectively. ### C Conversion Probabilities For definiteness we assume a two-flavour scenario. The oscillation probabilities are obtained by solving the Schröedinger evolution equation of the $`\nu _\mu \nu _X`$ system in the Earth–matter background (in our case, $`X=\tau `$ or $`s`$). For neutrinos this equation reads: $$i\frac{\text{d}}{\text{d}t}\left(\begin{array}{c}\nu _\mu \\ \nu _X\end{array}\right)=\left(\begin{array}{cc}H_\mu & H_{\mu X}\\ H_{\mu X}& H_X\end{array}\right)\left(\begin{array}{c}\nu _\mu \\ \nu _X\end{array}\right)$$ (8) where $`H_\mu `$ $`=`$ $`V_\mu +{\displaystyle \frac{\mathrm{\Delta }m^2}{4E_\nu }}\mathrm{cos}2\theta ,`$ (9) $`H_X`$ $`=`$ $`V_X{\displaystyle \frac{\mathrm{\Delta }m^2}{4E_\nu }}\mathrm{cos}2\theta ,`$ (10) $`H_{\mu X}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }m^2}{4E_\nu }}\mathrm{sin}2\theta ,`$ (11) with the following definition of the neutrino potentials in matter $`V_\mu =V_\tau `$ $`=`$ $`{\displaystyle \frac{\sqrt{2}G_F\rho }{M}}({\displaystyle \frac{1}{2}}Y_n),`$ (12) $`V_s`$ $`=`$ $`0.`$ (13) Here $`G_F`$ is the Fermi constant, $`\rho `$ is the matter density in the Earth, $`M`$ is the nucleon mass, and $`Y_n`$ is the neutron fraction. For anti-neutrinos the signs of potentials $`V_X`$ should be reversed. In the previous Eqs., $`\theta `$ is the mixing angle between the two mass eigenstate neutrinos of masses $`m_1`$ and $`m_2`$. We define $`\mathrm{\Delta }m^2=m_2^2m_1^2`$ in such a way that if $`\mathrm{\Delta }m^2>0`$ $`(\mathrm{\Delta }m^2<0)`$ the neutrino with largest muon-like component is heavier (lighter) than the one with largest $`X`$–like component. Since for the $`\nu _\mu \nu _\tau `$ case there is no matter effect, the solution of Eq.(8) is straightforward and the probability takes the well-known vacuum form $$P_{\mu \mu }=1\mathrm{sin}^2(2\theta )\mathrm{sin}^2\left(\frac{\mathrm{\Delta }m^2L}{2E_\nu }\right).$$ (15) where $`L`$ is the path-length traveled by neutrinos of energy $`E_\nu `$. In the case of $`\nu _\mu \nu _s`$, the presence of the matter potentials requires a numerical solution of the evolution equations in order to obtain the oscillation probabilities $`P_{\alpha \beta }`$, which are different for neutrinos and anti-neutrinos because of the reversal of sign of the $`V_X`$’s. In our calculations, for the matter density profile of the Earth we have used the approximate analytic parametrization given in Ref. of the PREM model of the Earth . Notice that for the $`\nu _\mu \nu _s`$ case, there are two possibilities depending on the sign of $`\mathrm{\Delta }m^2`$. For $`\mathrm{\Delta }m^2>0`$ the matter effects enhance neutrino oscillations while depress anti-neutrino oscillations, whereas for the other sign ($`\mathrm{\Delta }m^2<0`$) the opposite holds. Finally, we have not considered oscillations of $`\nu _\mu `$’s into electron neutrinos. This possibility is nowadays ruled out since $`\nu _\mu \nu _e`$ oscillations cannot describe the measured angular dependence of muon-like contained events . Moreover the most favoured range of masses and mixings for this channel has already been excluded by the negative results from the CHOOZ reactor experiment . ## III Atmospheric Neutrino Data Fits In this Section, we describe the fitting procedure we employ in order to determine the atmospheric neutrino parameters $`\mathrm{sin}^2(2\theta )`$ and $`\mathrm{\Delta }m^2`$ for the two conversion channels of interest and we present our results. ### A Statistical analysis We first notice that we rely on the separate use of the event numbers (not of their ratios), paying attention to the correlations between the sources of errors in the muon and electron predictions as well as the correlations among the errors of the different energy data samples. A detailed description of the steps required in order to determine the allowed regions of oscillation parameters was given in Ref. . Following Refs. , we define a $`\chi ^2`$–function as $$\chi ^2\underset{I,J}{}(N_I^{DATA}N_I^{TH})(\sigma _{DATA}^2+\sigma _{TH}^2)_{IJ}^1(N_J^{DATA}N_J^{TH}),$$ (16) where $`I`$ and $`J`$ stand for any combination of experimental data sets and event-types considered, i.e, $`I=(A,\alpha )`$ and $`J=(B,\beta )`$. The latin indexes $`A,B`$ stand for the different experiments or different data samples in a given experiment. The greek indexes denote electron–type or muon–type events, i.e, $`\alpha ,\beta =e,\mu `$. In Eq. (16), $`N_I^{TH}`$ stands for the predicted number of events (or for the predicted value of the flux, in the case of upgoing muons) calculated as discussed in the previous Section, whereas $`N_I^{DATA}`$ is the corresponding experimental measurement. In Eq. (16), $`\sigma _{DATA}^2`$ and $`\sigma _{TH}^2`$ are the error matrices containing the experimental and theoretical errors, respectively. They can be written as $$\sigma _{IJ}^2\sigma _\alpha (A)\rho _{\alpha \beta }(A,B)\sigma _\beta (B),$$ (17) where $`\rho _{\alpha \beta }(A,B)`$ is the correlation matrix containing all the correlations between the $`\alpha `$-like events in the $`A`$-type experiment and $`\beta `$-like events in $`B`$-type experiment, whereas $`\sigma _\alpha (A)`$ and $`\sigma _\beta (B)`$ are the errors for the number of $`\alpha `$ and $`\beta `$-like events in $`A`$ and $`B`$ experiments, respectively. The dimensionality of the error matrix varies depending on the combination of experiments included in the analysis. We have computed the correlation matrix $`\rho _{\alpha \beta }(A,B)`$ as described in Ref. . In the case of contained events, a detailed discussion of the errors and correlations used in our analysis can be found in Ref. , which, for the sake of clarity and completeness, we summarize in the Appendix. In the same Appendix, we include the discussion of the errors and correlations employed for the upgoing muon data analysis. With all the definitions discussed above, we can calculate the $`\chi ^2`$ in Eq. (16) as a function of the neutrino parameters. By minimization of the $`\chi ^2`$ with respect to $`\mathrm{sin}^2(2\theta )`$ and $`\mathrm{\Delta }m^2`$, we determine our best fit results, while the allowed regions are determined by the following conditions: $`\chi ^2\chi _{\mathrm{m}in}^2+4.61(6.1)[9.21]`$ for a confidence level (CL) of $`90(95)[99]`$ %, respectively. The data sets employed in the statistical analysis are the following: (i) Frejus, Nusex, IMB and Soudan2 data sets, which refer to low-energy, contained events. For each experiment, the total rate for $`e`$-like and for $`\mu `$-like events is reported. We jointly analyze these data (and hereafter we collectively denote them as FISN); (ii) $`e`$-like and $`\mu `$-like Kamiokande data, including the sub-GeV event rate and a 5-bin zenith-angle distribution for the multi-GeV data; (iii) Super–Kamiokande data, again comprising $`e`$-like and $`\mu `$-like contained events, arranged into sub-GeV and multi-GeV samples, each of which given as a 5-bin zenith-angle distribution; and the up–going data including the stopping muon flux (5 bins in zenith-angle) and the through-going muon flux (10 angular bins); (iv) MACRO and Baksan upgoing muons samples, each one with 10 angular bins. In the discussion of our results, along with presenting the results of the separate analyses of the single data sets, we analyze various possible combinations of data samples in order to develop an understanding of the relevance of the different sets of data, which, we recall, refer to events characterized by different properties. For instance, contained events are produced by neutrinos of relatively low energies (below a few GeV), while the upgoing muon fluxes are originated by neutrinos whose energies cover a much wider range, from a few GeV to hundreds of GeV. This feature, combined also with the angular distributions, allows one to test the energy dependence of the oscillation probabilities in the different channels. ### B Results of the analysis As a first result of our analysis, we show the (in)consistency of the data with the Standard Model hypothesis. The first column of Table I reports the values of the $`\chi ^2`$ function in the absence of new physics, as obtained with our prescriptions and calculated for different combinations of atmospheric data sets. We notice that all the data sets clearly indicate deviation from the standard model. The global analysis, which refers to the full combination of all data sets, gives a value of $`\chi _{SM,\mathrm{g}lobal}^2=`$ 214/(75 d.o.f) corresponding to a CL of $`3\times 10^{15}`$. This indicates that the Standard Model can be safely ruled out. Instead, the $`\chi ^2`$ for the global analysis decreases to $`\chi _{min}^2=`$ 74/(73 d.o.f) (45 % CL) when assuming the $`\nu _\mu \nu _\tau `$ oscillation hypothesis. Table I reports the minimum values of $`\chi ^2`$ and the resulting best fit points for the various oscillation channels and data sets considered. The corresponding allowed regions for the oscillation parameters at 90, 95 and 99 % CL are depicted in Figs. 1345 and 7. The zenith-angle distributions referring to the Super–Kamiokande data and to the upgoing muon measurements of MACRO and Baksan are plotted in Figs. 26 and 8. Note that no uncertainties are shown in the plots for the theoretical predictions, while experimental data errors are explicitly displayed. We now turn to the discussion of the main results. #### 1 Contained Events The allowed regions for contained events are displayed in Figs. 1 and 3. In all figures the best fit points are marked with a star and can be read from Table I together with the corresponding value of the $`\chi _{\mathrm{m}in}^2`$ for each case. In the first column of Fig. 1 we show the results for the combination of the “unbinned” FISN data, i.e. the total rates from the Frejus, IMB, Nusex and Soudan experiments. Since no angular information in present in this case, no lower limit on the oscillation length can be derived and the regions extend to arbitrary large mass differences. We notice, from Table I, that the best fit results, although definitely improved with respect to the SM case, do not show, for each oscillation channel, a CL better than 4%. We remind that the FISN data sample contains the Frejus and Nusex data, which are by themselves compatible with the SM. This fact plays a role in maintaining a relatively high $`\chi ^2`$ even in the oscillation cases. The second column of Fig. 1 corresponds to the combination of contained sub-GeV (unbinned) and multi-GeV (including the angular dependence) data from the Kamiokande experiment. In this case, the angular distribution of the multi-GeV events plays an important role in determining an upper limit for $`\mathrm{\Delta }m^2`$, for a given value of the mixing angle. The region which is obtained overlaps with the previous one from the FISN data sample, and indicates a preference for values of $`\mathrm{\Delta }m^210^3`$ for all the oscillation channels. For the Kamiokande data, the agreement with the oscillation hypothesis is at the level of $`70`$% CL. Finally, in the third column of Fig. 1 we plot the allowed regions for the combination of the angular distributions of the sub-GeV and multi-GeV 52 kton-yr Super–Kamiokande data. The best fit for the $`\nu _\mu \nu _\tau `$ hypothesis has a very high confidence level: $`\chi _{\mathrm{m}in}^2=`$ 8.9/(18 d.o.f.), which translates into 96% CL, substantially improved with respect to the Kamiokande data sets. In the case of oscillation to sterile neutrinos, the CL is 79%, slightly higher than for the Kamiokande data alone. The angular distributions of the Super–Kamiokande sub-GeV and multi-GeV $`\mu `$-like events are shown in Fig. 2. We also show in the figure the $`e`$–like distributions to illustrate that they do no show any evidence of deviation from the standard model prediction. The agreement of the $`e`$–like distributions with the expectations from the standard model has become more tight as the Super–Kamiokande collaboration has been increasing their data sample. This has translated into an overall improvement of the CL for the oscillation hypothesis into channels not involving electron neutrinos. Conversely solutions involving oscillations into electron neutrinos have become more disfavoured . From Fig. 2 one sees that there is a strong evidence of a depletion in the event rates with respect to the SM expectation. We notice that the zenith-angle distributions obtained with the best-fit neutrino parameters are able to reproduce the data to a high level of agreement. From Fig. 1, one can notice that in all the $`\nu _\mu \nu _s`$ channels where matter effects play a role, the range of acceptable $`\mathrm{\Delta }m^2`$ is slightly shifted towards larger values, when compared with the $`\nu _\mu \nu _\tau `$ case. This follows from looking at the relation between mixing in vacuo and in matter. In fact, away from the resonance region, independently of the sign of the matter potential, there is a suppression of the mixing inside the Earth. As a result, there is a lower cut in the allowed $`\mathrm{\Delta }m^2`$ value, and this lies higher than what is obtained in the case of the $`\nu _\mu \nu _\tau `$ channel. Also, for the $`\nu _\mu \nu _s`$ case with $`\mathrm{\Delta }m^2>0`$ the matter effects enhance neutrino oscillations while depress anti-neutrino oscillations. Since atmospheric fluxes are dominantly neutrinos, smaller mixing angle values can lead to the same $`\nu _\mu `$ suppression and the region extends to smaller mixing angles in the $`\mathrm{\Delta }m^2`$ region where the matter effects are important for the relevant neutrino energies. The opposite holds for $`\mathrm{\Delta }m^2<0`$. In this case the matter effects depress neutrino oscillations, and therefore larger mixing angles are needed to account for the observed deficit. As a consequence, the allowed regions become smaller (in angle) for this channel. When comparing our results with the corresponding analysis presented by the Super–Kamiokande Collaboration on their data sets, we find good agreement on the allowed parameters although in general our regions are slightly larger. We also find that, for the $`\nu _\mu \nu _\tau `$ channel, the allowed regions for the contained events are shifted towards slightly lower $`\mathrm{\Delta }m^2`$ values with respect to Super–Kamiokande analysis. We have traced this difference back to the sub-GeV sample and to our use of different neutrino fluxes. Notice that sub-GeV events are most sensitive to details in the neutrino fluxes from the different calculations, such as the modeling of the geo–magnetic cut–off . The expected angular distribution in the absence of oscillations from the Super–Kamiokande Monte–Carlo using Honda fluxes is slightly “flatter” then our calculations for the sub–GeV event distributions. We have verified that if we normalize our results to the Super–Kamiokande Monte–Carlo predictions, then the allowed region we obtain is shifted to slightly larger $`\mathrm{\Delta }m^2`$ values in agreement with the Super–Kamiokande analysis. Finally, the global analysis of the combination of all the above data sets is shown in the first column of Fig. 3 (notice a change of scale between this figure and Fig. 1). The allowed regions, as well as the best fit points, are clearly driven by the high statistics Super–Kamiokande data. The results show, as expected from the above discussion, that the $`\nu _\mu \nu _\tau `$ case prefers slightly lower values of $`\mathrm{\Delta }m^2`$ (from a few 10<sup>-4</sup> to a few 10<sup>-3</sup> eV<sup>2</sup>) when compared to the $`\nu _\mu \nu _s`$ cases, where values in excess of 10<sup>-3</sup> eV<sup>2</sup> are obtained. In all the cases, almost-maximal mixing is statistically preferred. Concerning the quality of the fits, as seen in Table I, in the full combination the $`\nu _\mu \nu _\tau `$ channel gives a slightly better fit ($`\chi _{min}^2=37/38`$) than $`\nu _\mu \nu _s`$ ($`\chi _{min}^2=40/38`$) although the difference is not statistically significant (51% CL versus 38% CL). #### 2 Upward Going Muons Our results for the allowed regions for upgoing muon events are displayed in Figs. 4 for the Super–Kamiokande data on stopping and through-going muons and in Fig. 5 for the MACRO and Baksan experiments. The global analysis on all the data samples on upgoing muons is shown in the second column of Fig. 3. From the analysis of Super–Kamiokande data on stopping muons we find that the $`\chi ^2`$-function is substantially flat for $`\mathrm{\Delta }m^2`$ values above 10<sup>-3</sup> eV<sup>2</sup> and $`\mathrm{sin}^2(2\theta )\stackrel{>}{}\mathrm{\hspace{0.25em}0.5}`$ and therefore the allowed regions are open from above, also at 90 % CL. This is a consequence of the fact that the stopping muons data sample by its own is consistent with a global reduction of the neutrino flux with no specific angular dependence. This feature can be observed by comparing the angular distribution of the data points in the first panel of Fig. 6 and the corresponding prediction in the absence of oscillations. This behaviour is what is expected from oscillations with short oscillation lengths (high $`\mathrm{\Delta }m^2`$). In contrast, the through-going Super–Kamiokande muon sample indicates a somewhat angular-dependent suppression. The second panel of Fig. 6 shows that the reduction is clearly larger for vertically-coming muons than for those arriving from the horizon. As a consequence, the corresponding allowed regions turn out to be closed from above, i.e. for large $`\mathrm{\Delta }m^2`$. This is plotted in the second column of Fig. 4, whose comparison with the first column shows the general agreement between the stopping and through–going muons data samples. This is an interesting property in favour of the neutrino oscillation hypothesis, since, as noted before, stopping and through–going muon events originate from very different parent–neutrino energies. Finally, the combination of both stopping and thru-going muon events in Super–Kamiokande is shown in the third column of Fig. 4. We find a very good agreement between our results and the corresponding ones given by the Super–Kamiokande Collaboration. The best fit points in ours and the Super–Kamiokande analyses are very similar, although also in this case our regions are slightly larger. From the angular distributions of Fig. (6), we can observe that due to matter effects the distribution for upgoing muons in the case of $`\nu _\mu \nu _s`$ are flatter than for $`\nu _\mu \nu _\tau `$ . Since the data show a somehow steeper angular dependence, a better description in terms of $`\nu _\mu \nu _\tau `$ oscillations is found. From Table I, we see that the upward going muon fit in this channel is indeed better, although the statistical preference of this channel over the sterile case (45% CL for $`\nu _\mu \nu _\tau `$ as compared to about 30% CL for $`\nu _\mu \nu _s`$) is not overwhelming. In the case of both MACRO and Baksan experiments we generally find a less significant statistical agreement between the data and the theoretical evaluations. Both data sets are clearly inconsistent with the SM predictions, as indicated by their high $`\chi _{SM}^2`$ values in Table I: the CL are $`7\times 10^4`$ for MACRO and $`5\times 10^3`$ for Baksan. However, the quality of the fits does not strongly improve when interpreted in terms of neutrino oscillations. We find that the best agreement occurs for the oscillation in the $`\nu _\mu \nu _\tau `$ channel for the MACRO experiment. In this case we obtain a 5% CL and an allowed region in the parameter space, which is shown in the first panel of Fig. 5. This result is fully consistent with the analysis performed by the MACRO Collaboration on their data set . We notice that, as can be seen from the angular distribution of Fig. 6, the flux in the fourth bin ($`0.7<\mathrm{cos}\theta _Z<0.6`$) is at least 2–$`\sigma `$ above the expectation in the presence of oscillations. When removing this point from the statistical analysis, the quality of the fit for the $`\nu _\mu \nu _\tau `$ oscillation channel improves to 19% CL, while affecting very little the allowed region. In contrast, for the sterile cases, we encounter a very low statistical confidence (about 1% CL), even after removing the fourth bin data point. As a result the corresponding allowed oscillation parameter regions are not very meaningful, since the best fit point has a large $`\chi _{\mathrm{m}in}^2`$ value. For this reason we have not reproduced them here. We only comment that the $`\chi ^2`$ is very flat in the neutrino oscillation parameters, without a clear indication of a statistically preferred region. In the $`\mathrm{\Delta }m^2>0`$ case, two almost degenerate minima are found, one for small and one for large mixing angle, as can be seen in Table I). For Baksan we find no clear preference for the oscillation hypothesis with respect to the Standard Model case, although in the sterile channels the best fits with oscillation turn out to be slightly better than for the no-oscillation case. We reproduce in Fig. 5 the allowed region for the $`\nu _\mu \nu _s`$ ($`\mathrm{\Delta }m^2>0`$) case, which is the one with lowest $`\chi ^2`$. In the case of $`\nu _\mu \nu _\tau `$ and $`\nu _\mu \nu _s`$ ($`\mathrm{\Delta }m^2<0`$), the allowed regions are very similar to the one depicted in Fig. 5. We wish to warn that, due to the low statistical significance of the best fit results, these regions should be taken only as indicative. To conclude this Section, we show in the right column of Fig. 3 the result of our analysis for the combination of all the data on upgoing muons discussed above. In the case of $`\nu _\mu \nu _\tau `$, one notices that the Super–Kamiokande and MACRO results give consistent and similar allowed regions, while Baksan gives compatible but not strongly constraining results. As a consequence, the combined analysis gives a region which is intermediate to the Super–Kamiokande and MACRO ones. The best fits of the combined upgoing neutrino analysis have a low CL (around 1% for all the oscillation channels), a result which is mainly driven by the high $`\chi _{\mathrm{m}in}^2`$ values of MACRO and Baksan (we recall that Super–Kamiokande upgoing-muon data alone indicate a preference for neutrino oscillation at a level always better than 25%). However, the global analysis of the upgoing muons data disfavours the Standard Model at the $`3\times 10^5`$ level. #### 3 Global Analysis Let us now move to the discussion of the comparison and combination of contained events with upgoing muons fluxes. We observe from Fig. 3 that the allowed regions obtained from both types of analyses are fully consistent between themselves for all the oscillation channels. For the cases of $`\nu _\mu \nu _s`$ the allowed region for the contained events lies always inside the corresponding regions allowed by the upgoing muon analysis. For the $`\nu _\mu \nu _\tau `$ channel, we find that the region for contained events extends to lower values of $`\mathrm{\Delta }m^2`$, when compared to the region for upgoing muons. In figure 7 we display the allowed regions after combining together contained and upward going muon data. In the first column we show the results when only including the Super–Kamiokande data. The second column shows the corresponding results when data from all experiments are included, while in the third column we show the allowed regions when only the results from experiments observing some evidence of neutrino oscillation are included. The general behaviour is that when including the results from all experiments the regions become slightly larger than those obtained from the analysis of the Super–Kamiokande data alone. On the other hand, as expected, the results become more restrictive when only the data from experiments observing some evidence for oscillations are included. In Table I we list the values of the best fit points for the various cases. Our results from the combined analysis show that the channel $`\nu _\mu \nu _\tau `$ gives a better fit to the data than oscillations into sterile neutrinos. The difference for the global analysis is $`\chi _{min}^2=`$ 74/(73 d.o.f.) (45 % CL) for the active case versus $`\chi _{min}^2=`$ 90/(73 d.o.f.) (8.5 % CL) for oscillations into sterile neutrinos with $`\mathrm{\Delta }m^2<0`$ and $`\chi _{min}^2=`$ 86/(73 d.o.f.) (14 % CL) for $`\nu _\mu \nu _s`$ with $`\mathrm{\Delta }m^2>0`$. This difference in the quality of the description is still maintained when some of the negative-result experiments are excluded from the analysis (by “negative-result experiments” we mean experiments which, from our statistical analysis, do not give a clear evidence of neutrino oscillation, i.e. have relatively high $`\chi _{\mathrm{m}in}^2`$ values). All of these features can be seen in Table I. When removing from the analysis the Frejus, Nusex and Baksan data points, we do not obtain an improvement for the sterile cases, while for the active case the CL is increased to 58 % ($`\chi _{min}^2=`$ 58/(61 d.o.f.)). When also MACRO is removed, we obtain higher CL also for the sterile cases, but the $`\nu _\mu \nu _\tau `$ hypothesis remains as the best option: $`\chi _{min}^2=`$ 41/(51 d.o.f.) (84 % CL) for the active case, $`\chi _{min}^2=`$ 51/(51 d.o.f.) (47 % CL) for $`\nu _\mu \nu _s`$ with $`\mathrm{\Delta }m^2<0`$ and $`\chi _{min}^2=`$ 50/(51 d.o.f.) (51 % CL) for $`\nu _\mu \nu _s`$ with $`\mathrm{\Delta }m^2>0`$. In conclusion, we find that the quality of the global description is better for the $`\nu _\mu \nu _\tau `$ channel although oscillations into $`\nu _\mu \nu _s`$ cannot be statistically ruled out on the basis of the global fit to the full set of observables. In Fig. 8 we show the zenith-angle distributions for the Super–Kamiokande data sets, calculated for the best–fit points obtained in the global analysis of the data, when only “positive results” experiments are considered. From the figure, we notice that the $`\nu _\mu \nu _s`$ cases have more difficulties in reproducing the distribution of the data points because for $`\nu _\mu \nu _s`$ the survival probability of muon–neutrinos has a less steep angular behaviour as compared with the $`\nu _\mu \nu _\tau `$ case due to the Earth matter effects present in the $`\nu _\mu \nu _s`$ channels. ## IV Discussion In this paper we have performed a global analysis of the atmospheric neutrino data in terms of neutrino oscillations. We have compared the relative statistical relevance of the active-active and active-sterile channels as potential explanations of the atmospheric neutrino anomaly. In the analysis we have included, for the contained events, the latest data from Super–Kamiokande, corresponding to 52 kton-yr, together with all other available atmospheric neutrino data in the sub-GeV and multi-GeV range. Specifically, we included the data sets of the Frejus, Nusex, IMB, Soudan2 and Kamiokande experiments. Our analysis also includes the results on neutrino induced upgoing muons from the Super–Kamiokande 52 kton-yr sample, from MACRO and from Baksan experiments. We have determined the best-fit neutrino oscillation parameters and the resulting allowed regions in $`\mathrm{sin}^2(2\theta )`$ and $`\mathrm{\Delta }m^2`$ for $`\nu _\mu \nu _X`$ ($`X=\tau ,s`$) channels. For oscillations into sterile neutrinos we have considered both positive and negative $`\mathrm{\Delta }m^2`$, since the two cases differ in the matter effect for neutrinos propagation in the Earth. The results of the combined analysis indicate that the channel $`\nu _\mu \nu _\tau `$ gives a better fit to the data than oscillations into sterile neutrinos: $`\chi _{min}^2=`$ 74/(73 d.o.f.) (45 % CL) for the active case versus $`\chi _{min}^2=`$ 90/(73 d.o.f.) (8.5 % CL) for oscillations into sterile neutrinos with $`\mathrm{\Delta }m^2<0`$ and $`\chi _{min}^2=`$ 86/(73 d.o.f.) (14 % CL) for $`\nu _\mu \nu _s`$ with $`\mathrm{\Delta }m^2>0`$. Since for some experiments we obtain large $`\chi _{min}^2`$ values also for neutrino oscillations, values not clearly better than the analysis for the standard model case, we decided to perform additional analyses by removing these “negative-result experiments”. When excluding Frejus, Nusex and Baksan data points, we do not obtain an improvement for the sterile cases, while for the active case the CL is increased to 58 % ($`\chi _{min}^2=`$ 58/(61 d.o.f.)). When also MACRO is removed, we obtain higher CL also for the sterile cases, but the $`\nu _\mu \nu _\tau `$ hypothesis remains as the best option: $`\chi _{min}^2=`$ 41/(51 d.o.f.) (84 % CL) for the active case, $`\chi _{min}^2=`$ 51/(51 d.o.f.) (47 % CL) for $`\nu _\mu \nu _s`$ with $`\mathrm{\Delta }m^2<0`$ and $`\chi _{min}^2=`$ 50/(51 d.o.f.) (51 % CL) for $`\nu _\mu \nu _s`$ with $`\mathrm{\Delta }m^2>0`$. We thus conclude that the quality of the global description of the atmospheric neutrino data in terms of neutrino oscillation is better for the $`\nu _\mu \nu _\tau `$ channel, although oscillations into $`\nu _\mu \nu _s`$ cannot be statistically ruled out, on the basis of the global fit to the full set of observables. We have also presented a sample of predicted zenith-angle distributions for the best-fit points corresponding to the various oscillation channels. Specifically, we showed the angular distribution corresponding to the four Super–Kamiokande data sets (sub-GeV, multi-GeV, stopping muons and through-going muons) and the angular distributions for upgoing muons at MACRO and Baksan. By using the zenith-angle distribution expected for contained events and up-going muon data at Super-Kamiokande in the presence of oscillations, we have compared the relative goodness of the three possible oscillation channels. This allowed one to understand why the $`\nu _\mu \nu _\tau `$ channel gives a better description than the $`\nu _\mu \nu _s`$ cases. On the one hand, due to matter effects, one finds that for the sterile case the up–down asymmetry in the multi-GeV sample is slightly smaller than observed. Moreover, also due to matter effects, the upgoing-muon distributions in the case of $`\nu _\mu \nu _s`$ are flatter than for $`\nu _\mu \nu _\tau `$, while the data show a slightly steeper angular dependence which can be better described by $`\nu _\mu \nu _\tau `$. To conclude, we compare our determinations of the allowed neutrino oscillation parameters from the analyses of the atmospheric neutrino data with the expected sensitivities of future long-baseline experiments such as K2K and MINOS (Fig. 7). We notice that, for the oscillations $`\nu _\mu \nu _\tau `$, K2K can cover most of the 90% CL allowed region while the MINOS test of NC/CC is sensitive to the complete 99% CL region of oscillation parameters. For oscillations into sterile neutrinos, a situation where only a CC disappearance test can be performed at long baseline experiments, K2K can cover a portion of the region allowed at 90% CL, while the full 99% CL allowed regions are completely accessible to the MINOS experiment. ## Appendix: Statistical Analysis In this Appendix, we discuss the errors and correlations employed in our analysis. ### A Errors Data errors contain the experimental statistical and systematic errors as well as the uncertainties arising from event mis-identification, as quoted by the experimental Collaborations. In the theoretical calculations of event rates and upgoing muon fluxes, we take into account the uncertainty in the atmospheric neutrino flux and the uncertainties in the charged-current neutrino cross-sections . We also include, the Monte Carlo (MC) statistical errors estimated by the experimental Collaborations with the simulations of their detectors. This uncertainty depends on the number of simulated Monte–Carlo events. The flux uncertainty is taken to be $`30\%`$ at lower energies, relevant for contained events, and $`20\%`$ at higher energies, characteristic of upward going muons. Nuclear cross-section uncertainties are taken to be 10% for all the contained event samples, except for Soudan2 for which we used the values $`7.5\%`$ and $`6.4\%`$ for $`e`$-like and $`\mu `$-like events, respectively . For stopping and through-going muons, as cross-section uncertainties we use 11.4% and 14.1 %, respectively . The Monte–Carlo statistical errors are estimated from the simulated exposure, as given by the experimental Collaborations, under the assumption that the $`e`$\- and $`\mu `$-like contained events follow a binomial distribution. ### B Correlations There is a large number of correlations, both from experimental and from theoretical sources. In our analysis, data errors between different experiments are assumed to be uncorrelated, i.e. $`\rho _{\alpha \alpha }^{DATA}(A,A)=1`$ $`(\alpha =e,\mu )`$ for all data samples (18) $`\rho _{\alpha \beta }^{DATA}(A,B)=0`$ $`(\alpha ,\beta =e,\mu )`$ $`\text{if}AB,`$ (19) while the correlations in the theoretical quantities referring to different experiments (i.e., for $`AB`$) are assumed to arise from the correlations amongst the theoretical errors for the neutrino flux normalization and amongst the theoretical uncertainties of the neutrino interaction cross section. In order to determine the different entries, we have classified the experiments in four samples: \| $``$ \| sub-GeV (SG): \| Frejus, IMB, Nusex, Soudan2, Kam sub-GeV \| \| --- \| --- \| --- \| \| \| \| and Super–Kam sub-GeV \| \| $``$ \| multi-GeV (MG): \| Kam multi-GeV and Super–Kam multi-GeV \| \| $``$ \| stopping muons (STOP): \| Super–Kam \| \| $``$ \| thru-going muons (THRU): \| Super–Kam, MACRO and Soudan. \| We then estimate the correlations as follows, $$\rho _{\alpha \beta }^{TH}(A,B)=\rho _{\alpha \beta }^{flux}(A,B)\times \frac{\sigma _\alpha ^{flux}(A)\sigma _\beta ^{flux}(B)}{\sigma _\alpha ^{TH}(A)\sigma _\beta ^{TH}(B)}+\rho _{\alpha \beta }^{cross}(A,B)\times \frac{\sigma _\alpha ^{cross}(A)\sigma _\beta ^{cross}(B)}{\sigma _\alpha ^{TH}(A)\sigma _\beta ^{TH}(B)}$$ where the correlation $`\rho _{\alpha \beta }^{flux}`$ and $`\rho _{\alpha \beta }^{cross}`$ can be computed from the allowed uncertainties on the corresponding ratios by means of: $`\sigma _{Rf}^2(A:\alpha ,B:\beta )`$ $`\sigma ^2\left({\displaystyle \frac{\mathrm{\Phi }_\alpha (A)}{\mathrm{\Phi }_\beta (B)}}\right)`$ $`=\sigma _\alpha ^{2flux}(A)+\sigma _\beta ^{2flux}(B)2\rho _{\alpha \beta }^{flux}(A,B)\sigma _\alpha ^{flux}(A)\sigma _\beta ^{flux}(B)`$ (20) $`\sigma _{Rc}^2(A:\alpha ,B:\beta )`$ $`\sigma ^2\left({\displaystyle \frac{\sigma _\alpha ^{CC}(A)}{\sigma _\beta ^{CC}(B)}}\right)`$ $`=\sigma _\alpha ^{2cross}(A)+\sigma _\beta ^{2cross}(B)2\rho _{\alpha \beta }^{cross}(A,B)\sigma _\alpha ^{cross}(A)\sigma _\beta ^{cross}(B)`$ (21) We assume that the errors in the above defined ratios of the neutrino fluxes and the neutrino cross-sections between the different experiments arise from three sources: the flavour dependence, the energy dependence and the angular dependence . $`\sigma _{Rf}^2(A:\alpha ,B:\beta )`$ $`=`$ $`\sigma _{Rf,flav}^2(A:\alpha ,B:\beta )+\sigma _{Rf,en}^2+(A:\alpha ,B:\beta )+\sigma _{ang}^2(A:\alpha ,B:\beta )`$ (22) $`\sigma _{Rc}^2(A:\alpha ,B:\beta )`$ $`=`$ $`\sigma _{Rc,flav}^2(A:\alpha ,B:\beta )+\sigma _{Rc,en}^2(A:\alpha ,B:\beta )+\sigma _{ang}^2(A:\alpha ,B:\beta )`$ (23) For sub-GeV data samples, there are additional sources of uncertainties with respect to the other data sets. For instance, the form of the geo–magnetic cut–off used in the neutrino flux, the nuclear modeling for the neutrino interaction cross section, the neutrino production in the atmosphere. These additional uncertainties allow a large variation of the expectations for sub-GeV experiments, without affecting the predictions for higher energy events. Thus we neglect the correlations in the theoretical errors between sub-GeV events and any of the higher energy samples. Moreover the different experimental collaborations use different calculations for the neutrino interaction cross section, thus we neglect correlations between the rΦtheoretical errors of the interaction cross sections for different detectors. For all other cases we use the following values: \| $`\sigma _{Rf,flav}^2(A:\alpha ,B:\beta )=0`$ \| for $`\alpha =\beta `$ \| \| --- \| --- \| \| $`\sigma _{Rf,flav}^2(A:e,B:\mu )=5\%`$ \| for A, B unbinned (SG) \| \| $`\sigma _{Rf,flav}^2(A:e,B:\mu )=10\%`$ \| for A, B binned (SG) or (MG) \| \| $`\sigma _{Rf,flav}^2(A:\mu ,B:\mu )=10\%`$ \| for A in (MG) and B in (STOP) \| \| $`\sigma _{Rf,flav}^2(A:\mu ,B:\mu )=14\%`$ \| for A in (MG) and B in (THRU) \| \| $`\sigma _{Rf,en}^2(A:\alpha ,B:\beta )=0`$ \| for A, B in the same sample \| \| $`\sigma _{Rf,en}^2(A:\mu ,B:\mu )=5\%`$ \| for A in (MG) and B in (STOP) \| \| $`\sigma _{Rf,en}^2(A:\mu ,B:\mu )=10\%`$ \| for A in (STOP) and B in (THRU) \| \| $`\sigma _{Rf,en}^2(A:\mu ,B:\mu )=11\%`$ \| for A in (MG) and B in (THRU) \| \| $`\sigma _{Rc,flav}^2(A:\alpha ,B:\beta )=0`$ \| for $`\alpha =\beta `$ \| \| $`\sigma _{Rc,flav}^2(A:e,B:\mu )=`$ 3 to 10% \| for A, B (SG) \| \| $`\sigma _{Rc,fl}^2(A:e,B:\mu )=`$ 2 to 4% \| for A, B in (MG) \| \| $`\sigma _{Rc,en}^2(A:\alpha ,B:\beta )=0`$ \| for A, B in the same sample \| \| $`\sigma _{Rc,en}^2(A:\alpha ,B:\mu )=5\%`$ \| for A in (MG) and B in (STOP) \| \| $`\sigma _{Rc,en}^2(A:\alpha ,B:\mu )=7\%`$ \| for A in (MG) and B in (THRU) \| \| $`\sigma _{Rc,en}^2(A:\mu ,B:\mu )=5\%`$ \| for A in (STOP) and B in (THRU) \| \| $`\sigma _{ang}^2(A:\alpha ,B:\beta )=5\%\|\mathrm{cos}(\theta _A)\mathrm{cos}(\theta _B)\|`$ \| for all angular distributions \| With this we get that the smallest non-vanishing theoretical correlation within the Super–Kamiokande samples occurs between the theoretical errors for most vertical multi-GeV electron bin and the most horizontal thru-going muon bin and it takes the value $`\rho ^{TH}=0.735`$ while, for example, the correlation between the theoretical errors of the two most horizontal bins of the thru-going muon sample is $`\rho ^{TH}=0.989`$. ###### Acknowledgements. We thank Ricardo Vazquez for providing us with the codes for the evaluation of the muon energy loss, and Todor Stanev who provided us with his atmospheric neutrino fluxes. It is a pleasure also to thank Mark Messier, Teresa Montaruli and Olga Suvorova for useful discussions. This work was supported by DGICYT under grants PB98-0693 and PB97-1261, and by the TMR network grant ERBFMRXCT960090 of the European Union.
warning/0002/physics0002043.html	ar5iv	text	# The decay of multiscale signals – deterministic model of the Burgers turbulence ## 1 Introduction The nonlinear diffusion equation $$\frac{v}{t}+v\frac{v}{x}=\nu \frac{^2v}{x^2};v(x,t=0)=v_0(x).$$ (1.1) was originally introduced by J.M.Burgers in (1939) as a model for hydrodynamical turbulence. Burgers’ equation (1.1) describes two fundamental effects characteristic of any turbulence : the nonlinear redistribution of energy over the spectrum and the action of viscosity in small scales. Burgers’ equation was used later to describe a large class of physical systems in which the nonlinearity is fairly weak (quadratic) and the dispersion is negligible compared to the linear damping . The most important example of such waves are acoustical waves with finite amplitude . Another class of problems, arising, e.g., in surface growth, also leads to Burgers’ equation ,,. The three dimensional form of (1.1) has been used in cosmology to describe the formation of large scale structures of the Universe at a nonlinear stage of gravitational instability ( see e.g. ). In the physically important case of large Reynolds number, the action of viscosity is significant only in the small regions with high gradient of the velocity field. In the limit $`\nu 0`$, the solution of Burgers’ equation has the following form (see , , ): $$v(x,t)=\frac{xy(x,t)}{t},$$ (1.2) where $`y(x,t)`$ is the coordinate of the maximum of the function $$G(x,y,t)=\mathrm{\Psi }_0(y)\frac{(xy)^2}{2t},v_0(x)=\frac{\mathrm{\Psi }_0(x)}{x}.$$ (1.3) Strong interaction between coherent harmonics leads to the appearance of local self-similar structures in Burgers’ equation. A periodic initial perturbation with zero mean velocity is transformed asymptotically into a sawtooth wave with gradient $`_xv=1/t`$ and with the same period $`l_0`$. It is important that at this stage the amplitude $`a(t)=l_0/t`$ and the energy density $`\sigma ^2(t)l_0^2/12t^2`$ do not depend on the initial amplitude. An initial one-signed pulse with the area $`m>0`$, localised at $`t=0`$ in the neighbourhood of the point $`x=0`$, also has asymptotically a universal form: it transforms into a triangular pulse with the gradient $`_xv=1/t`$ and increasing coordinate of the shock $`x_s(2mt)^{1/2}`$. Due to the increase of the integral scale the amplitude of such a pulse $`a(t)=x_s(t)/tm^{1/2}t^{1/2}`$ and its energy will decrease more slowly than for a periodic signal, like $`t^{1/2}`$. Continuous random initial fields are also transformed into sequences of regions with the same gradient $`_xv=1/t`$, but with random locations of the shocks separating them. Due to the multiple merging of the shocks the statistical properties of such random fields are also self-similar and may be characterised by the integral scale of the turbulence $`L(t)`$. The merging of the shocks leads to an increase of the integral scale $`L(t)`$, and because of this the energy $$\sigma ^2(t)L^2(t)/t^2$$ (1.4) of a random wave decreases more slowly than the energy of periodic signals. The type of turbulence evolution is determined by the behaviour of the large scale part of the initial energy spectrum $$E_0(k)=\alpha ^2k^nb_0(k);$$ (1.5) $$E_0(k)=\frac{1}{2\pi }v_0(x),v_0(x+z)e^{ikz}𝑑z.$$ (1.6) Here $`b_0(k)`$ is a function which falls off rapidly for $`k>k_0l_0`$, and $`b_0(0)=1`$. For $`n>1`$ the law of enrgy decay strongly depends on the statistical properties of the initial field (see e.g. and references therein). For the initial Gaussian perturbation the integral scale $`L(t)t^{1/2}`$ times logarithmic correction obtains and is determined by two integral characteristics of the initial spectrum: the variances of the initial potential $`\mathrm{\Psi }_0`$ and the velocity $`v_0(x)`$ ,,,. For $`n<1`$ the structure function of the initial potential increases as a power law in space. Then the initial potential field is Brownian, or fractional Brownian motion, and some scaling may be used ,, , , ,, . In this case the turbulence is also self-similar and the integral scale $`L(t)`$ increases as $$L(t)=(\alpha t)^{2/(3+n)}.$$ (1.7) The energy of the turbulence is derived from (1.4): $$\sigma ^2(t)t^p,p=\frac{2(n+1)}{n+3}.$$ (1.8) The difference between these two cases ( $`n<1`$ and $`n>1`$ ) is connected to the process of parametric generation of low frequency component of the spectrum. For the case $`n<1`$ the newly generated low frequency components are relatively small and we have the conservation of large scale part of the spectrum: $$E(k,t)=E_0(k)=\alpha ^2k^n,fork<<1/L(t).$$ (1.9) Thus, the laws of turbulent decay are more complex than for simple signals, which can be attributed to multiple merging of the shocks. In a model of a regular fractal signal with decay lower than for single one-signed pulse was introduced. The initial signal $`v_0(x)`$ was constructed as a sequence of one-signed pulses whose positions form a Cantor set with capacity (fractal dimension) $`D=\mathrm{ln}N/\mathrm{ln}\beta `$, where $`N^p`$ is the number of pulses in the scale $`L_pL_1N\beta ^{p1}`$$`0<D<1`$. Multiple merging makes the decay of the wave slower and the general behaviour of the energy decay may be approximated by the power law with the exponent in (1.8): $$p=\frac{1D}{2D},0<p<1.$$ (1.10) The evolution proves to be self-similar in successive time periods $`(t_i,t_{i1})`$ and $`(t_{i+1},t_{i+2})`$, where $`t_{i+1}/t_i=\beta ^2/N`$. This shows log-periodical self-similarity of the field evolution. Linear and non-linear decay of fractal and spiral fields given by the sequences of regular pulses was also investigated in . It was shown that the power law (1.8) with the exponent given by the formula (1.10) holds also true for homogeneous fractal pulse signal with capacity $`D`$. Another model of a multiscale signal, which has the same general behaviour on the external scale $`L(t)`$ (1.7), and energy of the Burgers turbulence, was also discussed in . It was assumed therein that the initial signal is a discrete set of modes - the spatial harmonics $$v_0(x)=\underset{p=0}{\overset{\mathrm{}}{}}a_p\mathrm{sin}(k_px+\phi _p),$$ (1.11) with wavenumbers $`k_p`$ and amplitudes $`a_p`$ given in terms of a parameter $`ϵ`$ by $`k_p`$ $`=`$ $`k_0ϵ^p,a_p=a_0ϵ^{hp},`$ $`h`$ $`=`$ $`({\displaystyle \frac{n+1}{2}}),a_0=\alpha k_0^{(n+1)/2}.`$ (1.12) Amplitudes $`a_p`$ and the scaling exponent $`h`$ are chosen from the condition that the mean energy of harmonics in the interval $`\mathrm{\Delta }_p=k_pk_{p+1}`$ be identical to that corresponding to the spectral density (1.5): $`a_p^2=E(k_p)\mathrm{\Delta }_p`$. For $`ϵ<<1`$ and $`n>1`$ the harmonics are spread over the spatial spectrum and accumulate at the point $`k=0`$ with decreasing amplitude. The main approach in this model was that the energy of the wave is the sum of energies of independent modes. The approach is nontrivial, but nevertheless leads to the same laws of the integral scale $`L(t)`$ (1.7) and the energy decay (1.8) as in the case of continuous spectrum (1.5). Let us point out that representation of the field given by the formula (1.11) is similar to shell models, which were introduced as useful models addressing the problem of analogous scaling in fully developed turbulence (see, e.g., , and references in there). In present paper we consider the evolution of a regular signal whose behaviour in general is similar to the evolution of the Burgers turbulence with continuous spectrum (1.5). The main difference,as compared to the model discussed above, is that we construct the exact solution of Burgers’ equation using as an initial mode the ”reverse” sawtooth wave. The frequency ratio in our model is $`ϵ=1/N`$, where $`N`$ is an integer and $`N2`$. These perturbation are similar to the well known Weierstrass and Weierstrass-Mandelbrot fractal functions (see , ). For the analysis of Burgers’ equation it is convenient to use a mechanical interpretation. There is a one-to-one correspondence between the solution of Burgers’ equation and the dynamic of a gas of inelastically interacting particles , . Let us take a one-dimensional particle flux with a contact interaction: as long as the particles do not run into each other they move with constant velocity. In the collision they stick together, forming a delta-function singularity in the matter density. This leads to the appearance of gas of two species: a hydrodynamical flux of the ”light” initial particles, and a gas of ”heavy” particles arising in the adhesion process of light particles. The evolution of the particle velocity field will be described by the solution of Burgers’ equation if we assume that the initial density of the light particles is $`\rho _0=const`$, the velocity of particles is equal to the initial velocity in (1.1), and that the collision of the particles conserve their mass and momentum. This analogy permits construction of a very fast (linear time) algorithm of solution of Burgers’ equation . In our case, for the initial reverse sawtooth wave, all the matter turns into heavy particles at the same moment of time. Thus, after this time the evolution of the Burgers turbulence is fully determined by the motion of heavy particles, whose positions are positions of shocks, and masses are equal to $`\mathrm{\Delta }vt`$, where $`\mathrm{\Delta }v`$ are the amplitudes of the shocks. The paper is organized as follows. In Section 2 we consider the evolution and interaction of “reverse” sawtooth modes. In section 3 we consider the interaction of small scale mode with large scale structures. In Section 4 we investigate the properties of the sawtooth Weierstrass-Mandelbrot fractal function. In section 5 we show that deterministic model has logarithmic periodic self-similarity. We also discuss here the multi-dimensional generalization of this model. Section 5 presents concluding remarks. ## 2 Evolution and interaction of ”reverse” sawtooth modes in Burgers’ equation Let us introduce the $`p`$th ”reverse” saw tooth mode as $$v_{(p)}(x,0)=a_pA(k_px+\phi _p)$$ (2.1) Here $`a_p,k_p`$ are amplitude and wavenumber of the mode, $`\phi _p`$ is its phase. The function $`A(x)`$ is $`2\pi `$ periodic function, given on its first period by the following expression $$A(x)=\pi x,x[0,2\pi [$$ (2.2) The set of ”reverse sawtooth” functions is not orthogonal, but nevertheless we will introduce a set of modes satisfying equation (1.12), whose wavenumbers and amplitudes satisfy the same relations as the sinusoidal modes of , i.e. relation (1.11). We introduce the term ”reverse sawtooth” because this signal is a sawtooth with teeth facing to the right, but the term ”sawtooth” by itself is widely used in Burgers turbulence literature to refer to the late stage of the evolution of the wave profile a sequence of sawteeth with positive slope $`1/t`$. The solution of Burgers’ equation with a linear velocity profile $`v_0=\gamma (xx_+)`$ is well-known (see, e.g., ): $$v(x,t)=\frac{\gamma (xx_+)}{1\gamma t}$$ (2.3) The value $`\gamma ^1`$ has the dimension of time, and for $`\gamma >0`$, at the finite time $`t=\gamma ^1`$ the gradient $`_xv`$ becames infinite. For $`\gamma <0`$ the gradient becomes equal to $`_xv=t^1`$, independent of the value of $`\gamma `$ at times $`t\|\gamma \|^1`$. Thus, we have from (1.12), (2.1),(2.2), that the evolution of the $`p`$th mode is characterised by the nonlinear time $$t_p=\gamma _p^1=a_pk_{p}^{}{}_{}{}^{1}=t_0/(ϵ^{(n+3)/2})^p,t_0=1/\alpha k_0^{n+3/2}$$ (2.4) Based on solution (2.3 ) it is easy to see that for the ”first” period (if $`\phi _p=0`$) the evolution of the $`p`$th reverse mode at the initial stage $`(t<t_p)`$ may be described as $$v_{(p)}(x,t)=\{\begin{array}{cc}x/t,\hfill & \text{if }0<x<\frac{t}{t_p}\frac{\pi }{k_p}\text{;}\hfill \\ \frac{1}{t_p}(\frac{\pi }{k_p}x)/\frac{1}{1t/t_p},\hfill & \text{if }\|\frac{\pi }{k_p}x\|<\frac{t}{t_p}\frac{\pi }{k_p}\text{ ;}\hfill \\ (x\frac{\pi }{k_p})/t,\hfill & \text{if }\frac{t}{t_p}\frac{\pi }{k_p}<x\frac{\pi }{k_p}<\frac{\pi }{k_p}\text{ .}\hfill \end{array}$$ (2.5) On the other hand, at time $`t=t_p`$, the mode transforms into a ”direct” sawtooth wave with slope $`_xv=1/t`$ , independent of the amplitude and wavenumber of the mode: $$v_{(p)}(x,t)=\{\begin{array}{cc}x/t,\hfill & \text{if }0<x<\frac{\pi }{k_p}\text{;}\hfill \\ (x\frac{\pi }{k_p})/t,\hfill & \text{if }\frac{\pi }{k_p}<x<\frac{2\pi }{k_p}\text{ .}\hfill \end{array}$$ (2.6) The density of energy $`\sigma ^2(t)=v^2(x,t)_L`$, where $`_L`$ denotes averaging over the period, is conserved before $`t<t_p`$, and decreases like $`(k_pt)^2`$ after $`t>t_p`$. Consider now the evolution of a gas of sticky particles in the case of independent evolution of the $`p`$th mode. In the general case, the density of the gas is calculated by using the Jacobian of the transformation from Lagrangian to Eulerian coordinates and may be written in the form ( see, e.g., ) $$\rho (x,t)=\rho _0(1t_xv(x,t))$$ (2.7) Then it is obvious that at the initial stage $$\rho (x,t)=\rho _0\frac{1}{1t/t_p},\|\frac{\pi }{k_p}x\|<\frac{t}{t_p}\frac{\pi }{k_p}$$ (2.8) while $`\rho `$ is zero outside this interval in each period. At time $`t=t_p`$ all the light particles in each period collide into a single heavy particle with mass $$m_p=\rho _0L_p=\rho _0\frac{2\pi }{k_p},$$ (2.9) and the heavy particles have positions $$x_{p,l}=\frac{\pi }{k_p}\frac{\phi _p}{k_p}+\frac{2\pi }{k_p}l;l=0,\pm 1,\pm 2,\mathrm{}$$ (2.10) equal to the zero positions of the initial $`p`$th mode. The process of light merging particles and the evolution of the velocity is shown in Fig.1. Consider now the joint evolution of two successive modes: $`p`$th and $`(p+1)`$th. From (2.4) one can see that the ratio of nonlinear times of the successive modes is $$\frac{t_{p+1}}{t_p}=ϵ^{(n+3)/2}ϵ^{1h},$$ (2.11) which does not depend on $`p`$ and increases if the exponent $`n`$ is greater than $`3`$. The gradient of the initial field $`v_{(p)}(x)+v_{(p+1)}(x)`$ is $`(\gamma _p+\gamma _{p+1})`$, so the effective nonlinear time for such a sum is $$t_{p,eff}=t_{p,p+1}=\frac{1}{\gamma _p+\gamma _{p+1}}=\frac{t_p}{1+t_p/t_{p+1}},$$ (2.12) Because all parts of the initial perturbation have the same slope, all light particles will collide at the same time $`t=t_{p,p+1}`$. The mass of heavy particles after merging are $`m_{p,i}=\rho _0\mathrm{\Delta }_{i,i+1}`$, where $`\mathrm{\Delta }_{i,i+1}`$ is the distance between adjacent shocks in the initial perturbation. The ratio of periods of two adjacent modes is $`L_{p+1}/L_p=k_p/k_{p+1}=ϵ^1`$. If $`N=ϵ^1`$ is an integer larger than $`1`$, there will be $`N+1`$ heavy particles on the period of the larger scale mode $`(p+1)`$th: $`(N1)`$ with the mass $`m_p=\rho _0L_p`$ (2.9), and two particles with total mass equal to $`m_p`$. These two particles only exist when the shock of the $`(p+1)`$th mode is located in the interval between shocks of the $`p`$th mode. For simplicity, we will consider the case where the spatial relations $$k_{p+1}\phi _{p+1}=k_p\phi _p+2\pi r/N,$$ (2.13) between the phases of successive modes, hold. In this case the discontinuities of the $`(p+1)`$th mode do not produce new shocks in the total perturbation $`v_{(p)}(x)+v_{(p+1)}(x)`$. Thus, at times $`t`$ larger than $`t>t_{r,eff}`$, the masses of all heavy particles will be the same, as would be the case without the large scale modes (2.9). The positions of these heavy particles at time $`t=t_{p,eff}`$ are $$X_{(p,l)}(t_{p,eff})=x_{p,l}+v_{(p+1)}(x_{p,l})t_{p,eff},$$ (2.14) where $`x_{p,l}`$ are the zero positions of the $`p`$th mode (2.10). The velocity of this particle is equal to $`v_{(p+1)}(x_{p,l})`$. Equation (2.14) is obvious if we use the trivial equality $`v_{(p)}(x_{p,l})+v_{(p+1)}(x_{p,l})=v_{(p+1)}(x_{p,l})`$, and also note that the position of the heavy particle $`x_{p,l}(t_{p,eff})`$ is equal to the position at the same time of all light particles with initial coordinate $`x=x_{p,l}`$. From (2.14) we immediately have that after time $`t_{p,eff}`$ the positions of the particles are $$X_{(p,l)}(t)=x_{(p,l)}+v_{(p+1)}(x_{p,l})t.$$ (2.15) The difference between the coordinates of the adjacent particles $`X_{p,l}(t)`$ and $`X_{p,l+1}(t)`$ decreases with time, proportionally to the gradient of $`v_{(p+1)}(x)`$: $`X_{(p,l+1)}(t)X_{p,l}(t)=(x_{p,l+1}x_{p,l})t{\displaystyle \frac{v_{(p+1)}(x)}{x}}(x_{p,l+1}x_{p,l})`$ $`(x_{p,l+1}x_{p,l})(1t/t_{p+1}).`$ (2.16) These particles collide at time $`t=t_{p+1}`$ (2.4) and the newly created heavy particles will have masses $$m_{(p+1)}=\rho _0L_{p+1}$$ (2.17) and positions $$x_{p+1,l}=\frac{\pi }{k_{p+1}}\frac{\phi _{p+1}}{k_p}+\frac{2\pi l}{k_{p+1}};l=0,\pm 1,\pm 2,\mathrm{}$$ (2.18) The velocity of these particles is zero. Thus, at times $`t`$ larger then $`t_{p+1}`$, the evolution of the initial perturbation $`v_0(x)=v_{(p)}(x)+v_{(p+1)}(x)`$ will be the same as the evolution of only the large scale mode $`v_{(p+1)}(x)`$. The process of particles merging and the evolution of the velocity for the sum of two successive modes with the periods ratio $`N=2`$ are shown in Fig. 2. By recurrence, it is evident that for finite number of modes $`v(x)=v_{(p)}(x)+v_{(p+1)}(x)+\mathrm{}+v_{(M)}(x)`$ the evolution of the field after $`t_M`$ will be the same as the evolution of only the largest mode $`v_{(M)}(x)`$. The reason for this, is of course, the special relation between the phases $`\phi _p`$ and the wavenumbers $`k_p`$ of all interacting modes: $`k_p=k_0/N^p`$ ( see equation (1.12) ). For integer $`N`$ the minimal value of any combination of these wavenumbers is equal to the largest mode wavenumber $`k_M=k_0/N^M`$. So the nonlinear interaction does not produce new components at frequencies less than $`k_M`$. ## 3 Interaction of small scale ”reverse” sawtooth mode with large scale structures Let us now consider the interaction of the $`p`$th mode with an infinite series of larger scale modes $$W_p(x)=\underset{r=p+1}{\overset{\mathrm{}}{}}v_r(x)\underset{r=p+1}{\overset{\mathrm{}}{}}a_rA(k_rx+\phi _r),$$ (3.1) assuming that the phases of the modes satisfy the relations ( 2.13 ) and that $`k_r=k_0ϵ^r,a_r=a_0ϵ^{hr},h=(n+1)/2`$. From (3.1) and (2.4) we have for the gradient of the initial perturbation $`v_0(x)=v_p(x)+W_p(x)`$ $`_xv_0(x)`$ $`=_xv_p(x)+_xW_p(x)={\displaystyle \underset{r=p}{\overset{\mathrm{}}{}}}\gamma _p=`$ (3.2) $`=\gamma _0{\displaystyle \underset{r=p}{\overset{\mathrm{}}{}}}(ϵ^{\frac{(n+3)}{2}})^r=\gamma _0ϵ^{\frac{(n+3)}{2}p}{\displaystyle \frac{1}{1ϵ^{\frac{(n+3)}{2}}}},`$ the condition $`n>3`$ ($`h<1`$) being necessary for the series to converge. From (3.2), we have for the effective time of nonlinearity of $`p`$th mode $$\stackrel{~}{t}_p=1/_xv_0(x)=t_p(1ϵ^{\frac{(n+3)}{2}}),$$ (3.3) with the original $`t_p`$ determined by the equation (2.4). Thus, after the time of collision $`\stackrel{~}{t}_p`$, heavy particles with mass $`m_p=\rho _0L_p`$ (3.2) appear. The coordinates of these particles will be determined by an equation similar to (2.15) $$x_{p,l}(t)=x_{p,l}(t)+W_p(x_{p,l})t,$$ (3.4) with the velocity of particles determined by the function $`W_p(x)`$ (see 3.1), which is a sum of all larger modes, and $`x_{p,l}`$ are the coordinates of the zeros of the $`p`$th modes. The difference between the coordinates of adjacent particles $`x_{p,l}`$ and $`x_{p,l}`$ will then decrease with time like $`(1t/\stackrel{~}{t}_{p+1})`$, where $`\stackrel{~}{t}_{p+1}=t_{p+1}(1ϵ^{(n+3)/2})`$ is the inverse of the gradient of the function $`W_p(x)`$, see (3.1) and (3.2). Thus, the time of particle collision for this generation will be described by equation (3.3) with $`p=p+1`$, and the new masses will be determined by the period of the $`(p+1)`$th mode - see equation (2.17). The extrapolation of this particle merging process to the next generations is evident by recurrence. The $`q`$th collision of heavy particles takes place at time $`\stackrel{~}{t}_{p+q}=t_{p+q}(1ϵ^{\frac{n+3}{2}}`$), the masses of these particles at this time are determined by the period of the $`(p+q)`$th mode $`m_{p+q}=\rho _0L_{p+q}=2\pi \rho _0/k_{p+q}`$ (2.9), (2.17). In the time interval $`t[\stackrel{~}{t}_{p+q},\stackrel{~}{t}_{p+q+1}]`$, $$\frac{\stackrel{~}{t}_{p+q+1}}{\stackrel{~}{t}_{p+q}}=ϵ^{\frac{(n+3)}{2}}=N^{\frac{n+3}{2}};$$ (3.5) the coordinates of particles will be determined by the equation (3.4) with $`p=p+q`$. Here $`W_{p+q}(x)`$ is the sum of the velocities of all larger modes with $`r>p+q`$, and $`x_{p+q,l}`$ are the zeros of $`(p+q)`$-th mode. It is important to note that at time $`t>\stackrel{~}{t}_{p+q}`$ the evolution of the particles is solely determined by the modes with $`rp+q`$. It means, that at times $`t>\stackrel{~}{t}_{p+q}`$ the position of the particles does not depend on the presence in the initial condition of the small scale modes with $`r<p+q`$. Thus, two processes with different initial velocities: $`\stackrel{~}{v}_0(x)`$, the field with small scales, and $`v_0(x)`$, the field without small scales modes: $$v_0(x)=W_{p+q1}(x);\stackrel{~}{v}_0(x)=W_{p1}(x)$$ (3.6) will have the same evolution after $`t>\stackrel{~}{t}_{p+q1}`$. Even if $`p\mathrm{}`$ (when modes with very small scales $`L_pϵ^p=N^p`$ and very large amplitudes $`a_pa_0(ϵ^{(n+1)/2})^p=a_0(N^{(n+1)/2})^p)`$ are present in the initial perturbation) the multiple merging of the particles will lead to the independence of the evolution of large scale modes with respect to the small scale modes. This effect is similar to the self–preservation of large scale structures in Burgers turbulence ,. When the initial field $`v_0(x)`$ is noise, the highly nonlinear structures continuously interact and due to the merging of shocks, their characteristic scale $`L(t)`$ constantly increases. The presence of small scale noise perturbation $`v_h(x)`$ results in additional fluctuations in the shock coordinates $`\mathrm{\Delta }x_k(t)`$, and these fluctuations increase in strength with the passage of time. Thus, the final result of the evolution of the field is determined by the competition of two factors, the increase in the external scale $`L(t)`$ of the structures and the increase in the strength $`\mathrm{\Delta }x_k(t)`$ of shock coordinates fluctuations, the later being related to the perturbation $`v_h(x)`$. In a turbulence, having power index $`n<1`$ (1.5), multiple merging of shocks leads to self-preservation of the large scale structures independently of the presence of small scale components. For the model signal this effect appears for arbitrary $`n`$ due to the special choice of wavenumbers and phases of interacting modes. It was stressed in the introduction that the solution of Burgers’ equation has a one-to-one correspondence with the dynamics of the gas of inelastically interacting particles (). The stage when all light particles collide, forming heavy particles, corresponds to the solution of Burgers’ equation with a well-defined slope $`_xv=1/t`$. In this case the profile of the field $`v(x,t)`$ is fully determined by the coordinates and amplitudes of the shocks. Their coordinates $`X_s(t)`$ are equal to the coordinates of heavy particles, their velocity $$v_s(t)=\frac{dX_s(t)}{dt}=(v_s(x_s0,t)+v_s(x_s+0,t))/2$$ (3.7) is equal to the velocity of the particles, and the amplitude of the shock $$\mathrm{\Delta }v_s(x)=(v(x_s0,t)v(x_s+0,t))=m/t$$ (3.8) is determined by the mass of the particle $`(\rho _01)`$ (see, e.g., ). Thus, the investigation of the motion of heavy particles permits to fully reconstruct the properties of the velocity field $`v(x,t)`$ of Burgers’ equation. ## 4 The sawtooth Weierstrass-Mandelbrot fractal function It was shown in the previous section that the evolution of the particles (shocks) is determined by the function $`W_p(x)`$ (3.1). The basis functions of $`W_p(x)`$ are the reverse sawtooth periodic functions with wavenumbers $`k_r=k_0ϵ^r`$ and amplitudes $`a_r=a_0ϵ^{hr},h=(n+1)/2`$ satisfy relations (1.12). Wavenumbers form a geometrical progression like in the Weierstrass function (see ) and accumulate at the origin $`k=0`$. In the original Weierstrass function, the situation was the opposite with increasing frequencies, but nevertheless the function $`W_p(x)`$ has many properties of Weierstrass function and of its generalisation – the Weierstrass-Mandelbrot function (see , ). We consider here a deterministic function $`W_p(x)`$ with the special phase relation $`\phi _p=(2\pi k/N)p`$ $`(k=1,2,\mathrm{},N;N=1/ϵ`$ ), thus, the discontinuities in the largest modes $`r>p+1`$ coincide with some of the discontinuities of the smaller mode $`r=p+1`$. The function $`W_p(x)`$ is continuous in the intervals $`2\pi /k_{p+1}=2\pi /(k_0ϵ^{p+1})`$ with the same slope in each interval. The inverse value of this slope $`\stackrel{~}{t}_{p+1}`$ $$\stackrel{~}{t}_{p+1}=t_{p+1}(1ϵ^{\frac{n+3}{2}});t_{p+1}=t_0(ϵ^{\frac{n+3}{2}})^{p+1}$$ (4.1) is proportional to the nonlinear time $`t_{p+1}`$ of the smallest mode. Of course, we need $`n>3`$, so that the convergence of (3.2) is assured and the inequalities $`t_{p+1}>t_p`$ hold. The amplitudes of the modes are proportional to $`ϵ^{(n+1)/2}`$ and for $`n>1`$ the function $`W_p(x)`$ is bounded $$W_p(x)\underset{r=p+1}{\overset{\mathrm{}}{}}a_r=a_0(ϵ^{\frac{n+1}{2}})^{p+1}\frac{1}{1ϵ^{\frac{n+1}{2}}}.$$ (4.2) Thus, for finite $`p`$ the energy of $`W_p(x)`$ is also finite. For the case of the phase relation introduced above, the functions $`W_p(x)`$ also have scaling properties, so that for instance for $`k=0`$, we have $$W_p(x)=ϵ^{hp}W_0(ϵ^px);W_p(ϵ^mx)=ϵ^{hm}W_{m+p}(x).$$ (4.3) The case $`1<n<1`$ is similar to the initial conditions with generalized white noise in Burgers turbulence. The energy of the initial signal in such turbulence is determined by the largest cutoff wavenumber, so in our model by the smallest scale $`p`$. If $`p\mathrm{}`$ the energy of the model signal (as the energy of white noise) will tend to infinity. But from the considerations in the previous section we have that at the finite time $`t`$ all the modes with $`t_p<t`$ have finite energy $`L_p^2/t^2`$ due to the nonlinear dissipation, so that the whole energy of the turbulence is also finite. Thus, even in the case of ”divergent” initial conditions $`(p\mathrm{})`$, we will have a ”convergent” solution for any time $`t>0`$. The case $`n<1`$ is similar to having fractional Brownian motion initial condition in Burgers turbulence. In this case, the series (4.2) diverges and the initial signal $`W_p(x)`$ is unbounded. But for Burgers turbulence ( for the process of particles motion and collisions ) only relative velocity of the particles matters. So we can use the same regularisation procedure with $`W_p(x)`$. Such a procedure was done with the Weierstrass function in . In our case, taking $`\phi _p0`$, we can introduce the function $`W_p^{\mathrm{}}(x)=W_p(x)W_p(0)`$, according to , which is finite in all finite spatial intervals. The other way to get a bounded function is to use special phase relations for the modes. ## 5 Self-similarity properties of deterministic model in one and two dimension Here we summarise the properties of the evolution of the multiscale deterministic signal using some additional information about scaling characteristics of $`W_p(x)`$, and compare them with the properties of the Burgers turbulence. Let us consider the evolution of the multiscale signal $$v_0(x)=v_p(x)+W_p(x).$$ (5.1) It was shown that at times $`t`$ for which $`t>\stackrel{~}{t}_p`$, heavy particles with mass $`M_p=\rho _0L_p`$ appear and their coordinates are determined by the relation (3.5). These particles collide at time $`\stackrel{~}{t}_{p+1}`$, ($`\stackrel{~}{t}_{p+1}/\stackrel{~}{t}_p=N^{\frac{(n+3)}{2}},N=ϵ^1`$), and new particles with masses $`m_{p+1}=\rho _0L_{p+1}=m_pN`$ appear. Their motion will be determined by the same law (3.5) with substitution $`pp+1`$. Using the scaling properties of $`W_p(x)`$ (4.3) we have, that the motion of the particles in this interval will be similar to the motion of the particles in the interval $`[\stackrel{~}{t}_p,\stackrel{~}{t}_{p+1}]`$ if we rescale the time $`t/\stackrel{~}{t}_pt/\stackrel{~}{t}_{p+1}`$. Since the ratio $`t_{p+1}/t_p`$ does not depend on $`p`$, one can speak about the logarithmic periodic self-similarity of the motion of the particles. This means that at arbitrary interval $`[\stackrel{~}{t}_q,\stackrel{~}{t}_{q+1}]`$ the motion of the particles will be similar to the motion of the particles in the interval $`\stackrel{~}{t}_p,\stackrel{~}{t}_{p+1}`$, by the scaling factor $`x_p/x_q=ϵ^{pq}`$ in space, and the scaling function $`t_p/t_q=(ϵ^{\frac{n+3}{2}})^{pq}`$ in time. The coordinates and masses of the particles fully determine the velocity field, and so the solution of Burgers’ equation is also logarithmic periodic self–similar. With each collision, the mass $`M(t)`$ of the particles increases $`N=1/ϵ`$ times. The time interval between the two successive collisions increases as $`t_{p+1}/t_p=N^{\frac{n+3}{2}}`$. Thus, by the approximation of piecewise constant function $`m(t)`$ by the power law $$m(t)m_0(t/t_0)^{(n+3)/2}.$$ (5.2) we obtain the same result as for the Burgers turbulence. In our case, $`m`$ is proportional to the period of the smallest mode at time $`t`$, and is analogous to the integral scale in Burgers turbulence. In the case $`n>1`$ we can also estimate the energy decay of the model signal. For $`n>1`$ and $`ϵ1`$ the main energy of the signal at time $`t`$ is in the smallest mode and is proportional to $`L^2(t)/t^2`$. Thus, we have here again the same law for the energy decay as for Burgers turbulence. The numerical simulation based on the algorithm was done to illustrate the process of particles merging and velocity field evolution. The trajectories of the particles and profile of the field at different times are plotted for the initial ”white noise” signal ($`n=0,h=1/2`$) in Fig.3, and for the initial ”Brownian” motion ($`n=2,h=1/2`$) in Fig.4. Ten modes with the ratio of successive wavenumbers $`ϵ=1/N=1/2`$ were used. The plots show the initial stage of the evolution in some relatively small region where the finiteness of the number of modes is not significant. In Fig. 3 one can see that for $`n=0`$ the initial ”sawtooth” multiscale function oscillates near $`v=0`$ like a ”white noise” with finite variance. After the collision of light particles, when the reverse sawtooth function transforms into a sawtooth wave with positive gradient $`_xv=1/t`$, the structure of the signal is relatively simple, and even for $`N=2`$ the main energy remains in the mode with smallest wavenumber. In the case $`n=2`$ the initial profile has a large deviation behaviour which is typical for Brownian motion functions. After merging of light particles, the sawtooth profile has a set of small shocks with different amplitudes, which is also similar to the properties of Brownian signal in the Burgers turbulence . In Figs. 3(c) and 4(c) the velocity field at three successive merger times $`t_{}/t_{}=N^{(n+3)/2}`$ are plotted. These figures show the logarithmic periodic self-similarity of the evolution of multiscale signals. We notice now that such multiscale waves may be constructed for multidimensional Burgers’ equation. Let us assume that the initial vector field $`𝐕_p\left(𝐱\right)`$ is an infinite series of ”reverse” modes $`𝐯_r\left(𝐱\right)`$: $$𝐕_p\left(𝐱\right)=\underset{r=p}{\overset{\mathrm{}}{}}𝐯_r\left(𝐱\right),$$ (5.3) In the two dimensional case, the $`r`$th ”reverse” mode may be composed of piecewise linear functions defined on a system of regular triangles of size $`L_r`$ covering the plane. We consider here the special case when the ratio of the scales of two adjacent modes is $`L_{r+1}/L_r=ϵ^1=N=2`$. We assume also the special symmetry and phase relation between the different modes. In our case one big triangle is divided into four smaller triangles with vertices located at midpoints of its sides ( see Fig.5 ). We assume that inside each triangle in the $`r`$th mode the velocity has a linear profile $`𝐯_r\left(𝐱\right)=\gamma _r\left(𝐱𝐱_+\right)`$, where $`𝐱_+`$ is the coordinate of the center of the triangle. The solution of the multidimensional Burgers equation for such initial perturbation $$𝐯(𝐱,t)=\frac{\gamma \left(𝐱𝐱_+\right)}{1\gamma t}$$ (5.4) is now valid inside the triangle of size $`L_r(t)=L_r(1t\gamma _r)`$ . The value $`\gamma _r^1`$ has the dimension of time, and at the finite time $`t_r=\gamma _r^1`$ the velocity gradient becomes infinite. On the other hand, at time $`t=t_r`$, the mode transforms into a ”direct” sawtooth wave with the universal behaviour inside the new set of triangles and with the gradient $`1/t`$ independent of the amplitude and wavenumber of the mode: $$𝐯(𝐱,t)=\frac{𝐱𝐱_c}{t},$$ (5.5) where $`𝐱_c`$ is now the center of the triangle, coinciding with the top of the initial triangular set. Consider now the evolution of a gas of sticky particles in the case of independent evolution of the $`r`$th mode. Then, it is obvious that at the initial stage, inside the ”collapsing” triangle of size $`L_r(t)=L_r(1t\gamma _r)`$ the density increases as $$\rho (𝐱,t)=\rho _0\frac{1}{(1t/t_p)^2},$$ (5.6) while $`\rho `$ is zero outside ”collapsing” triangular in each cell. At time $`t=t_r`$, all the light particles in each cell collide into a single heavy particle with mass $$m_r=\rho _0L_r^2\sqrt{3}/4,$$ (5.7) and the heavy particles’ positions are equal to the center of the initial triangle $`𝐱_+`$. We assume also that the evolution of the $`r`$th mode is characterised by a nonlinear time $`t_r`$ the same as in the one dimension case (2.4): $$t_r=\gamma _r^1=t_0/(2^{(n+3)/2})^r.$$ (5.8) Let us now consider the evolution of the vector field $`𝐕_p\left(𝐱\right)`$ (5.3) which is an infinite series of ”reverse” modes. The evolution of the vector field is very similar to the evolution of the scalar field (3.1). For the gradient of the initial perturbation $`𝐕_p\left(𝐱\right)`$ we have the same relation (3.2) as in $`1D`$. The effective time of nonlinearity of the smallest $`p`$th mode in the vector field (5.3), in presence of all large scale modes, is determined by the equation (3.3). Thus, after the time of collision $`\stackrel{~}{t}_p`$, heavy particles with mass $`m_p`$ (5.7) appear. Velocities of these particles will be determined by the function $`𝐕_{p+1}\left(𝐱\right)(\text{5.3})`$, which is a sum of all larger modes, but the number of particle collisions is determined by the next $`p+1`$ mode. At time $`\stackrel{~}{t}_{p+1}`$ we have a collision of four heavy particles. The extrapolation of this particle merger process to next generations is evident by recurrence. The $`q`$th collision of heavy particles takes place at time $`\stackrel{~}{t}_{p+q}`$), (5.8) the masses of these particles at this time are determined by the scale of the $`(p+q)`$th mode $`m_{p+q}=\rho _0L_{p+q}^2\sqrt{3}/4`$ (5.7). Here also one can speak about the logarithmic periodic self-similarity of the motion of the particles. This means that at arbitrary interval $`[\stackrel{~}{t}_q,\stackrel{~}{t}_{q+1}]`$ the motion of the particles will be similar to the motion of the particles in the interval $`\stackrel{~}{t}_p,\stackrel{~}{t}_{p+1}`$, by the scaling factor $`x_p/x_q=2^{(qp)}`$ in space, and the scaling function $`t_p/t_q=(2^{\frac{n+3}{2}})^{(pq)}`$ in time. We used computer simulation for studying two dimensional case; results of the simulation were generated in so called VRML (Virtual Reality Modeling Language), which enables to handle three dimensional figure in different projections. Fig. 6 presents snapshots of this modeling. On the Fig. 6(a) one can see that particles formed by small triangles move towards the center of an embracing triangle; this center in its turn, moves towards the center of the next bigger triangle in the hierarchy; e.t.c. Fig.6(b) gives the side-view of this process. ## 6 Conclusion In conclusion, we would like to point out that the evolution of the multiscale signal with the Weierstrass spectrum simulates properties of Burgers turbulence such as self-similarity, conservation of large scale structures and has the same laws of the energy decay and integral scale. The difference between the deterministic model and Burgers turbulence is that here we have the exact solution for the evolution of multiscale signals and these properties are not stochastic but deterministic. The evolution of the multiscale signal is exactly self-similar in logarithmically spaced time intervals. The evolution of the large scale modes is completely independent of the small scales modes, even if these have very large amplitudes. These properties take place for Burgers turbulence in the stochastical sense and, moreover, for a signal with cutoff frequencies of small scales, only asymptotically. Of course, these properties of the multiscale signal are determined by the special form of modes (reverse sawtooth function), the special relations between wavenumbers of modes ($`k_{r+1}=k_r/N`$, where $`N`$ is an integer) and their phase relations. On the other hand, these model signals do not reflect such properties of Burgers turbulence as qualitative difference in the behaviour of the turbulence for $`n<1`$ and $`n>1`$ in the power spectrum (1.9) due to the process of generation of large scale components in the spectrum. For the deterministic model this process is not present due to the special relation of wavenumbers. Let as now move to the mechanical interpretation of solution of the Burgers equation. For the initial reverse sawtooth wave, all the matter turns into heavy particles at the same moment of time. Thus, after this time, the evolution of Burgers turbulence is fully determined by the motion of heavy particles. The trajectories of heavy particles form regular tree-like structure on the plane $`(X,t)`$, see Figs. 3 and 4. The properties of this structure depend on the parameters of our model. The integer $`ϵ^1=N`$ is the number of trajectories which intersect at one point and form a new branch of our structure. For $`ϵ^1=N=2`$ we, thus, obtain binary tree structure. Changing of the parameter $`h`$ stretches or contracts the structure in the $`t`$ direction. One can say that our structure is the plane representation of the $`N`$tree; the root of our tree is located at $`t=+\mathrm{}`$. This tree is similar to the flattened fractal model of botanical umbrella tree (see ). If we take some node $`(X,T)`$ of this structure as a root, and consider the trajectories of all heavy particles, which will merge at the moment of time $`T`$ at this point $`(X,T)`$ we shall also obtain an $`n`$-tree. The whole tree seems self-similar, because every branch plus the branches it carries is a reduced scale version of the whole. We also notice that such multiscale waves may be constructed for multidimensional Burgers’ equation. ## 7 Acknowledgements The authors are grateful to U. Frisch for useful discussions and for his hospitality at the Observatoire de la Côte d’Azur, to G.M. Molchan, A.I. Saichev, A. Noullez and W.A. Woyczynski for useful discussions. This work was partially supported by the French Ministry of Higher Education, by INTAS through grant No 97–11134, by RFBR through grant 99-02-18354. Figure captions Figure 1: Evolution of one-mode pulse. (a) particle trajectories; (b) evolution of the initial velocity field given in the same spatio-temporal scale. Bold lines on the time axis denote moments at which the profiles of the velocity are plotted. Figure 2: Evolution and interaction of two modes. (a) particle trajectories; (b) evolution of the initial velocity field given in the same spatio-temporal scale. Bold points on the time axis denote moments at which the profiles of the velocity are plotted. Figure 3: Evolution of the multiscale fractal signal with $`n=0(h=1/2)`$, corresponding to ”white noise” signal. (a) particle trajectories; (b) evolution of the initial velocity field; (b) velocity field taken at the initial moment of time and then at three successive time moments of self-similarity. Figure 4: Evolution of the multiscale fractal signal with $`n=2(h=1/2)`$, corresponding to ”Brownian motion” signal. (a) particle trajectories; (b) evolution of the initial velocity field; (b) velocity field taken at the initial moment of time and then at three successive time moments of self-similarity. Figure 5: Plane construction for two dimensional case. The hierarchy of triangles, used for the construction of the multiscale signal, with four layers shown. The initial signal is constructed as a series of signals piece-wise linear on triangles. Figure 6: Particle trajectories for the multiscale fractal signal in two dimensional case presented in spatio-time three dimensional space; the width of particle trajectory reflects its mass: (a) top view; (b) side-view.
warning/0002/quant-ph0002001.html	ar5iv	text	# Optimal States for Bell inequality Violations using Quadrature Phase Homodyne Measurements ## I Introduction There has been recently active interest in tests of quantum mechanics versus local realism in a high efficiency detection limit. Several authors including ourselves have considered detection schemes quadrature phase homodyne measurements. Such schemes use strong local oscillators and hence have very high detection efficiency. This removes one of the current loopholes and potentially allows a strong test of quantum mechanics to be performed. The original idea of Gilchrist et. al. was to use a circle or pair coherent state produced by nondegenerate parametric oscillation with the pump mode adiabatically eliminated. Using highly efficient quadrature phase homodyne measurements, the Clauser Horne strong Bell inequality could be tested in an all optical regime. A small (approximately $`1.5\%`$) but significant theoretical violation was found for this extremely idea system. While the mean photon number for the system may be low (approximately $`1.12`$), the use of homodyne measurements allow a macroscopic current to be detected. In this article, we take an unphysical but interesting approach and answer the following questions: * Given that your detection scheme is a quadrature phase homodyne measurement, what is the optimal input or correlated photon number state to maximize the potential violation? * What is the optimal Bell inequality to test? To begin we will restrict our attention to correlated photon number states of the form $`\|\mathrm{\Psi }={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}c_n\|n\|n`$ (1) Two main sources of correlated photon number currently exist, each having it own particular form of $`c_n`$. The most well known is simply the nondegenerate parametric amplifier specified by an ideal Hamiltonian of the form $`H=\mathrm{}\chi ϵ\left(ab+a^{}b^{}\right).`$ (2) where $`ϵ`$ is field amplitude of a nondepleting classical pump and $`\chi `$ is proportional to the susceptibility of the medium. $`a,b`$ are the boson operators for the orthogonal signal and idler modes. After a time $`\tau `$, the state of the system is given by (1) with $`c_n`$ specified by $`c_n={\displaystyle \frac{\mathrm{tanh}^n\left[\chi ϵ\tau \right]}{\mathrm{cosh}\left[\chi ϵ\tau \right]}}`$ (3) In the quadrature phase amplitude basis this state has a positive Wigner function. Hence it can be described as a local hidden variable theory and thus cannot violate a Bell inequality. The other source of highly correlated photon number states exists in nondegenerate parametric oscillation. In the limit of very large parametric nonlinearity and high Q cavities, a state of the form $`\|\mathrm{\Psi }={\displaystyle \frac{e^{r^2}}{\sqrt{4\pi ^2I_0\left(2r^2\right)}}}{\displaystyle _0^{2\pi }}𝑑\theta \|re^{i\theta }\|re^{i\theta }`$ (4) can be generated. Here $`r`$ is the size of the circle of the coherent states and $`I_0`$ is the zeroth order modified Bessel function. Equivalently this state can be written in the form of (1) with $`c_n`$ given by $`c_n={\displaystyle \frac{r^{2n}}{n!I_0\left(2r^2\right)}}`$ (5) This was the state considered by Gilchrist et. al. Given the general form of known correlated number states (1), the next fundamental question that should be initially addressed is what we mean by the Bell inequality. A number of Bell inequalities exist, and the particular one depends used heavy on your application and experimental setup. The Bell inequalities to be considered in this article are the Clauser Horne, the spin, and the information-theoretic Bell inequality. A detailed derivation of the various inequalities will not be given, the reader is referred to references . Here we will consider only strong inequalities, that is inequalities where auxiliary assumptions (not based on local realism) are not required. In Fig (1) we depict a very idealized setup for general Bell inequality experiment. Probably the most well known inequality is the Clauser Horne strong Bell inequality given by $`\|𝐁_{ch}\|1`$ (6) where $`B_{ch}={\displaystyle \frac{P_{11}(\theta ,\varphi )P_{11}(\theta ^{},\varphi )+P_{11}(\theta ,\varphi ^{})+P_{11}(\theta ^{},\varphi ^{})}{P_1\left(\theta ^{}\right)+P_1\left(\varphi \right)}}`$ (7) Here $`P_{11}`$ is the probability that a “1” results occurs at each analyzer $`A,B`$ given $`\theta ,\varphi `$. Similarly $`P_1`$ is the probability that a “1” occurs at a detector while having no information about the second. For many of the actual experimental considerations an angle factorization occurs so that $`P_{11}(\theta ,\varphi )`$ depends only on $`\theta +\varphi `$. Also $`P_1\left(\theta \right)`$ and $`P_1\left(\varphi \right)`$ are independent of $`\theta ,\varphi `$. In this case $`B_{ch}`$ can be simplified to $`B_{ch}={\displaystyle \frac{3P_{11}\left(\psi \right)P_{11}\left(3\psi \right)}{2P_1}}`$ (8) where $`\psi =\theta +\varphi =\theta ^{}\varphi ^{}=\theta +\varphi ^{}`$ and $`3\psi =\theta ^{}+\varphi `$. The second form of the Bell inequality (sometimes referred to as the spin or original Bell inequality) is $`B_s=\|E(\theta ,\varphi )`$ $``$ $`E(\theta ^{},\varphi )`$ (9) $`+`$ $`E(\theta ,\varphi ^{})+E(\theta ^{},\varphi ^{})\|2`$ (10) where the correlation function $`E(\theta ,\varphi )`$ is given by $`E(\theta ,\varphi )=P_{11}(\theta ,\varphi )`$ $`+`$ $`P_{00}(\theta ,\varphi )`$ (11) $``$ $`P_{10}(\theta ,\varphi )P_{01}(\theta ,\varphi )`$ (12) Here as discussed above $`P_{11}`$ is probability that a “1” results occurs at each analyzer $`A,B`$ given $`\theta ,\varphi `$. $`P_{00}`$ is probability that a “0” results occurs at each analyzer $`A,B`$, while $`P_{10}`$ ($`P_{01}`$) is probability that a “1” (“0”) results occurs at the analyzer $`A`$ and a “0” (“1”) at $`B`$. With the angle factorization given above, the inequality (9) can be rewritten as $`B_s=\left\|3E\left(\psi \right)E\left(3\psi \right)\right\|2`$ (13) Our final form of Bell inequality to be considered in this article was developed by Braunstein and Caves. This classical information-theoretic Bell inequality has the form $`B_{info}0`$ (14) where $`B_{info}=H\left(\theta \|\varphi \right)`$ $`+`$ $`H\left(\theta \|\varphi ^{}\right)`$ (15) $`+`$ $`H\left(\varphi ^{}\|\theta ^{}\right)+H\left(\theta ^{}\|\varphi \right)`$ (16) Here $`H\left(\theta \|\varphi \right)`$ is given by $`H\left(\theta \|\varphi \right)={\displaystyle \underset{a,b}{}}P(a,b)\mathrm{log}\left({\displaystyle \frac{P(a,b)}{P(a)}}\right)`$ (17) with $`\mathrm{log}\left(P(a,b)/P(a)\right)`$ being the information gained at $`B`$ given the result at $`A`$ is known. The conditional information is then given by $`H\left(\theta \|\varphi \right)`$. The base of the logarithm determines the units of the information (base 2 for bits, base e for nats). For quantum computing purposes, this inequality should prove highly useful as it directly deals with information content. Several other Bell inequality do exist such as the CHSH inequality, but these are not considered here due to there weaker nature. Auxiliary assumptions are necessary in there derivation which open up several loopholes. ## II Correlated States From (1) we need to find the optimal $`c_n`$ which gives the largest Bell inequality violation. Before determining the $`c_n`$ we need to briefly focus our attention on the quadrature phase homodyne measurement. A quadrature phase-amplitude homodyne measurement $`X(\theta )`$ at $`A`$ can achieved by combining a signal field (say $`\widehat{a}`$) with a strong local oscillator field (say $`ϵ`$) to form two new fields given by $`\widehat{c}_\pm =\left[\widehat{a}\pm ϵ\mathrm{exp}\left(i\theta \right)\right]/\sqrt{2}`$. Here $`\theta `$ is a phase shift which allows the choice of particular observable to be measured, for instance choosing $`\theta `$ as $`0`$ or $`\pi /2`$ allows the measurement of the conjugate phase variables $`X(0)`$ and $`X(\pi /2)`$ respectively. The homodyne measurement gives the photocurrent difference as $`I_d`$ $`=`$ $`c_+^{}c_+c_{}^{}c_{}`$ (18) $`=`$ $`ϵ\left(\widehat{a}e^{i\theta }+\widehat{a}^{}e^{i\theta }\right)=ϵX(\theta )`$ (19) Performing a measurement on the quadrature phase amplitude $`X(\theta )`$ at $`A`$ yields a result $`x_1(\theta )`$ which ranges in size and sign. Similarly a measurement on the quadrature phase amplitude $`X(\varphi )`$ at $`B`$ yields a result $`x_2(\varphi )`$. For our state given by (1), the probability of obtaining the result $`x_1(\theta ),x_2(\varphi )`$ is simply $`P_{x_1x_2}(\theta ,\varphi )=\left\|x_1(\theta )\|x_2(\varphi )\|\mathrm{\Psi }\right\|^2`$ (20) where $`x(\phi )\|n={\displaystyle \frac{1}{\sqrt{2^nn!\sqrt{\pi }}}}e^{in\phi }e^{x_i^2/2}H_n(x_i).`$ (21) Here $`H_n(x_i)`$ is the Hermite polynomial and $`\phi `$ is the phase of the local oscillator. Eqn (20) can be explicitly written as $`P_{x_1x_2}\left(\psi \right)`$ $`=`$ $`{\displaystyle \underset{n,m}{}}{\displaystyle \frac{c_nc_m^{}e^{i(nm)\psi }}{2^{n+m}n!m!\pi }}`$ (22) $`\times {\displaystyle \underset{i=1}{\overset{2}{}}}e^{x_i^2}H_n(x_i)H_m(x_i)`$ (23) where $`\psi =\theta +\varphi `$, that is our expression depends only on the sum of the individual local angles. The probability given by (22) is for continuous variables. The majority of the tests of quantum mechanics versus local realism require a binary result. Hence for a given quadrature measurement $`x_i`$ we classify the result as “1” if $`x_i0`$ and the mutually exclusive “0” if $`x_i<0`$. Here we have set the binning window about $`x_i=0`$. Where this binning window is located is quite arbitrary, but the maximum violation occurs for the value we have selected. The probability of obtaining both particles in the “1” bin is $`P_{11}\left(\psi \right)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑x_1𝑑x_2P_{x_1x_2}\left(\psi \right)`$ (24) while the probability of obtaining both particles in the “0” bin is $`P_{00}\left(\psi \right)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^0}{\displaystyle _{\mathrm{}}^0}𝑑x_1𝑑x_2P_{x_1x_2}\left(\psi \right)`$ (25) The other probabilities such $`P_{10}\left(\psi \right),P_{01}\left(\psi \right)`$ can be calculated in a similar fashion. The probabilities formulated above are joint probabilities. Various of the strong Bell inequalities also require marginal probabilities of the form $`P_1\left(\psi \right)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x_1𝑑x_2P_{x_1x_2}\left(\psi \right)`$ (26) The above integrals can be easy evaluated using the results $`{\displaystyle _0^{\mathrm{}}}e^{x^2}H_n(x)H_m(x)`$ $`=`$ $`{\displaystyle \frac{\pi 2^{n+m}}{nm}}\left[(n,m)(m,n)\right]`$ (27) $`(fornm)`$ (28) $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}e^{x^2}H_n(x)H_m(x)`$ $`=`$ $`2^nn!\sqrt{\pi }\delta _{n,m}`$ (29) where $`(n,m)`$ is given by $`^1(n,m)=\mathrm{\Gamma }\left({\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{2}}n\right)\mathrm{\Gamma }\left({\displaystyle \frac{1}{2}}m\right)`$ (30) with $`\mathrm{\Gamma }`$ being the Gamma function. Performing the integrals for (24) and (25) we find $`P_{11}\left(\psi \right)`$ $`=`$ $`P_{00}\left(\psi \right)`$ (31) $`=`$ $`{\displaystyle \frac{1}{4}}+{\displaystyle \underset{n>m}{}}{\displaystyle \frac{2^{n+m+1}\pi c_nc_m^{}}{n!m!(nm)^2}}`$ (32) $`\times \left[(n,m)(m,n)\right]^2\mathrm{cos}\left[(nm)\psi \right]`$ (33) Similarly Eqn (26) simplifies to $`P_1`$ $`=`$ $`1/2`$ (34) which is independent of the sum of the local oscillator angle $`\psi `$. It is also simple to calculate the correlation function $`E(\psi )`$ $`E(\psi )`$ $`=`$ $`{\displaystyle \underset{n>m}{}}{\displaystyle \frac{2^{n+m+3}\pi c_nc_m^{}}{n!m!(nm)^2}}\left[(n,m)(m,n)\right]^2`$ (35) $`\times \mathrm{cos}\left[(nm)\psi \right]`$ (36) Given the probabilities $`P_{11}`$, $`P_{00}`$, $`\mathrm{}`$ it is also possible to calculate the conditional information $`H\left(\theta \|\varphi \right)`$ $`H\left(\theta \|\varphi \right)=`$ $``$ $`P_{11}\mathrm{log}\left[2P_{11}\right]P_{00}\mathrm{log}\left[2P_{00}\right]`$ (37) $``$ $`P_{10}\mathrm{log}\left[2P_{10}\right]P_{01}\mathrm{log}\left[2P_{01}\right]`$ (38) It is now possible to calculate the Clauser Horne (6) and spin (9) and information-theoretic (14) Bell inequalities. Some insight into the problem can be achieved by a careful examination of the term $`{\displaystyle \frac{2^{n+m}\pi }{n!m!(nm)^2}}\left[(n,m)(m,n)\right]^2`$ (39) which is present in all the joint probability distributions. This expression has several interesting features. First, as the difference between $`n`$ and $`m`$ becomes large, the smaller that the above expression contributes to any of the probability distributions. The main contribution for the expression comes from the case $`m=n\pm 1`$. Second, when $`nm`$ is even, the above expression is zero. Finally, as $`n`$ is large the different between the $`n,m=n1`$ and $`n+1,m=n`$ elements for fixed large $`n`$ vanishes and they reach an asymptotic limit which is smaller than the $`n=1,m=0`$ case. If these higher order $`n`$ terms dominate due to the choice of the $`c_n`$ in the probability formula, then the various Bell inequalities can not violated. This also has the implication that the mean photon number cannot be high if a violation is to occur and hence it is not a macroscopic test of quantum mechanics. ## III A simple Case To begin our investigations of the Bell inequalities, consider the case of we have only two photon pair states that is, $`\|\mathrm{\Psi }=c_0\|0\|0+c_1\|1\|1`$ (40) where for convenience we choose $`c_n`$ real. We also require $`c_0^2+c_1^2=1`$. The joint probability distributions are readily calculated and in fact $`P_{11}\left(\psi \right)=P_{00}\left(\psi \right)`$ $`=`$ $`{\displaystyle \frac{1}{4}}+{\displaystyle \frac{c_0c_1}{\pi }}\mathrm{cos}\left[\psi \right]`$ (41) $`P_{10}\left(\psi \right)=P_{01}\left(\psi \right)`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{c_0c_1}{\pi }}\mathrm{cos}\left[\psi \right]`$ (42) Calculating $`B_{ch}`$ and $`B_s`$ from (6) and (9) we find $`B_{ch}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{c_0c_1}{\pi }}\left\{3\mathrm{cos}\left[\psi _0\right]\mathrm{cos}\left[3\psi \right]\right\}`$ (44) $`B_s`$ $`=`$ $`{\displaystyle \frac{4c_0c_1}{\pi }}\left\{3\mathrm{cos}\left[\psi _0\right]\mathrm{cos}\left[3\psi \right]\right\}`$ (45) Optimizing for the angle $`\psi `$ we find $`B_{ch}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+{\displaystyle \frac{2\sqrt{2}c_0c_1}{\pi }}`$ (47) $`B_s`$ $`=`$ $`{\displaystyle \frac{8\sqrt{2}c_0c_1}{\pi }}`$ (48) that is, $`\|B_{ch}\|1`$ and $`\|B_s\|2`$ for all $`c_0`$. No violation of the strong Clauser Horne or spin Bell inequality is possible. For the information theoretic case we find $`H\left(\psi \right)=`$ $``$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left[{\displaystyle \frac{1}{4}}\lambda ^2\right]\lambda \mathrm{log}\left[{\displaystyle \frac{1+2\lambda }{12\lambda }}\right]`$ (50) where $`\lambda ={\displaystyle \frac{2c_0c_1}{\pi }}\mathrm{cos}\left[\psi \right]`$ (51) The informational theoretic Bell inequality is given by $`B_{info}=3H\left(\psi \right)H\left(3\psi \right)0`$. A violation of this inequality is possible if $`B_{info}<0`$. Unfortunately for all $`c_0c_1`$ and $`\psi `$ we have $`B_{info}>0`$. No violation is possible for any of the Bell inequalities considered for the ideal state (40) when the detection scheme is based on homodyne quadrature phase measurements. If more correlated photon pairs are present can a violation be achieved? The obvious answer is yes, because of the recent work of Gilchrist et. al.. The real question is how large this violation is? ## IV Numerical Studies Considering the expression (1) for the correlated photon pairs, what are the optimal $`c_n`$ coefficients to maximize the violation. Because of the results indicated by Gilchrist et. al. and our previous discussion we anticipate that the mean photon number per mode must be low to obtain a violation. Hence we will truncate the number state basis at $`10`$ photon per mode. Performing a numerical optimization over all the $`c_n`$, the optimal set is found to maximize the Clauser Horne and spin Bell inequality (Table I). A plot of $`c_n`$ versus $`n`$ is depicted in Fig (2) It is interesting to now discuss some properties of these optimal $`c_n`$. First the general shape of the $`c_n`$ versus $`n`$ curve shown in (2) is similar to that considered in the circle state by Gilchrist et. al.. It is however not exactly the same (see (Table I). Given this optimal parameter set, what is the maximum violation of the Bell inequalities we are considering. In Fig (3) we plot both the Clauser Horne and spin Bell inequalities versus $`\psi `$. For the Clauser Horne Bell inequality the maximum violation corresponds to $`B_{ch}=1.019`$, while the maximum violation for the spin Bell inequality corresponds to $`B_s=2.076`$. Interesting here is that the percentage violation of the spin inequality is approximately $`3.8\%`$ compared with the $`1.9\%`$ for the Clauser Horne case. This significantly increases the potential for an experiment to be performed provided such an experiment were not significantly more difficult. Also the results for the optimal $`c_n`$ set give a Clauser Horne Bell inequality violation that is approximately $`20\%`$ greater than the circle state results of Gilchrist et. al.. It is interesting to consider whether a greater violation of the Bell inequality can be achieved with the state given by (5). To this end we show the effect of the variation of both $`r`$ and $`\psi `$ (sum of the local oscillator angles) for both the Clauser Horne and spin Bell inequalities in Fig (4). As can be seen the spin Bell inequality can be violated far more significantly than the similar Clauser Horne case. In fact, as occurred previously the percentage maximum violation in the spin inequality is twice that of the Clauser Horne result. In any of the analysis considered above we have not discussed errors, there sources and how they effect the potential violation. We will not present any significant details here in this article but refer the reader to for such a decision. Our final Bell inequality to be considered is the Braunstein and Caves information-theoretic case. In Fig (5) we plot $`B_{info}`$ versus $`\psi `$. No violation of the information-theoretic inequality is possible for any $`\psi `$. A question here to be addressed is why two of the strong inequalities can be violated while this information-theoretic Bell inequality is far from being violated. In the binning process to give a binary result for a quadrature measurement, information must be discarded. The information-theoretic inequality is much more sensitive to this information loss than the Clauser Horne inequality. Also why would we fundamentally expect all three inequalities to be violated. A violation of any of the inequalities indicate a discrepancy between quantum mechanics and local realism. ## V Conclusion In this article we have place strict bounds on the optimal $`c_n`$ coefficients for the state (1) which maximizes the Clauser Horne and Spin Bell inequalities when a homodyne quadrature phase measurements is performed. The spin Bell inequality is violated by approximately $`3.6\%`$ while the Clauser Horne inequality is violated by approximately $`1.9\%`$. The violation is small however due to the fact that we are discarding information in the binning process. In fact due to the information loss in the binning process the information theoretic Bell inequality is not violated in any regime. A larger violation cannot be obtained using homodyne measurements with the strong inequalities we have considered. While our optimal $`c_n`$ coefficient give a slightly better violation than the pair coherent state, it is difficult to see how such a state could be generated. Closely examining the spin Bell inequality with the pair coherent state still indicates that a greater violation (approximately twice the size) is possible than for the other inequalities. This would make the test much more feasible provided the pair coherent state could be generated. In such a system the mean photon number is small, so this is not strictly a macroscopic test of quantum mechanics. It does however have a macroscopic nature due the strong local oscillator which means large photodetector currents are obtained. To conclude, quadrature phase homodyne measurement provide a mechanism for performing tests of the Bell inequality with highly efficient detection. This allows one of the loopholes in current experiments to be closed. However, due to the inherent information loss in the binning process, the violations are small but should be achievable. .
warning/0002/astro-ph0002338.html	ar5iv	text	# Discovery of a Color-Selected Quasar at 𝑧=5.50Based on observations at the W.M. Keck Observatory, Kitt Peak National Observatory, and Palomar Observatory. Keck Observatory is operated as a scientific partnership among the University of California, the California Institute of Technology, and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation. ## 1 Introduction The past few years have witnessed a watershed in our direct observations of the high-redshift Universe. A decade ago, only a handful of galaxies were identified past a redshift of 3. These sources represented the rare beast in the cosmos: high-redshift radio galaxies, or galaxies associated with extremely distant, luminous quasars. New techniques and instruments allow us now to routinely identify normal, star-forming galaxies at these same epochs (e.g.,Steidel et al. 1996). Stern & Spinrad (1999) review modern search techniques for distant galaxies. Improved computing power and ambitious, large-area surveys also have pushed the frontier of distant quasar studies (e.g.,Djorgovski et al. 1999; Fan et al. 1999). Now we are regularly identifying objects which have collapsed only $`1`$ Gyr after the Big Bang. Such observations tell us about the earliest phases of galaxy and structure formation and probe the conditions of the early Universe. The identification of high-redshift quasars is especially important for several reasons. First, quasars at early cosmic epoch require the rapid formation of a supermassive black hole. Assuming black holes are not primordial, this requires the condensation of a large cloud of hydrogen, presumably embedded within a dark matter halo. Additionally, the presence of metal lines in quasars demand a previous generation of stars (two generations for nitrogen). High-redshift quasars thus constrain models of galaxy and structure formation (e.g.,Loeb 1993; Eisenstein & Loeb 1995). Also, quasars provide valuable probes of the intervening intergalactic medium (e.g.,Rauch 1998) and the intergalactic ionizing background. For example, the absence of a smooth depression in quasar continua short-ward of the Ly$`\alpha `$ emission strongly constrains the amount of neutral hydrogen in the intergalactic medium (Gunn & Peterson 1965). Songaila et al. (1999) find no Gunn-Peterson trough out to redshift 5 from deep spectroscopic observations of SDSSp J033829.31+002156.3 at $`z=5.00`$ (Fan et al. 1999). In this Letter, we report the discovery of a quasar at $`z=5.50`$, the most distant quasar identified to date. The previous most distant quasar was SDSSp J120441.73$``$002149.6 at $`z=5.03`$ (Fan et al. 2000). At $`z=5.50`$, an $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,\mathrm{\Lambda }=0,\mathrm{\Omega }=1`$ (0.1) universe is 790 Myr (1.51 Gyr) old, corresponding to a look-back time of 94.0% (90.9%) of the age of the universe. For the lambda cosmology supported by recent studies of distant supernovae, $`H_0=65\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,\mathrm{\Lambda }=0.7,\mathrm{\Omega }_\mathrm{m}=0.3`$, the universe is 1.11 Gyr old at $`z=5.50`$, corresponding to a look-back time of 92.4% of the age of the universe. ## 2 Observations and Target Selection RD J030117+002025 was identified from deep $`RIz`$-band imaging using a slightly redder version of the ‘dropout’ color selection techniques which have proved successful at identifying high-redshift galaxies (e.g.,Steidel et al. 1996; Dey et al. 1998; Spinrad et al. 1998) and quasars (e.g.,Kennefick, Djorgovski, & de Calvalho 1995; Djorgovski et al. 1999; Fan et al. 1999; Stern et al. 2000b). The selection criteria rely upon absorption from the Ly$`\alpha `$ and Ly$`\beta `$ forests attenuating the rest-frame ultraviolet continua. At $`z5`$, such objects will disappear from the $`R`$-band. Long-ward of the redshifted Ly$`\alpha `$, both quasars and star-forming galaxies display relatively flat (in $`f_\nu `$) continua. In concept, our survey is similar to established quasar surveys relying upon the digitized Palomar Sky Survey (e.g.,Djorgovski et al. 1999) or Sloan Digital Sky Survey (e.g.,Fan et al. 1999). In practice, we probe a much smaller area of sky (eventually a few $`\times `$ 100 arcmin<sup>2</sup>) to much fainter magnitudes. Our survey is designed to study the high-redshift, “normal” galaxy population, but is also sensitive to (low-luminosity) high-redshift quasars. The $`Iz`$ imaging was obtained using the Kitt Peak National Observatory 150″ Mayall telescope with its Prime Focus CCD imager (PFCCD) equipped with a thinned AR coated $`2048\times 2048`$ Tektronics CCD. This configuration gives a $`14.3\times 14.3`$ arcminute field of view with 0$`\stackrel{}{\mathrm{.}}`$43 pixels. The CCD was operated using “short scan”, where the CCD was mechanically displaced while its charge is shifted in the opposite direction to reduce fringing at $`I`$ and $`z`$ to very low levels. Two hours of Mould $`I`$-band ($`\lambda _c=8200`$ Å; $`\mathrm{\Delta }\lambda =1820`$ Å) data were obtained on UT 1995 August 31. The $`z`$-band (RG850, long-pass filter) data were obtained during UT 1997 November $`46`$, and the summed image represents 3.3 hours of integration. The combined, processed $`I`$ and $`z`$ images reach limiting magnitudes of 25.7 and 24.8 mag, respectively (3$`\sigma `$ limits in 3″ diameter apertures; AB magnitudes are used throughout this Letter) and have 0$`\stackrel{}{\mathrm{.}}`$9 and 1$`\stackrel{}{\mathrm{.}}`$2 seeing, respectively. These images comprise one field in the $`BRIzJK`$ Elston, Eisenhardt, & Stanford (2000) field galaxy survey. On UT 1999 November $`1112`$, we used the COSMIC camera (Kells et al. 1998) on the 200″ Hale telescope at Palomar Observatory to obtain extremely deep (4.4 hour) Kron-Cousins $`R`$-band ($`\lambda _c=6200`$ Å; $`\mathrm{\Delta }\lambda =800`$ Å) imaging of the same field, with the purpose of identifying high-redshift candidates. COSMIC uses a $`2048\times 2048`$ pixel SITe (formerly Tektronix) thinned CCD with 0$`\stackrel{}{\mathrm{.}}`$2846 pixels, yielding a $`9.7\times 9.7`$ arcmin field-of-view. Our combined, processed $`R`$-band image has 1$`\stackrel{}{\mathrm{.}}`$2 seeing and reaches a depth of 26.3 mag (3$`\sigma `$ limit in 3″ diameter aperture). High-redshift candidates for spectroscopy, designated RD for $`R`$-drop, were identified on the basis of a strong $`RI`$ color index and relatively flat $`Iz`$ color. No morphological criteria were implemented as the primary goal of this program is to study “normal,” star-forming galaxies at high redshift. Candidates were then screened by eye, yielding a total of six good targets over the central 74 arcmin<sup>2</sup> field. Fig. 1 presents a finding chart for RD J030117+002025, the brightest of our candidates and the subject of this Letter. Other candidates will be discussed in a future publication. We obtained spectra of several $`R`$-band dropouts through 1$`\stackrel{}{\mathrm{.}}`$5 wide, 13″ $``$ 44″ long slitlets using the Low-Resolution Imaging Spectrometer (LRIS; Oke et al. 1995) on the Keck II telescope on UT 2000 January 10 and 11. Observations were obtained at a position angle of $``$111.6 (east of north) with the 150 lines mm<sup>-1</sup> grating ($`\lambda _{\mathrm{blaze}}=7500`$ Å; $`\mathrm{\Delta }\lambda _{\mathrm{FWHM}}17`$ Å). The spectra sample the wavelength range 4000 Å to 1$`\mu `$m. Seeing was $``$ 1$`\stackrel{}{\mathrm{.}}`$1 during both nights and conditions were photometric. We performed $``$ 3″ spatial offsets between each 1800 s exposure in order to facilitate removal of fringing at long wavelength ($`\lambda \mathrm{}>7200`$ Å). All data reductions were performed using IRAF and followed standard slit spectroscopy procedures. We calculated the dispersion using a HgNeArKr lamp spectrum observed immediately subsequent to the science observations (RMS variations of 0.6 Å), and employed telluric emission lines to adjust the wavelength zero-point. The spectra were flux-calibrated using observations of Feige 67 and Feige 110 (Massey & Gronwall 1990). We corrected for foreground Galactic extinction using a reddening of $`E_{\mathrm{B}\mathrm{V}}=0.03`$ determined from the dust maps of Schlegel, Finkbeiner, & Davis (1998). The final composite spectrum of RD J030117+002025, presented in Fig. 2, represents 4.5 hours of integration. ## 3 Results and Discussion Though of moderate signal-to-noise ratio, the spectrum of RD J030117+002025 has the unambiguous signature of an extremely distant quasar. The broad emission with a sharp absorption at 7900 Å is consistent with Ly$`\alpha `$/N 5 $`\lambda `$1240 emission attenuated by the nearly opaque Ly$`\alpha `$ forest at $`z=5.50`$. An additional discontinuity is visible at 6690 Å, associated with the Ly$`\beta `$ forest. The Ly$`\alpha `$ forest absorption and poor detection of the long-wavelength Si 4/O 4\] $`\lambda `$1403 emission make centroiding on the emission features ill-advised; the redshift is instead determined from the sharp forest decrements. We estimate $`z=5.50\pm 0.02`$. This and other properties of RD J030117+002025 are given in Table 1. The spectral character of RD J030117+002025 is slightly atypical of high-redshift quasars, though it undoubtedly resides within the diverse category of quasars. The Ly$`\alpha `$/N 5 complex is unusually broad and distinguishing the emission lines is impractical. Many of the highest redshift quasars share similar spectroscopic shapes, e.g., SDSSp J033829.31+002156.3 at $`z=5.00`$, SDSSp J021102.72$``$000910.3 at $`z=4.90`$ (Fan et al. 1999), and GB 1428+4217 at $`z=4.72`$ (Hook & McMahon 1998). The strong continuum absorption associated with the Ly$`\alpha `$ forest is the dominant spectroscopic feature of RD J030117+002025. A robust determination of $`D_A`$, the standard parameter for describing the Ly$`\alpha `$ forest decrement (Oke & Korycansky 1982), requires knowledge of the continuum spectral slope long-ward of Ly$`\alpha `$. We estimate $`D_A`$ by assuming the standard quasar power law spectral index of $`0.5`$ for the continuum long-ward of Ly$`\alpha `$ (e.g.,Richstone & Schmidt 1980; Schneider et al. 1992), with the amplitude determined over the wavelength interval $`\lambda \lambda 84009000`$ Å (between the N 5 $`\lambda `$1240 and Si 4/O 4\] $`\lambda `$1403 emission complexes). We derive $`D_A=0.90\pm 0.02`$. We derive $`D_B=0.95\pm 0.04`$ for the strength of the Ly$`\beta `$ forest. The $`D_A`$ value is comparable to those measured for distant galaxies in the Hubble Deep Field at similar redshifts (Weymann et al. 1998; Spinrad et al. 1998) and models of the Ly$`\alpha `$ forest (Madau 1995; Zhang et al. 1997). At $`z=5.50`$, the features used to describe continuum properties are redshifted to challenging wavelengths. We estimate $`AB_{1450(1+\mathrm{z})}`$ using the continuum modeled above. Consistent with previous work in this field, $`M_B`$ is calculated for an Einstein-de Sitter universe with $`H_0=50\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1,q_0=0.5`$ and the standard quasar power law index of $`0.5`$. We find $`M_B22.7`$; see Stern et al. (2000a) for details of how we calculate $`M_B`$. Comparison with the 1.4 GHz FIRST survey (Becker, White, & Helfand 1995) reveals no radio source within 30″ of the quasar to a limiting flux density of $`f_{1.4\mathrm{GHz}}1`$ mJy (5$`\sigma `$). How unusual is it to find a quasar as distant and luminous as RD J030117+002025 in a $`100`$ arcmin<sup>2</sup> field? This is difficult to answer as RD J030117+002025 is the least luminous $`z>4`$ quasar known. The previously known high-redshift ($`z>4`$) quasars of lowest luminosity are PC 0027+0521 ($`z=4.21,M_B=24.0`$; Schneider, Schmidt, & Gunn 1994), which was discovered serendipitously, and the X-ray selected quasar RX J105225.9+571905 ($`z=4.45,M_B=23.9`$; Schneider et al. 1998). These luminosities are comparable to the lower luminosity objects in the $`z\mathrm{}<1`$ Bright Quasar Survey (Schmidt & Green 1983). Most high-redshift quasar luminosity functions (e.g.,Schmidt, Schneider, & Gunn 1995) have been derived from samples of quasars $`100`$ times as luminous as RD J030117+002025. To calculate the expected surface density of high-redshift, faint quasars, we follow the methodology outlined in Kennefick et al. (1995) and Boyle & Terlevich (1998): we adopt the Boyle et al. (1991) $`z=2`$ quasar luminosity function (for $`q_0=0.5`$), scaled down in density using the evolution predicted by Schmidt et al. (1995), namely, that the quasar space density falls off by a factor of 2.7 per unit redshift beyond $`z=3`$. The predicted surface density of $`R`$-drop ($`4.3\mathrm{}<z\mathrm{}<5.8`$) quasars with $`M_B22.5`$ is $`2\times 10^3`$ arcmin<sup>-2</sup>, implying $`0.15`$ such quasars should have been uncovered in our survey. It is dangerous, yet enticing, to draw conclusions from a single object. The discovery of the quasar RD J030117+002025 at $`z=5.50`$ in the modest sky coverage of our survey is suggestive of less dramatic evolution in the quasar luminosity function at faint magnitudes and high redshift. Such a change would have significant cosmological implications, including changing the budget of high-energy, ionizing photons in the early Universe. We also note that the low signal-to-noise ratio data are suggestive of strong hydrogen absorption near the quasar redshift. Is this due to a neutral hydrogen cloud near the quasar, at odds with the proximity effect? Or are we seeing the first glimpses of an object radiating prior to the reionization epoch, with neutral intergalactic hydrogen absorbing the rest-frame UV photons? Higher resolution, higher signal-to-noise ratio data will be essential for answering these questions. ###### Acknowledgements. We are indebted to the expertise of the staffs of Kitt Peak, Palomar, and Keck Observatories for their help in obtaining the data presented herein, and to the efforts of Bev Oke and Judy Cohen in designing, building, and supporting LRIS. We especially thank Barbara Schaeffer, Greg Wirth, and Jerome at Keck II for their assistance during the January 2000 observing run. We are grateful to Carlos DeBreuck and Richard McMahon for carefully reading the manuscript and to the referee, Ray Weymann, for prompt and helpful comments. Portions of this work were carried out by the Jet Propulsion Laboratory, California Institute of Technology, under a contract with NASA. Portions of this work was performed under the auspices of the U.S. Department of Energy by University of California Lawrence Livermore National Laboratory under contract No. W-7405-Eng-48. This work has been supported by the following grants: NSF grant AST 95$``$28536 (HS), the Cambridge Institute of Astronomy PPARC observational rolling grant ref. no. PPA/G/O/1997/00793 (AJB), and NSF CAREER grant AST 9875448 (RE).
warning/0002/cond-mat0002249.html	ar5iv	text	# Phase behaviour of a model of colloidal particles with a fluctuating internal state ## 1 Introduction Theories for the behaviour of systems of particles usually apply to a model for the interaction between particles in which the energy of interaction is pairwise additive. The Hamiltonian is of the form $$H(𝐫^N)=\underset{i,j=1}{\overset{N^{}}{}}u(r_{ij}),$$ (1) a sum of spherically-symmetric pairwise-additive potentials $`u(r)`$. $`𝐫^N`$ denotes the coordinates of the $`N`$ particles and $`r_{ij}`$ is the scalar separation of particles $`i`$ and $`j`$. But experiments are almost never on simple rigid particles. For example, the particles in a colloidal suspension are not simple rigid particles. Generally they are stabilised in the suspension either by being charged, in which case they are surrounded by a cloud of counterions, or by having short polymers grafted to their surface . In either case we have not a rigid particle but a system with many degrees of freedom, e.g., the positions of the counterions, which fluctuate. Each particle is a system, with a free energy, susceptibilities etc., in its own right. Here we will treat systems of particles, each of which is a weakly fluctuating statistical mechanical system in its own right. Each particle will have a Landau-like free energy which couples to that of neighbouring particles. We will use the simplest possible form of this free energy, and show exactly and analytically that it leads to many-body attractions. These cannot be expressed as a Hamiltonian with the form of Eq. (1). This work is partly inspired by recent experiments on highly charged colloidal particles under conditions of minimal amounts of added salt. The potential of mean force between an isolated pair of colloidal particles with minimal salt concentrations has been measured and is purely repulsive . Yet the particles form crystallites which appear to be metastable at close to zero osmotic pressure ; this is very hard to explain unless there is some sort of cohesive attraction between the particles in the crystallites. It does not seem possible to explain these findings with a Hamiltonian which depends only on the coordinates of the colloidal particles and has the conventional form of Eq. (1). Here we develop a simple phenomenological theory for colloidal particles which treats each particle as fluctuating statistical mechanical system. See Refs. for recent theoretical work on understanding this behaviour starting from a Hamiltonian which explicitly includes the counterions and the electrostatic interactions. The theory is phenomenological as we simply assume that the state of a particle can be described by a single coarse-grained scalar variable, we do not derive this. This scalar variable fluctuates and these fluctuations are perturbed by the presence of another particle nearby. The perturbations due to the particles which surround any given particle add up, meaning that the state of a particle changes with the number of its neighbours. Essentially, the more particles that surround any given particle the more the particle’s fluctuations are biased towards values that minimise the interaction free energy and so the more attractive is the interaction between the given particle and all the surrounding particles. It is possible for two of our model particles to repel each other but for particles in clusters of more than two particles to attract each other. This can lead to coexistence between a dilute fluid and a dense crystal phase even when a pair of particles repel each other. We will show that generally if the fluctuations can be described by a scalar variable then the interaction between a pair of particles within a cluster of several particles is more attractive than between an isolated pair. In the next section we define our model. In section 3 we show that it exactly corresponds to a potential between structureless particle which contains both pair and triplet terms. Then in section 4 we apply a perturbation theory to obtain an approximation to the free energy of both the fluid and crystalline phases, which is accurate for long-ranged interactions between the fluctuating internal states. We show results for the phase behaviour and for the zero wavevector structure factor in section 5. Section 6 is a conclusion. ## 2 Model Our model particles have an internal state specified by a single scalar variable. The value of the variable for a single isolated particle fluctuates weakly. Essentially, we view it as a coarse-grained variable obtained by averaging over some large number of degrees of freedom associated with each particle. An outline of how our mesoscopic Hamiltonian may be derived from a microscopic one is given in the Appendix. The interaction between a pair of the particles depends on the value of this variable. Thus, when two particles interact the mean values of the variables on the two particles will change; the particles ‘polarise’ each other, the interaction between them biases both internal variables towards values for which their interaction free energy is low. We included polarise in quotes because our model is phenomenological not electrostatic; cf Ref. which considers the polarisation of one a charged colloidal particle by another. The effective Hamiltonian $`H`$ of the system of particles has two parts: one part independent of the $`s`$ variables, $`U`$, and the other is the free energy of the particles as a function of the $`N`$ internal variables, $`F_N`$, $$H(𝐫^N,s^N)=U(𝐫^N)+F_N(𝐫^N,s^N),$$ (2) where $`𝐫^N`$ and $`s^N`$ symbolise the centre-of-mass and internal-variable coordinates, respectively, of all $`N`$ particles. We need not specify $`U`$ beyond saying that it should be such that the system has a well-defined thermodynamic limit . Later on we will set $`U`$ to be the sum of hard-sphere repulsions between pairs of the particles, $$U(𝐫^N)=\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}u_{hs}(r_{ij}),$$ (3) where $`u_{hs}`$ is the hard sphere potential $$u_{hs}(r)=\{\begin{array}{cc}0\hfill & r\sigma \hfill \\ \mathrm{}\hfill & r<\sigma \hfill \end{array}.$$ (4) $`r`$ is the separation of the centres of the interacting particles. The dash over the sum in Eq. (3) means that the sum is only over those terms for which $`ji`$. Note that not all the double and triple sums below will have a dash, in the undashed ones terms with equal subscripts are summed over. Each particle has a dimensionless internal variable $`s`$. Now, we assume that the coupling between these variables on different particles is pairwise additive and spherically symmetric. Then the free energy $`F_N`$ is, $$F_N(𝐫^N,s^N)=\underset{i=1}{\overset{N}{}}f^{(1)}(s_i)+\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}f^{(2)}(s_i,s_j,r_{ij}).$$ (5) $`f^{(1)}`$ is the free energy of a single isolated particle and is a function only of its $`s`$ variable. $`f^{(2)}`$ is the difference in free energy between an isolated pair of particles and two isolated particles, and is a function only of the two $`s`$ variables and the magnitude of the separation of the two particles. An isolated particle is a particle far from any other particle, an isolated pair of particles is a pair of particles far from any other particle. For weak fluctuations $`s`$ is so we Taylor expand $`f^{(1)}`$ and $`f^{(2)}`$ $$f^{(1)}(s)=\alpha s^2+\alpha _3s^3+\alpha _4s^4+\mathrm{},$$ (6) where the linear term is missing as the variable $`s`$ is defined so that if a particle is isolated its mean value is zero, and $$f^{(2)}(s,s^{},r)=\varphi _0(r)+\varphi _1(r)(s_i+s^{})+\varphi _2(r)(s^2+s^2)+\varphi _2^x(r)ss^{}+\mathrm{}.$$ (7) The coefficients $`\alpha `$, $`\alpha _3`$ etc. are derivatives of $`f^{(1)}`$ $$\alpha =\frac{1}{2}\frac{\mathrm{d}^2f^{(1)}(s)}{\mathrm{d}s^2},$$ (8) and $`\alpha _3`$ is $`1/6`$ times the third derivative etc.. Similarly, the coefficients $`\varphi _0(r)`$, $`\varphi _1(r)`$ etc. are derivatives of $`f^{(2)}`$ at a fixed separation of the particles. $`\varphi _0`$ is the zeroth derivative, $`\varphi _1`$ is the first derivative with respect to either of the two $`s`$ variables, $$\varphi _1(r)=\left(\frac{f^{(2)}(s,s^{},r)}{s}\right)_{s^{},r},$$ (9) and similarly for the higher derivatives $`\varphi _2`$ etc.. Our Taylor expansions are quite general for particles which have a state which can be described by a single weakly fluctuating scalar variable and in which the interaction between two particles can be expressed in terms of this variable. If the variable is a vector or a tensor not a scalar then clearly there will be expressions analogous to that of Eq. (5) but vectorial or tensorial variables will in general lead to results very different from those found here. We truncate the Taylor expansions, Eqs. (6) and (7), after their lowest nontrivial terms and substitute the resulting expansions into Eq. (5) to obtain $$F_N(𝐫^N,s^N)=\underset{i=1}{\overset{N}{}}\alpha s_i^2+\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\left[\varphi _0(r_{ij})+\varphi _1(r_{ij})\left(s_i+s_j\right)\right].$$ (10) The inverse susceptibility, $`\alpha `$, and the zeroth, $`\varphi _0`$, and first, $`\varphi _1`$, order coefficients all have the dimensions of energy. So, for our Hamiltonian, Eqs. (2) and (10), to describe a system of particles accurately we require: (1) that the fluctuations in $`s`$ of an isolated particle be sufficiently small that the higher order terms in $`f^{(1)}`$ may be neglected, (2) that the interactions between the particles can be described as a sum over interactions of pairs, (3) that this interaction is spherically symmetric, (4) that this interaction can be accurately described by function of a single scalar, coarse-grained, variable, and (5) that the interactions between fluctuations are sufficiently weak that the higher terms in $`f^{(1)}`$ and $`f^{(2)}`$ are small. For noninteracting particles the second sum of Eq. (10) may be neglected and $`F_N`$ is just a sum of independent quadratic terms. The distribution of the $`s`$’s is then Gaussian, which is correct for small fluctuations of a coarse-grained variable . $`\alpha `$ is an inverse susceptibility of an isolated particle, the smaller it is the larger are the fluctuations. The interaction between a pair of particles is expressed as a Taylor expansion in the two $`s`$’s truncated after the linear term. One particle feels an interaction due to another nearby particle, which couples to its order parameter $`s`$ with a strength $`\varphi _1(r)`$. The $`s`$-dependence of the free energy of Eq. (10) is what we would expect if the particles were weakly interacting, weakly fluctuating thermodynamic systems. ## 3 Exact theory We start from the configurational integral $`Z_N`$ for $`N`$ particles in a volume $`V`$ and at a temperature $`T`$ $$Z_N=d𝐫^Nds^N\mathrm{exp}(H(𝐫^N,s^N)/kT),$$ (11) with the Hamiltonian given by Eqs. (2) and (10). The Helmholtz free energy $`A`$ is then $$A=kT\mathrm{ln}\left(Z_N\mathrm{\Lambda }^N/N!\right),$$ (12) where $`\mathrm{\Lambda }`$ derives from the integration over the momenta. Using Eqs. (2) and (10) in Eq. (11), $`Z_N`$ $`=`$ $`{\displaystyle d𝐫^N\mathrm{exp}(U(𝐫^N)/kT)ds^N\mathrm{exp}\left(\underset{i=1}{\overset{N}{}}(\alpha /kT)s_i^2\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\left[\varphi _0(r_{ij})/kT+\varphi _1(r_{ij})\left(s_i+s_j\right)/kT\right]\right)}`$ (13) $`=`$ $`{\displaystyle d𝐫^N\mathrm{exp}\left(\frac{U(𝐫^N)}{kT}\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\frac{\varphi _0(r_{ij})}{kT}\right)\underset{i=1}{\overset{N}{}}\left\{_{\mathrm{}}^{\mathrm{}}ds_i\mathrm{exp}\left(\frac{\alpha }{kT}s_i^2\underset{j=1ji}{\overset{N}{}}\frac{\varphi _1(r_{ij})}{kT}s_i\right)\right\}},`$ where to obtain the second line we expressed the integral over the $`s`$ variables as a product of integrations over each one and then grouped all the terms which depend on each $`s_i`$ together. Each integration over an $`s`$ variable is independent and can be done easily $`Z_N`$ $`=`$ $`{\displaystyle d𝐫^N\mathrm{exp}\left(\frac{U(𝐫^N)}{kT}\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\frac{\varphi _0(r_{ij})}{kT}\right)\underset{i=1}{\overset{N}{}}\left\{\left(\frac{\pi kT}{\alpha }\right)^{1/2}\mathrm{exp}\left(\frac{1}{4\alpha kT}\left[\underset{j=1ji}{\overset{N}{}}\varphi _1(r_{ij})\right]^2\right)\right\}}`$ (14) $`=`$ $`\left({\displaystyle \frac{\pi kT}{\alpha }}\right)^{N/2}{\displaystyle d𝐫^N\mathrm{exp}\left(\frac{U(𝐫^N)}{kT}\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\frac{\varphi _0(r_{ij})}{kT}\right)\underset{i=1}{\overset{N}{}}\left\{\mathrm{exp}\left(\frac{1}{4\alpha kT}\underset{j=1ji}{\overset{N}{}}\underset{k=1ki}{\overset{N}{}}\varphi _1(r_{ij})\varphi _1(r_{ik})\right)\right\}}`$ $`=`$ $`\left({\displaystyle \frac{\pi kT}{\alpha }}\right)^{N/2}{\displaystyle d𝐫^N\mathrm{exp}\left(\frac{U(𝐫^N)}{kT}\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\frac{\varphi _0(r_{ij})}{kT}+\frac{1}{4\alpha kT}\underset{i,j,k=1jiki}{\overset{N}{}}\varphi _1(r_{ij})\varphi _1(r_{ik})\right)},`$ where to obtain the second from the first line we expressed the square of the sum as a double sum. To obtain the third line we converted the product of exponentials to the exponential of a sum. The factor in front of the integration is of course very familiar: it is just the partition function of $`N`$ independent simple harmonic oscillators. It does not depend on density and so has no effect on the phase behaviour. Note that in the triple sum although neither $`j`$ nor $`k`$ can be equal to $`i`$, $`j`$ can be equal to $`k`$. The restrictions on $`j`$ and $`k`$ derive from the fact that a particle cannot interact with itself, which would correspond to $`j,k=i`$ but as both $`j`$ and $`k`$ in the triple sum come from the square of a single sum they are in effect from the same interaction and therefore are allowed to be equal. We can rewrite Eq. (14) by extracting the $`j=k`$ terms from the triple sum, $$Z_N=\left(\frac{\pi kT}{\alpha }\right)^{N/2}d𝐫^N\mathrm{exp}\left(\frac{U(𝐫^N)}{kT}\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\left[\frac{\varphi _0(r_{ij})}{kT}\frac{\varphi _1(r_{ij})^2}{2\alpha kT}\right]+\frac{1}{4\alpha kT}\underset{i,j,k=1}{\overset{N^{}}{}}\varphi _1(r_{ij})\varphi _1(r_{ik})\right),$$ where the triple sum has a dash to indicate that terms in which any of the three subscripts are the same are excluded. The configurational integral, Eq. (3), is an integral only over the positions of the $`N`$ particles. Neglecting the irrelevant (for the phase behaviour) prefactor it is nothing but the configurational integral of the Hamiltonian, $`H_{eff}`$, $$H_{eff}(𝐫^N)=U(𝐫^N)+\frac{1}{2}\underset{i,j=1}{\overset{N^{}}{}}\left[\varphi _0(r_{ij})\frac{\varphi _1(r_{ij})^2}{2\alpha }\right]\frac{1}{6}\underset{i,j,k=1}{\overset{N^{}}{}}\frac{\left[\varphi _1(r_{ij})\varphi _1(r_{ik})+\varphi _1(r_{ij})\varphi _1(r_{jk})+\varphi _1(r_{ik})\varphi _1(r_{jk})\right]}{2\alpha },$$ where we have rewritten the summand of the triple sum to make it symmetric with respect to the three indices. The phase behaviour of our model particles will be identical to that of structureless, spherically symmetric particles with interactions described by the Hamiltonian $`H_{eff}`$ of Eq. (3). The Hamiltonian is that of a triplet or 3-body potential; the first two sums are conventional sums over pair potentials but the last sum is over a 3-body potential. We started with a pair potential which depended on the $`s`$ variables as well as the positions, integration over the $`s`$ variables resulted in a potential which depends only on the positions but is no longer a simple pair potential. The effective Hamiltonian, $`H_{eff}`$, can lead to behaviour qualitatively different from that of a Hamiltonian which is a simple sum over a pair potential. To see this we will compare the interaction between a pair of particles to that between a triplet of particles. For a pair of particles, 1 and 2, a distance $`r`$ apart, the interaction is from Eq. (3), $$H_{eff}(𝐫^2)=U(𝐫^2)+\varphi _0(r)\frac{1}{2}\frac{\varphi _1(r)^2}{\alpha }.$$ (15) Note that the contribution due to the fluctuations, the part inversely proportional to $`\alpha `$ is always negative regardless of the sign of $`\varphi _1`$, fluctuations always produce an effective attraction. Now, consider 3 particles, at the corners of an equilateral triangle of side $`r`$ for simplicity. We denote the set of 3 coordinates marking the corners of an equilateral triangle by $`𝐫_e^3`$. The interaction is, from Eq. (3), $$H_{eff}(𝐫_e^3)=U(𝐫_e^3)+3\varphi _0(r)3\frac{\varphi _1(r)^2}{\alpha }.$$ (16) If we compare this with the interaction of a pair, Eq. (15), we see that the ratio of the $`\varphi _1^2/\alpha `$ to $`\varphi _0`$ term has doubled. For a pairwise additive potential, the interaction free energy would be simply three times Eq. (15): $`3\varphi _0(3/2)\varphi _1^2/\alpha `$. The relative contribution from the fluctuations in the $`s`$ variables has doubled. This contribution is always attractive so the free energy of attraction is always more negative than for a pairwise additive potential. This is a general result for a weakly fluctuating scalar variable. Our result for the interaction between three particles may be compared with the more microscopic work on the interaction between three particles of Löwen and coworkers. They considered triplet interactions between charged colloidal particles and between star polymers . When $`\varphi _0(r)>0`$ and $`\varphi _0(r)<\varphi _1(r)^2/\alpha <2\varphi _0(r)`$, the interaction minus the part from $`U`$ is repulsive, i.e., greater then zero, for a pair $`r`$ apart but attractive for a triplet at the corners of an equilateral triangle of side $`r`$. If $`U`$ is some short-range repulsion, e.g., a hard sphere repulsion, which is zero at some sufficiently large value of $`r`$, then it is possible for a pair of our model colloidal particles to repel each other, but for a triplet of them to attract. For other arrangements of 3 particles the sign of the interaction free energy will depend on the values of the 3 separations of the centres of the particles and on the distance dependence of $`\varphi _0`$ and $`\varphi _1`$. However, this should not obscure the basic fact that the interaction free energy of 3 particles can be negative when the interaction of 2 particles is positive. For 4 particles at the corners of a tetrahedron the interaction minus the part from $`U`$ is $`6\varphi _09\varphi _1^2/\alpha `$ Again, it is possible to obtain negative interaction free energies even when the interaction free energy of a pair is always positive. Indeed it is possible to obtain negative free energies even when they are positive for three particles at the corners of an equilateral triangle. ## 4 Mean-field free energy In order to able to keep the theory simple we now specialise to an interaction between the internal variables which is both weak and long-ranged. We also set $`U`$ to be the hard-sphere interaction. Weak in the sense that $`\|\varphi _0(r)\|,\|\varphi _1(r)\varphi _1(r^{})/\alpha \|kT`$ for all values of $`r`$, $`r^{}`$, and long-ranged as both $`\varphi _0(r)`$ and $`\varphi _1(r)`$ decay to zero over a characteristic length scale which is much longer than the hard-sphere diameter $`\sigma `$. So the interaction Hamiltonian, Eq. (3), which is a function of the positions only, consists of a hard-sphere interaction, $`U`$, and a weak long-ranged interaction which has two parts: a pair potential and a triplet potential. Weak but long-range many-body interactions have been considered by the author in Ref. . We view the long-range part of the interactions as a perturbation . The free energy is expressed as that of hard spheres, $`A_{hs}`$, plus the difference in free energy between a system of our particles and that of hard spheres, $`\mathrm{\Delta }A`$, $$A=A_{hs}+\mathrm{\Delta }A.$$ (17) For the hard sphere free energy we use the expression of Carnahan and Starling for the fluid phase and the fit to simulation data of Hall for the solid phase. We approximate $`\mathrm{\Delta }A`$ by an average of the perturbing part of the Hamiltonian, $`H_{eff}U`$, $$\mathrm{\Delta }A=\frac{d𝐫^N\left(H_{eff}(𝐫^N)U(𝐫^N)\right)\mathrm{exp}\left(U(𝐫^N)/kT\right)}{Z_{hs}},$$ (18) where $`Z_{hs}`$ is the configurational integral for $`N`$ hard spheres. Substituting Eq. (3) into Eq. (18), $$\mathrm{\Delta }A=Z_{hs}^1d𝐫^N\mathrm{exp}\left(U(𝐫^N)/kT\right)\left\{\frac{1}{2}N(N1)\left(\varphi _0(r_{12})\frac{\varphi _1(r_{12})^2}{2\alpha }\right)\frac{1}{6}(N1)(N2)\frac{3\varphi _1(r_{12})\varphi _1(r_{13})}{2\alpha }\right\}.$$ The $`n`$-particle density of hard spheres, $`\rho _{hs}^{(n)}`$, and their $`n`$-particle distribution function, $`g_{hs}^{(n)}`$, are defined as $`\rho _{hs}^{(n)}(𝐫^n)`$ $`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\rho _{hs}^{(1)}(𝐫_i)\right)g_{hs}^{(n)}(𝐫^n)`$ (19) $`=`$ $`{\displaystyle \frac{N!}{(Nn)!}}{\displaystyle \frac{\mathrm{exp}\left[\beta U(𝐫^N)\right]d𝐫^{Nn}}{Z_{hs}}}.`$ The 1-particle density, $`\rho _{hs}^{(1)}(𝐫)`$, is not assumed to be uniform so that the theory applies to crystalline as well as fluid phases. Using, Eq. (19) in Eq. (4), we obtain $$\mathrm{\Delta }A=\frac{1}{2}d𝐫^2\rho _{hs}^{(2)}(𝐫_1,𝐫_2)\left(\varphi _0(r_{12})\frac{\varphi _1(r_{12})^2}{2\alpha }\right)\frac{1}{4\alpha }d𝐫^3\rho _{hs}^{(3)}(𝐫_1,𝐫_2,𝐫_3)\varphi _1(r_{12})\varphi _1(r_{13}).$$ As $`\varphi _0,\varphi _1`$ decay to zero only when the separations of the particles are much larger than $`\sigma `$, the integrals in Eq. (4) are dominated by configurations when the spheres are far apart, i.e., where the pair separations are much larger than $`\sigma `$. In a fluid the one particle density is a constant, $`\rho ^{(1)}=\rho `$, and at separations large with respect to $`\sigma `$ the distribution function is close to one, $`g_{hs}^{(n)}1`$. Thus in the fluid phase the integrands of Eq. (4) can be simply approximated by $`\rho ^2(\varphi _0(r_{12})\varphi _1(r_{12})^2/2\alpha )`$ and $`\rho ^3\varphi _1(r_{12})\varphi _1(r_{13})`$. In a crystalline phase, although there are long-range correlations in $`\rho _{hs}^{(n)}`$ the one particle density, $`\rho _{hs}^{(1)}`$, averaged over a unit cell is just $`\rho `$. The attractive interaction between particles has a range much larger than the lattice constant of the lattice (a little larger than $`\sigma `$) and so as $`\varphi _0`$ and $`\varphi _1`$ vary little across a unit cell we can regard $`\rho _{hs}^{(1)}(𝐫)`$ as approximately constant at its average value, $`\rho `$. Similarly, as we change any one of the $`n`$ position vectors upon which the $`n`$-body distribution function, $`\rho _{hs}^{(n)}`$, depends the density oscillates rapidly over each unit cell but averages to $`\rho `$. So, in the crystalline phase as well as in the fluid phase we approximate $`\rho _{hs}^{(n)}`$ by $`\rho ^n`$. Then the integrands are the same as in a fluid phase. We have $$\mathrm{\Delta }A=\frac{1}{2}V\rho ^2_\sigma ^{\mathrm{}}d𝐫\left(\varphi _0(r)\frac{\varphi _1(r)^2}{2\alpha }\right)\frac{1}{4\alpha }V\rho ^3\left(_\sigma ^{\mathrm{}}d𝐫\varphi _1(r)\right)^2.$$ (20) Now the term coming from the square of $`\varphi _1`$ is negligible if $`\varphi _1`$ is long-ranged. To see this consider an explicit functional form for $`\varphi _1`$. We choose a Kac potential, $$\varphi _1(r)=ϵ\gamma ^3\mathrm{exp}(\gamma r/\sigma ).$$ (21) The range of this function is $`\sigma /\gamma `$, i.e., $`\gamma ^1`$ is a range in units of $`\sigma `$. Now the integral of Eq. (21) over 3-dimensional space is $`𝒪(ϵ\sigma ^3)`$; essentially the integrand is of order $`ϵ\gamma ^3`$ within a volume of order $`(\sigma /\gamma )^3`$, and negligible outside of this volume. However, the integral of the square of Eq. (21) over 3-dimensional space is $`𝒪(ϵ\sigma ^3\gamma ^3)`$. For a long-range $`\varphi _1`$, the inverse range $`\gamma 1`$ and so the integral of the square is negligible. We chose a specific functional form for $`\varphi _1`$ simply for clarity, our conclusion that the integral of the square is negligible holds for any long-range slowly decaying function. So, in Eq. (20), we neglect the $`\varphi _1^2`$ term, integrate over the remaining terms in Eq. (20) and obtain $$\frac{\mathrm{\Delta }A}{N}=\frac{1}{2}\rho \nu _0\frac{1}{4}\frac{(\rho \nu _1)^2}{\alpha },$$ (22) where $`\nu _0`$ and $`\nu _1`$ are the integrals $$_\sigma ^{\mathrm{}}d𝐫\varphi _i(r)=\nu _i,i=0,1.$$ (23) The pressure $`p`$ can be easily derived from the free energy, Eq. (22), $$p=p_{hs}+\frac{1}{2}\rho ^2\nu _0\frac{1}{2}\frac{\rho ^3\nu _1^2}{\alpha },$$ (24) where $`p_{hs}`$ is the pressure of hard spheres. The contribution from the fluctuations in the $`s`$ variables to the pressure is cubic in the density as it must, it is equivalent to a long-range 3-body attraction . The chemical potential $`\mu `$ is then simply given by $`\mu =A/N+p/\rho `$. The pressure and chemical potential as functions of the temperature and density allow us to calculate phase diagrams. ## 5 Results The free energy, Eq. (22), depends on a single parameter, the dimensionless ratio $`R=\nu _0\alpha \sigma ^3/\nu _1^2`$. It ranges from $`\mathrm{}`$ to $`+\mathrm{}`$; the sign of $`R`$ is determined by the sign of $`\nu _0`$ as $`\alpha `$ must be positive. The larger the magnitude of $`R`$ the more dominant is the part of the interaction potential which does not depend on $`s`$ and so the closer the interaction is to a pairwise additive potential. In the $`R\mathrm{}`$ limit the interaction is a simple pairwise additive attractive potential. Then the free energy is just a modification of the free energy of van der Waals . The modification being the replacement of his approximate free energy of hard spheres by that of Carnahan and Starling in the fluid, and that of Hall in the crystal. In the $`R+\mathrm{}`$ limit the potential is again pairwise additive but it is repulsive. A long-range repulsion tends to cause not bulk vapour-liquid separation but microphase separation, see Refs. . The $`R\pm \mathrm{}`$ limits result from either the interaction between particles not depending on the value of the $`s`$ variables, $`\nu _10`$, or the fluctuations of the $`s`$ variables tending to zero, $`\alpha \mathrm{}`$. When $`R`$ is finite then the interactions depend on the fluctuations in $`s`$ and are no longer pairwise additive. For $`R<0`$ the interactions are attractive even in the absence of fluctuations in $`s`$; the fluctuations in $`s`$ merely make the attractions stronger and not pairwise additive. For $`R>0`$ the pairwise interactions are repulsive while the 3-body interactions from the fluctuations are, as always, attractive. The part of the free energy due to fluctuations, the last term in Eq. (22), is one power higher in the density than the part independent of the fluctuations, the second in Eq. (22). Thus, if $`\nu _0>0`$ the interactions are always repulsive at sufficiently low densities but become attractive, i.e., contribute a negative amount to the free energy of the system, at a density $`2R/\sigma ^3`$. We now discuss some example phase diagrams. The diagrams are in the density-temperature plane. We use the reduced density $`\eta =(\pi /6)\rho \sigma ^3`$, which is equal to the fraction of the volume occupied by the hard cores, and a reduced temperature $`T^{}`$ which is either $`kT\sigma ^3/\|\nu _0\|`$ or $`kT\alpha \sigma ^6/\nu _1^2`$. For reference we plot the familiar phase diagram of hard spheres with a long-range attraction, often called the van-der-Waals fluid ; the $`R\mathrm{}`$ limit. This has a free energy given by Eqs. (17) and (22) with $`\nu _1=0`$ and $`\nu _0<0`$. The first condition means that the internal variables on different particles do not couple and so the attraction is a simple pair potential and the second condition makes the long-range interaction attractive so there is a vapour-liquid transition, i.e., phase separation into two fluid phases of different densities. The phase diagram is shown in Fig. 1; there is a large temperature range over which there is stable vapour-liquid coexistence. This is of course not new, in the $`\nu _1=0`$ limit our model reduces to the most basic model of particles which form a liquid. In Fig. 2 we show the phase diagram of particles with $`R=0`$. The long-range interaction is purely proportional to the internal variable $`s`$, i.e., $`\nu _0=0`$. There is vapour-liquid coexistence over a range of temperatures but this range is much smaller than for the simple pair potential of Fig. 1. Also, the density at the critical point is higher; it is almost double that in Fig. 1. The interaction between $`s`$ variables is effectively a long-range 3-body attraction, i.e., an attraction which is van-der-Waals like except for the fact that it is between triplets not pair of particles . The phase diagram, Fig. 2, differs from Fig. 2 of Ref. only in that the temperature scale is different. In Fig. 3 we show the phase diagram of particles with $`R=2`$. There is no vapour-liquid coexistence at equilibrium, only a fluid-crystal coexistence region which broadens dramatically at low temperatures. The liquid phase has disappeared. For other examples of liquid phases disappearing see Ref. ; the most studied system in which the liquid disappears is that of particles with a short-range, pairwise additive attraction . Within the fluid-crystal coexistence region the pressure and chemical potential have van-der-Waals loops so our bulk free energy predicts vapour-liquid coexistence within the fluid-crystal coexistence region. Note that the density at the critical point is even higher than in Fig. 2. We have only considered fluid phases and crystalline phases. However, sufficiently strong long-range repulsions can transform a bulk phase separation into microphase separation . Essentially, microphase separation occurs when the interaction between particles is repulsive at the largest separations at which they interact separately and the repulsion is sufficiently strong and long ranged. Thus our present theory, which neglects the possibility of microphase separation, should not be used if the long-range interaction is predominantly repulsive, i.e., if $`\nu _0>0`$, and either $`\varphi _0`$ is longer ranged than $`\varphi _1`$, or $`R1`$. ### 5.1 The structure of the fluid phase We have shown that the phase behaviour of our particles can be very different from that of a simple pair potential but what about the structure? Of course the pair distribution function or equivalently the structure factor depends on the functional form of the interactions, i.e., on the precise forms of the functions $`\varphi _0(r)`$ and $`\varphi _1(r)`$. However, we have only considered long-range interactions and there the bulk phase behaviour is insensitive to the precise details of $`\varphi _0(r)`$ and $`\varphi _1(r)`$, it only depends on their integrals $`\nu _0`$ and $`\nu _1`$. Also, long-range interactions only affect the structure factor $`S(q)`$ at small wavevectors $`q`$; small meaning of order the reciprocal of their range or less. Thus, we will consider only the zero wavevector limit, $`S(0)`$, of the structure factor. This is simply related to a thermodynamic quantity, the isothermal compressibility $`\chi _T`$, by $$S(0)=\rho kT\chi _T,$$ (25) where $$\chi _T^1=\rho \left(\frac{p}{\rho }\right)_T.$$ (26) Using Eq. (24) for the pressure $$\chi _T^1=\rho \left(\frac{p_{hs}}{\rho }\right)_T+\rho ^2\nu _0\frac{3}{2}\frac{\rho ^3\nu _1^2}{\alpha }.$$ (27) Eqs. (25) and (27) give us the zero wavevector structure factor of our particles. Below we will consider the density dependence of $`S(0)`$. It will therefore be useful to obtain $`S(0)`$ as a density expansion. To do this we start by inserting the virial expansion for the pressure into the definition of the isothermal compressibility $`\chi _T`$, Eq. (26). Inserting this expansion into Eq. (25) for $`S(0)`$ yields $$S(0)=12B_2\rho +𝒪(\rho ^2),$$ (28) the initial slope of $`S(0)`$ is equal to minus twice the second virial coefficient $`B_2`$. For our particles $$B_2=\frac{2}{3}\pi \sigma ^3\frac{1}{2}\nu _0,$$ (29) where the first term on the right hand side is the second virial coefficient of hard spheres. Note that the fluctuations in the $`s`$ variables do not contribute to the second virial coefficient. This is only the case in the long-range interaction limit. In section 3 and in Eq. (20) we found that in general the fluctuations did contribute to the interaction between pairs and hence to the second virial coefficient but that this contribution was weaker than its contribution to the interaction between larger numbers of particles and to the third virial coefficient. The term in the free energy per unit volume proportional to $`\nu _1^2/\alpha `$ which comes from the fluctuations varies with density as $`\rho ^3`$ and so it contributes (only) to the third virial coefficient. In Fig. 4 we have plotted the zero wavevector structure factor for a van-der-Waals fluid, hard spheres plus a long-range pairwise additive attraction (the solid curve), and for an attraction which is proportional to $`s`$ (the dashed curve). They are both at temperatures just above the critical temperature. The peaks in $`S(0)`$ are due to the nearby critical points, where $`S(0)`$ diverges of course. The striking difference between the curves is their opposite slopes at low densities. The limiting slope at vanishing densities is equal to $`2B_2`$, Eq. (28). For a van-der-Waals fluid the second virial coefficient is negative at the critical temperature so the slope of $`S(0)`$ is positive until above the critical density. However, interactions between fluctuations do not contribute to the second virial coefficient and so it is equal to that of hard spheres which is positive. Thus the slope of $`S(0)`$ is negative at low densities and $`S(0)`$ goes through a minimum below the critical density. ## 6 Conclusion Particles which have a fluctuating internal state interact in a way which is different from structureless particles. The interaction is intrinsically many body and cannot be described using a pair potential. It is even possible for pairs of particles with a fluctuating state to repel but larger clusters of particles to attract each other. We have studied a rather generic model of a particle with a fluctuating state. Remarkably, because of the very simple form of our Hamiltonian’s dependence on the variables which describe the states of the particles, see Eq. (10), we have been able to integrate over these variables exactly and analytically. The result is an interaction between the particles which has both pair and triplet interactions, Eq. (3). Our model only included the leading order terms in Taylor expansions for the free energy associated with the internal states of the particles. It is possible to include higher order terms although then the $`s`$ variables will have to integrated over approximately or numerically, alternatively computer simulation could be used. In any case the effective potential will then contain not only pair and triplet terms but terms of all orders, 4-body, 5-body, etc.. It is perhaps of interest to contrast the particles with mesoscopic fluctuations studied here to the rods with microscopic fluctuations studied by Ha and Liu , see also Refs. . They studied the interactions between parallel rods and allowed the charge density along each rod to fluctuate; if a rod is along the $`z`$-axis then we can define a 1-dimensional charge density $`\rho (z)`$, which will fluctuate. The Coulomb interaction between charges will then couple the fluctuations in $`\rho (z)`$ in nearby rods. Ha and Liu found that not only can this coupling generate an attraction but that it is strongly nonpairwise additive. In contrast with their model our is rather general and simpler. The results for our rather non-specific model show clearly that these findings of Ha and Liu are rather generic. The disadvantage of our model with respect to that of Ha and Liu is that to relate it to experiment requires determining or at least guessing the appropriate values of the parameters of our phenomenological interaction. We hope that our model will be useful for determining generic differences in behaviour between rigid particles interacting via a pair potential and particles with a fluctuating state. For example, although all our calculations have been for a single component they are easily generalised to mixtures. When this is done, it is straightforward to see that, if the fluctuations between the internal states of particles of different components are only weakly coupled then this will tend to drive phase separation of these components. However, differences between the inverse susceptibility, $`\alpha `$, cannot drive phase separation. At least within our theory which includes only leading order terms. It is a pleasure to acknowledge discussions with B.-Y. Ha, Y. Levin, R. van Roij and P. B. Warren. ## Appendix Rather generally, when we are evaluating a configurational integral, with the aim of calculating the free energy, we can split the variables over which we are integrating into two sets. Then integrating over one set is always possible in principle, and it results in an effective Hamiltonian which is a function only of the remaining set of variables. However, this effective Hamiltonian is not in general pairwise additive, simply because there is no reason for it to be. If we denote the set of all variables by $`v^{N+M}`$, and the two subsets by $`v^N`$ and $`v^M`$ then the configurational integral is $$Z=dv^{N+M}\mathrm{exp}\left(H(v^{N+M})/kT\right)=dv^N\mathrm{exp}\left(H^{}(v^N)/kT\right),$$ (30) where $$\mathrm{exp}\left(H^{}(v^N)/kT\right)=dv^M\mathrm{exp}\left(H(v^{N+M})/kT\right).$$ (31) $`H^{}`$ will not in general be expressible as a sum over a pair potential even if $`H`$ can be. Physically, this is most often useful when the two sets of coordinates, $`v^N`$ and $`v^M`$, are very different. We consider a model in which the set of coordinates $`v^M`$ is coarse grained , i.e., the integration over $`v^M`$ is left as a function of a set of coarse-grained variables, so instead of Eq. (31) we have, $$\mathrm{exp}(H^{}(v^N,c^L))/kT)=\mathrm{d}v^M\mathrm{exp}(H(v^{N+M})/kT)_{a=1}^L\delta (c_a(v^M)c_a),$$ (32) where $`c^L`$ is a set of $`LM`$ coarse-grained variables. $`c_a`$ is the $`a`$th coarse-grained variable and $`c_a(v^M)`$ is the definition of the $`a`$th coarse-grained variable in terms of the microscopic coordinates, $`v^M`$. The configurational integral is then $$Z=dv^Ndc^L\mathrm{exp}\left(H^{}(v^N,c^L)/kT\right).$$ (33) Our model has a configurational integral, Eq. (11), of this form, with $`L=N`$ as there is one coarse-grained variable per particle. The $`v^N`$ and $`c^N`$ coordinates are the positions and $`s`$ variables of the particles.
warning/0002/math-ph0002013.html	ar5iv	text	# Theorem 1.1 MINIMAL ESCAPE VELOCITIES W. Hunziker Institut für Theoretische Physik, ETH Zürich hunziker@itp.phys.ethz.ch I.M. Sigal Departement of Mathematics, University of Toronto sigal@math.toronto.edu A. Soffer Departement of Mathematics, Rutgers University soffer@math.rutgers.edu (March 1997) Abstract. We give a new derivation of the minimal velocity estimates \[SiSo1\] for unitary evolutions. Let $`H`$ and $`A`$ be selfadjoint operators on a Hilbert space $``$. The starting point is Mourre’s inequality $`i[H,A]\theta >0`$, which is supposed to hold in form sense on the spectral subspace $`_\mathrm{\Delta }`$ of $`H`$ for some interval $`\mathrm{\Delta }R`$. The second assumption is that the multiple commutators $`ad_A^{(k)}(H)`$ are well-behaved for $`k=1\mathrm{}n(n2)`$ . Then we show that, for a dense set of $`\psi `$’s in $`_\mathrm{\Delta }`$ and all $`m<n1,`$ $`\psi _t=exp(iHt)`$ is contained in the spectral subspace $`A\theta t`$ as $`t\mathrm{}`$, up to an error of order $`t^m`$ in norm. We apply this general result to the case where $`H`$ is a Schrödinger operator on $`R^n`$ and $`A`$ the dilation generator, proving that $`\psi _t(x)`$ is asymptotically supported in the set $`\|x\|t\sqrt{\theta }`$ up to an error of order $`t^m`$ in norm. 1. INTRODUCTION Before posing the problem in abstract form we describe it in its original concrete setting. Consider the Schrödinger equation $$i_t\psi _t=H\psi _t;H=\frac{1}{2}p^2+V(x)\mathrm{on}L^2(R^n)$$ $`(1.1)`$ for a particle in $`R^n`$ under the influence of a potential $`V(x)`$. We are interested in the long–time behavior of orbits $`t\psi _t`$ in the continuous spectral subspace $`_c`$ of $`H`$. Under mild conditions on $`V`$, $`H`$ is selfadjoint and $$p^2_\psi \mathrm{const}.H+c_\psi $$ $`(1.2)`$ for some constant $`cR`$. Here $`\mathrm{\Phi }_\psi =(\psi ,\mathrm{\Phi }\psi )`$ denotes the expectation value of an observable $`\mathrm{\Phi }`$ in the state $`\psi `$. By Ruelle’s theorem (\[Rue\], see also \[CFKS\], \[HuSi1\]) any orbit $`\psi _t`$ in $`_c`$ is escaping in a mean ergodic sense : $$\underset{T\mathrm{}}{lim}\frac{1}{T}_0^T𝑑t_{\|x\|R}𝑑x\|\psi _t(x)\|^2=0$$ $`(1.3)`$ for any finite $`R`$. The question is : how fast ? The answer will of course depend on the initial state $`\psi `$. The simplest example is the free particle ($`V=0`$): if the Fourier transform $`\widehat{\psi }`$ of $`\psi `$ is smooth and supported outside some ball of radius $`v>0`$, then a standard asymptotic expansion gives the result $$_{\|x\|vt}𝑑x\|\psi _t(x)\|^2=O(t^m)(t\mathrm{})$$ $`(1.4)`$ for any $`m`$. In this sense the orbits $`\psi _t`$ of the given type are said to have a minimal escape velocity $`v`$ given by the support of $`\widehat{\psi }`$. To obtain a similar result for $`V0`$ we study the long–time behavior of the expectation values $`A_t=(\psi _t,A\psi _t)`$ for suitable observables $`A`$, which evolve according to $$_tA_t=i[H,A]_t.$$ $`(1.5)`$ Mourre’s very fruitful idea \[Mou\] was to find observables $`A`$ such that the commutator in (1.5) is conditionally positive , in the sense that $$E_\mathrm{\Delta }i[H,A]E_\mathrm{\Delta }\theta E_\mathrm{\Delta }(\theta >0)$$ $`(1.6)`$ for some interval $`\mathrm{\Delta }R`$, where $`E_\mathrm{\Delta }`$ is the corresponding spectral projection of $`H`$. This implies that $$A_t\theta t+O(1)\mathrm{}(t\mathrm{})$$ $`(1.7)`$ for orbits $`\psi _t`$ in the spectral subspace $`_\mathrm{\Delta }=\mathrm{Ran}E_\mathrm{\Delta }`$. Evidently $`A`$ must be unbounded, so that domain questions arise. Also $`_\mathrm{\Delta }`$ must be a subspace of $`_c`$ since $`A_t`$ is constant and $`i[H,A]_\psi =0`$ for any eigenvector $`\psi `$ of $`H`$. (1.6) is a special case of a more general inequality due to Mourre, which was first proven for Schrödinger operators (including $`N`$–body systems), where $`A`$ was taken as the dilation generator : $$A=\frac{1}{2}(px+xp);i[H,A]=p^2xV(x)$$ $`(1.8)`$ (\[Mou\], \[PSS\], see also \[CFKS\], \[HuSi1\]). In this case the intervals $`\mathrm{\Delta }`$ for which (1.6) holds fill the continous spectrum of $`H`$, in the sense that the corresponding subspaces $`_\mathrm{\Delta }`$ span $`_c`$. Moreover, since $$A=i[H,\frac{1}{2}x^2]$$ $`(1.9)`$ is itself a commutator, (1.5) can be written as $`_{t}^{}{}_{}{}^{2}x^2_t2\theta `$ for orbits $`\psi _t`$ in $`_\mathrm{\Delta }`$, which implies that $$x^2_t\theta t^2+O(t)(t\mathrm{}).$$ $`(1.10)`$ This is of course weaker than (1.4) : it only says that the mean value of $`x^2`$ for the probability distribution $`\|\psi _t(x)\|^2`$ diverges like $`\theta t^2`$, whereas we want to prove that the support of this distribution is asymptotically contained in $`\|x\|t\sqrt{\theta }`$ as $`t\mathrm{}`$ . The first step is to derive a corresponding result for the spectral support of $`\psi _t`$ with respect to $`A`$. We state this result in abstract form for a pair $`(H,A)`$ of selfdjoint operators on a Hilbert space $``$. ###### Theorem 1.1 Suppose that $`ad_A^{(k)}(f(H))`$ is bounded for any $`fC_0^{\mathrm{}}(R)`$ and $`k=1\mathrm{}n`$, $`n2`$, and that the Mourre inequality (1.6) holds for some open interval $`\mathrm{\Delta }R`$. Let $`\chi ^\pm `$ be the characteristic function of $`R^\pm `$. Then $$\chi ^{}(Aa\vartheta t)e^{iHt}g(H)\chi ^+(Aa)\mathrm{const}.t^m$$ $`(1.11)`$ for any $`gC_0^{\mathrm{}}(\mathrm{\Delta })`$, any $`\vartheta `$ in $`0<\vartheta <\theta `$ and any $`m<n1`$, uniformly in $`aR`$ . Since $`gC_0^{\mathrm{}}(\mathrm{\Delta })`$ and $`aR`$ are arbitrary, the vectors of the form $$\psi =g(H)\chi ^+(Aa)\phi ;\phi $$ $`(1.12)`$ are dense in $`_\mathrm{\Delta }`$. (1.11) says that, for any such $`\psi `$, $`\psi _t=\mathrm{exp}(iHt)\psi `$ has spectral support in $`[t\vartheta ,+\mathrm{})`$ with respect to $`A`$ , up to a remainder of order $`t^m`$ in norm. Remarks Commutators. The hypothesis $`ad_A^{(k)}(f(H))L(H)`$ may be replaced by conditions on $`ad_A^{(k)}(H)`$ which are more subtle to formulate since the operators $`A`$ and $`H`$ are generally unbounded (see e.g. \[ABG\], \[JMP\]). For the special case (1.8) this is further discussed below. Resolvent smoothness and local decay. We indicate briefly how minimal velocity estimates are related to resolvent smoothness \[JMP\] and to local decay \[PSS\]. Let $`\rho (A)=(1+A^2)^{1/2}`$. Setting $`a=\theta t/2`$ and using that $$\rho (A)^\alpha =\rho (A)^\alpha \chi ^\pm (A\pm ct)+O(t^\alpha ),$$ we obtain from (1.11) $$\rho (A)^\alpha e^{iHt}g(H)\rho (A)^\alpha \mathrm{const}.(1+t)^{\mathrm{min}(\alpha ,m)}.$$ For $`\alpha ,m>1`$ this is integrable over $`\mathrm{}<t<+\mathrm{}`$ , which (by Fourier transform) leads to the resolvent estimate $$\underset{zR}{sup}\rho (A)^\alpha (zH)^1g(H)\rho (A)^\alpha \mathrm{}.$$ $`(1.13)`$ Similar estimates for the derivatives with respect to $`z`$ (higher powers) of the resolvent $`(zH)^1`$ are obtained using correspondingly higher values of $`\alpha ,m`$ (resolvent smoothness). Replacing $`g`$ by $`g^2`$ in (1.13) it also follows that the operator $`\rho (A)^\alpha `$ is $`H`$–smooth for $`\alpha >1`$ and therefore (\[RSIV\] Theorem XIII.25 and corollary) $$_{\mathrm{}}^+\mathrm{}𝑑t\rho (A)^\alpha e^{iHt}\psi ^2\mathrm{const}.\psi ^2\psi _\mathrm{\Delta }^{},$$ $`(1.14)`$ where $`\mathrm{\Delta }^{}`$ is any fixed compact subset of $`\mathrm{\Delta }`$ (local decay). By an independent argument of Mourre (given in \[PSS\]) the estimate (1.13) and therefore (1.14) can be improved from $`\alpha >1`$ to $`\alpha >1/2`$. Our second result is an application of Theorem 1.1 to the Schrödinger equation (1.1). Here $`A`$ is given by (1.8), and $$i^kad_A^{(k)}(H)=2^{k1}p^2+(x)^kV(x).$$ $`(1.15)`$ A simple way to satisfy the hypothesis of Theorem 1.1 in this case is to assume that $`V`$ has relative bound less than 1 with respect to $`\frac{1}{2}p^2`$, and that the distributions $`(x)^kV(x)`$ are locally $`L^2`$ and (as multiplications operators) bounded relative to $`p^2`$ for $`k=1\mathrm{}n`$. Then the operators (1.15) are bounded relative to $`H`$ and it is straightforward to compute and to estimate the the norms of $`ad_A^{(k)}((zH)^1)`$ for $`\mathrm{Im}(z)0`$. Representing $`f(H)`$ by the resolvent $`(zH)^1`$ (e.g. using the Helffer-Sjöstrand formula (\[HeSj\], \[Dav\]) it then follows that $`ad_A^{(k)}(f(H))`$ is bounded for $`k=1\mathrm{}n`$. The result is that fast decay (large $`m`$) in (1.11) must be paid for by high smoothness of $`V(x)`$ for all $`x`$. This is unnatural, and in fact there is a better way to construct $`A`$ which requires only smoothness of $`V(x)`$ for arbitray large $`\|x\|`$. The idea is to replace $`x^2`$ in (1.9) by a smooth, convex function $`G(x)`$ which is equal to $`x^2`$ for large $`\|x\|`$ but constant in some abitrary large ball $`\|x\|R`$. Then $`A`$ changes to $$A=\frac{1}{2}(G(x)p+pG(x)),$$ and the Mourre inequality can be established as before. Then the operators $`ad_A^{(k)}(p^2)`$ remain second order in $`p`$ with bounded coefficients, while $$i^kad_A^{(k)}(V)=(G(x))^kV(x)$$ requires only derivatives of $`V(x)`$ in the region $`\|x\|>R`$. A more carful construction of $`G(x)`$ due to Graf \[Gra1\] is specially adapted to the $`N`$-body case, requiring only smoothness of the pair potentials for large separations (see \[Skib\], \[Gri\]). The following result can also be proven in this more general setting. ###### Theorem 1.2 Let $`H`$ and $`A`$ be given by (1.1) and (1.8). Then, under the hypothesis of Theorem 1.1 , $$\chi ^{}(x^22at\vartheta t^2)e^{iHt}g(H)\chi ^+(Aa)\mathrm{const}.t^m$$ $`(1.16)`$ for any $`\vartheta `$ in $`0<\vartheta <\theta `$, any $`m<n1`$ and any $`aR`$. We remark that for initial states of the form (1.13) this is equivalent to $$_{\|x\|vt}𝑑x\psi _t(x)^2\mathrm{const}.t^{2m}$$ $`(1.17)`$ for any $`v<\sqrt{\theta }`$. The proofs of these results are given in section 2. The main tool is the method of commutator expansions summarized in section 3. We conclude the introduction with some (not exhaustive) bibliographical notes. Minimal velocity estimates were first given by Sigal and Soffer \[SiSo1\] and then extended by Skibsted \[Ski\] and by Gérard and Sigal \[GeSi\], with applications to scattering theory (\[SiSo2\], \[Sig\], \[HeSk\]) and to the theory of resonances (\[GeSi\], \[SoWe\], \[Nier\]). Our derivation is similar in spirit to \[Ski\] and incorporates remarks by Froese and Loss \[FrLo\]. The related subject of resolvent smoothness and local decay is more fully treated by time–independent methods e.g. in \[PSS\], \[JMP\] and \[GIS\], where further references can be found. The generalization of Mourre’s theorem using the construction of Graf \[Gra1\] first appears in \[Ski\]. A simpler proof due to Graf \[Gra2\] is given in \[Gri\]. For other applications of Mourre’s inequality to wave equations and spectral geometry see e.g. \[DHS\], \[DBiPr\], \[FHP\]. 2. PROOFS The following lemma gives the basic estimate for a proof by bootstrap of Theorem 1.1. A smooth function $`f`$ on $`R`$ is said to be of order $`p`$ if for each $`k`$ $$\|f^{(k)}(x)\|\mathrm{const}.\|x\|^{pk}.$$ ###### Lemma 2.1 Suppose that $`A,H`$ and $`g`$ satisfy the hypothesis of Theorem 1.1. Let $`f`$ be a positive $`C^{\mathrm{}}`$-function on $`R`$ of order $`<4`$ with $`f^{}0`$ and $`f(x)=0`$ for $`x0`$. Let $`1s<\mathrm{}`$; $`aR`$, $`A_s=s^1(Aa)`$, and $`\epsilon 1`$. Then $$\begin{array}{cc}\hfill g(H)i[H,f(A_s)]g(H)& s^1\theta g(H)f^{}(A_s)g(H)\hfill \\ & +s^{(1+\epsilon )}g(H)f_1(A_s)g(H)\hfill \\ & +\mathrm{const}.s^{(2n1\epsilon )}g^2(H)\hfill \end{array}$$ $`(2.1)`$ uniformly in $`aR`$, where $`f_1`$ is a function on $`R`$ which again satisfies the hypothesis stated for $`f`$. Proof. In the commutator $`i[H,f(A_s)]`$ occurring in (2.1) we can replace $`H`$ by a bounded operator $$H_b=Hb(H)$$ $`(2.2)`$ where $`bC_0^{\mathrm{}}(R);b1`$ on $`\mathrm{supp}(g)`$. Then the commutators $$B_k=iad_A^{(k)}(H_b),k=1\mathrm{}n,$$ are bounded by hypothesis. As a first step we show that $$i[H_b,f(A_s)]=s^1(f^{}(A_s))^{1/2}B_1(f^{}(A_s))^{1/2}+\mathrm{remainder}.$$ $`(2.3)`$ Here and in the rest of the proof a remainder is defined as a quadratic form rem$`(s)`$ with an estimate $$\pm \mathrm{rem}(s)s^{(1+\epsilon )}f_1(A_s)+\mathrm{const}.s^{(2n1\epsilon )},$$ $`(2.4)`$ uniformly in $`a`$, where $`f_1`$ is a function on $`R`$ satisfying the hypothesis for $`f`$. Any such remainder clearly fits into (2.1) and needs no further consideration. To prove (2.3) we factorize $`f=F^2`$ and then expand the commutator $$\begin{array}{cc}\hfill i[H_b,f]& =i[H_b,F]F+Fi[H_b,F]\hfill \\ & =\underset{k=1}{\overset{n1}{}}\frac{1}{k!}s^k(F^{(k)}B_kF+FB_k^{}F^{(k)})\hfill \\ & +s^n(RF+FR^{})\hfill \end{array}$$ $`(2.5)`$ using (3.1). Since $`n2`$ and since $`F`$ is of order $`<2`$ it follows from (3.2) that $`R`$ is bounded uniformly in $`s`$ and $`a`$. Now we observe that all the terms in the expansion (2.5) except the leading term $`(k=1)`$ are remainders. In particular $$\|(\psi ,F^{(k)}B_kF\psi )\|B_kF^{(k)}\psi F\psi \mathrm{const}.(\psi ,f_1\psi ),$$ $`(2.6)`$ where $`f_1`$ is a common upper bound for $`F^2=f`$ and $`(F^{(k)})^2`$ which satisfies the hypothesis for $`f`$. The last term in (2.5) is estimated using the operator inequality $$\pm (P^{}Q+Q^{}P)P^{}P+Q^{}Q$$ for $`Q=s^{\frac{1}{2}(1+\epsilon )}F`$; $`P^{}=s^{n+\frac{1}{2}(1+\epsilon )}R`$, with the result $$\pm s^n(RF+FR^{})s^{(1+\epsilon )}f+s^{(2n1\epsilon )}R^2.$$ $`(2.7)`$ Therefore it remains to consider the leading term $`(k=1)`$ in (2.5), which is rewritten as $$s^1(F^{}B_1F+FB_1F^{})=s^1(v^2B_1u^2+u^2B_1v^2)$$ by factorizing $`F=u^2`$, $`F^{}=v^2`$. Since $`u`$ is of order $`<1`$ it follows from (3.2) that $`[B_1,u]\mathrm{const}.s^1`$ uniformly in $`a`$, and similarly for $`[B_1,v]`$. As in (2.6) this leads to the form estimate $$\begin{array}{cc}& s^1\|v^2B_1u^2uvB_1uv\|=s^1\|v[B_1,u]uv+v[v,B_1]u^2\|\hfill \\ & \mathrm{const}.s^2(v^2+u^2v^2+u^4)\mathrm{const}.s^2f_1(A_s),\hfill \end{array}$$ where $`f_1`$ shares the properties of $`f=u^4`$ (note that $`v`$ is of lower order than $`u`$). Since the same estimate holds with $`u`$ and $`v`$ interchanged, we conclude that $$s^1(F^{}B_1F+FB_1F^{})=2s^1uvB_1uv$$ plus a remainder. (2.3) now follows since $`f^{}=2FF^{}=2(uv)^2`$. To complete the proof of Lemma 2.1. we multiply (2.3) from both sides with $`g(H)=g(H)G(H)`$, where $`GC_0^{\mathrm{}}(\mathrm{\Delta })`$ is real and $`G1`$ on $`\mathrm{supp}(g)`$. We also adjust the function $`b`$ in (2.2) such that $`b1`$ on $`\mathrm{supp}(G)`$. Multiplying (1.6) by $`G(H)`$ from both sides we obtain the Mourre inequality : $$G(H)B_1G(H)=G(H)i[H,A]G(H)\theta G^2(H).$$ $`(2.8)`$ Abbreviating $`G(H)=G`$ and $`(f^{}(A_s))^{1/2}=j`$ we show that $$s^1GjBjGs^1jGBGj=s^1\left(jGB[j,G]+[G,j]BGj+[G,j]B[j,G]\right)$$ $`(2.9)`$ is a remainder for any bounded $`B=B^{}`$, using for $`[G,j]`$ the expansion $$[G,j]=\underset{k=1}{\overset{n1}{}}\frac{1}{k!}s^kj^{(k)}(A)ad_A^{(k)}(G)+s^nR$$ and its adjoint for $`[j,G]`$. Since $`ad_A^{(k)}(G)`$ is bounded for $`kn`$, the right hand side of (2.9) then becomes a sum of terms of the following types : $$\begin{array}{ccc}& s^{(k+l+1)}\left(j^{(k)}Cj^{(l)}+j^{(l)}C^{}j^{(k)}\right);(0k,ln1;k+l1);\hfill & (a)\hfill \\ & s^{(k+n+1)}\left(j^{(k)}C+C^{}j^{(k)}\right);(0kn1);\hfill & (b)\hfill \\ & s^{(2n+1)}C,\hfill & (c)\hfill \end{array}$$ where in each case $`C`$ stands for some operator which is bounded uniformly in $`a`$ and $`s`$. By the same arguments as in (2.6) and (2.7) these terms have corresponding upper and lower bounds $$\begin{array}{ccc}& \pm s^2\mathrm{\hspace{0.17em}2}C\left(j^{(k)2}+j^{(l)2}\right);\hfill & (a)\hfill \\ & \pm \left(s^2j^{(k)2}+s^{2n}C^2\right);\hfill & (b)\hfill \\ & \pm s^{(2n+1)}C.\hfill & (c)\hfill \end{array}$$ For any bounded $`B=B^{}`$ we therefore obtain $$s^1GjBjGs^1jGBGj,$$ $`(2.10)`$ meaning that the difference of the two expresssions is a remainder (2.4). For $`B=B_1=i[H_b,A]`$ we use Mourre’s inequality (2.8) to obtain from (2.3) : $$\begin{array}{cc}\hfill Gi[H,f]G& s^1GjB_1jGs^1jGB_1Gj\hfill \\ & s^1\theta jG^2js^1\theta Gj^2G\hfill \\ & =s^1\theta Gf^{}G,\hfill \end{array}$$ with remainders arising from (2.3) and twice from (2.10). Multiplying from both sides with $`g(H)`$ removes $`G(H)`$ and leads directly to (2.1). $`\text{ }\text{ }\text{ }\text{ }\text{ }`$ Proof of Theorem 1.1. We prove a slightly stronger version of (1.11), which will serve later in the proof of Theorem 1.2. Let $$A_{ts}=s^1(Aa\theta t),$$ and suppose that $`F`$ is a positive $`C^{\mathrm{}}`$-function of order $`1/2`$ on $`R`$ with $`F^{}0`$ and $`F(x)=0`$ for $`x0`$. Instead of (1.11) we show under the same hypothesis that $$F(A_{ts})e^{iHt}g(H)\chi ^+(Aa)\mathrm{const}.s^m$$ $`(2.11)`$ for $`m<n1`$, uniformly in $`0ts`$ and in $`aR`$. (1.11) then follows by setting $`t=s`$ and by observing that, since $`\vartheta <\theta `$, $`\chi ^{}(s^1(Aa)\vartheta )F(s^1(Aa)\theta )`$ for some $`F`$ of the required type. To prove (2.11) we consider the operator $$\varphi _s(t)=g(H)f(A_{ts})g(H);f=F^2,$$ and the evolution $$\psi _t=e^{iHt}\chi ^+(Aa)\phi ;\phi .$$ $`(2.12)`$ Then the estimate (2.11) to be proved reads $$\varphi _s(t)_t=(\psi _t,\varphi _s(t)\psi _t)\mathrm{const}.\phi ^2s^{2m},$$ $`(2.13)`$ uniformly in $`1<s`$ , $`0ts`$ and in $`aR`$. We compute $$\begin{array}{ccc}\hfill _t\varphi _s(t)_t& =(\psi _t,D_t\varphi _s(t)\psi _t);\hfill & (2.14)\hfill \\ \hfill D_t\varphi _s(t)& =i[H,\varphi _s(t)]+_t\varphi _s(t)\hfill & \\ & =g(H)i[H,f(A_{ts})]g(H)s^1\theta g(H)f^{}(A_{ts})g(H).\hfill & (2.15)\hfill \end{array}$$ First we conclude that $$D_t\varphi _s(t)\mathrm{const}.$$ $`(2.16)`$ uniformly in $`s`$, $`t`$, $`a`$, since $`f`$ is of order $`1`$. (cf. the remark after (3.4)). Secondly, by (3.8), $$\begin{array}{cc}\hfill \varphi _s(0)_0& \phi ^2F(s^1(Aa))g(H)\chi ^+(Aa)^2\hfill \\ & \mathrm{const}.s^{2n}\phi ^2.\hfill \end{array}$$ $`(2.17)`$ Integrating (2.14) over $`t`$ and using (2.16) and (2.17) we find the crude estimate $$\varphi _s(t)_t\mathrm{const}.\phi ^2(s^{2n}+s)$$ for $`0ts`$, which proves (2.13) for $`m=1/2`$. Now we bootstrap this estimate. First we note that by (2.15) and Lemma 2.1 $$D_t\varphi _s(t)s^{1\epsilon }g(H)f_1(A_{st})g(H)+\mathrm{const}.s^{(2n1\epsilon )}g^2(H).$$ $`(2.18)`$ As an induction assumption, suppose that (2.13) holds for some $`m<n1`$. Since $`f_1`$ also satisfies the hypothesis for $`f`$ it then follows from (2.18) that $$\|D_t\varphi _s(t)_t\|\mathrm{const}.\phi ^2s^1\left(s^{(2m+\epsilon )}+s^{(2(n1)\epsilon }\right),$$ and again by integrating over $`t`$ : $$\varphi _s(t)_t\mathrm{const}.\phi ^2\left(s^{2n}+s^{(2m+\epsilon )}+s^{(2(n1)\epsilon )}\right)$$ uniformly in $`0ts`$ and in $`aR`$. Recalling that $`\epsilon 1`$, the best decay for $`s\mathrm{}`$ is obtained by setting $$\epsilon =\mathrm{min}(1,(n1)m),$$ which boosts the exponent $`m`$ in (2.13) to $`m^{}=m+\epsilon /2`$. Therefore (2.13) holds for any $`m<n1`$. $`\text{ }\text{ }\text{ }\text{ }\text{ }`$ Proof of Theorem 1.2. It suffices to prove that $$\chi (t^2x^22t^1a\vartheta )e^{iHt}g(H)\chi ^+(Aa)\mathrm{const}.t^m$$ $`(2.19)`$ if $`\chi `$ is a smoothed characteristic function of $`(\mathrm{},\epsilon )`$ for some $`\epsilon >0`$ : $`\chi (x)=1`$ for large negative $`x`$, $`\chi ^{}0`$, and $`\chi (x)=0`$ for $`x\epsilon `$. Following the line of the previous proof we consider the operators $$\varphi _s(t)=f(x_{ts}^2);f=\chi ^2;x_{ts}^2=\frac{1}{s^2}(x^22at\vartheta t^2)$$ $`(2.20)`$ and the evolution $$\psi _t=e^{iHt}g(H)\chi ^+(Aa)\phi ,\phi ,$$ $`(2.21)`$ for $`0ts`$. The desired inequality (2.19) then reads $$\varphi _s(s)_s\mathrm{const}.s^{2m}.$$ $`(2.22)`$ Writing $`x^22at\vartheta ^2=(xa)^2(a^2+\vartheta )t^2`$ we see that $$\varphi _s(t)=0\text{ for }ts\sqrt{\frac{\epsilon }{a^2+\vartheta }}\alpha s,$$ and therefore $$\varphi _s(s)_s=\underset{\alpha s}{\overset{s}{}}𝑑tD_t\varphi _s(t)_s.$$ $`(2.23)`$ To find $`D_t\varphi _s(t)`$ we first compute $$\begin{array}{cc}\hfill i[H,\varphi _s(t)]& =\frac{i}{2}[p^2,\varphi _s(t)]=s^2\left(Af^{}(x_{ts}^2)+f^{}(x_{ts}^2)A\right)\hfill \\ & =2s^2u(x_{ts}^2)Au(x_{ts}^2),\hfill \end{array}$$ where we have factorized $`f^{}=u^2`$ and used that $`\left[[A,u]u\right]=0`$. Adding the term $$_t\varphi _s(t)=2s^2(a+\vartheta t)u^2(x_{ts}^2)$$ we arrive at $$D_t\varphi _s(t)=2ts^2u(x_{ts}^2)(t^1(Aa)\vartheta )u(x_{ts}^2).$$ $`(2.24)`$ Now we use $`\vartheta <\theta `$ to estimate $$\begin{array}{cc}\hfill (t^1(Aa)\vartheta )& (t^1(Aa)\vartheta )\chi ^{}(t^1(Aa)\vartheta )\hfill \\ & \left(F(t^1(Aa)\theta )\right)^2\hfill \end{array}$$ by some smooth function $`F`$ of order $`1/2`$ supported in $`R^{}`$ and with $`F^{}0`$. Setting $`A_t=t^1(Aa)\theta `$ we find $$D_t\varphi _s(t)_t2ts^2F(A_t)u(x_{ts}^2)e^{iHt}g(H)\chi ^+(Aa)\phi ^2.$$ $`(2.25)`$ On the other hand, it follows from (2.11) by setting $`t=s`$ that $$F(A_t)e^{iHt}g(X)\chi ^+(Aa)\mathrm{const}.t^m.$$ $`(2.26)`$ Before we can use this estimate in (2.25) we must commute the factor $`u(x_{ts}^2)`$ to the left. The required commutator can be expanded to any order $`n`$ : $$[u(x_{ts}^2),F(A_t)]=\underset{k=1}{\overset{n1}{}}\frac{1}{k!}t^kad_A^{(k)}(u)F^{(k)}(A_t)+t^nR,$$ where $`R\mathrm{const}.ad_A^{(n)}(u)`$. Since $`uC_0^{\mathrm{}}(R)`$, the commutators $`ad_A^{(k)}(u)`$ are easily bounded : $$\begin{array}{cc}\hfill iad_A^{(1)}(u)& =xu(x_{ts}^2)=2s^2x^2u^{}(x_{ts}^2)\hfill \\ & =2x_{ts}^2u^{}(x_{ts}^2)+2s^2(2at+\vartheta t^2)u^{}(x_{ts}^2)\hfill \end{array}$$ and so forth, with the result that $`ad_A^{(k)}(u)\mathrm{const}.`$ uniformly in $`1s<\mathrm{}`$ and $`0ts`$. Since (2.26) also holds if $`F`$ is replaced by a derivative $`F^{(k)}`$, we find the estimate $$\begin{array}{cc}\hfill D_t\varphi _s(t)_s& \mathrm{const}.ts^2(t^m+\underset{k=1}{\overset{n1}{}}t^{km}+t^n)^2\hfill \\ & \mathrm{const}.s^2t^{2m+1}\hfill \end{array}$$ for $`1s<\mathrm{}`$, $`0ts`$, uniformly in $`a`$. Therefore, by (2.23), $$\varphi _s(s)_s\mathrm{const}.s^2\underset{\alpha s}{\overset{s}{}}𝑑tt^{2m+1}\mathrm{const}.s^{2m}.$$ $`\text{ }\text{ }\text{ }\text{ }\text{ }`$ 3. COMMUTATOR EXPANSIONS Let $`H`$ and $`A`$ be self-adjoint operators on a Hilbert space $``$ and suppose that $`H`$ is bounded. To say that the commutor $`i[H,A]`$ is bounded means that the quadratic form $$i[(H\psi ,A\psi )(A\psi ,H\psi )]$$ on $`D(A)`$ is bounded and thus defines a bounded, symmetric operator called $`i[H,A]`$. In the same sense we assume that the higher commutators $$ad_A^{(k)}(H)=[ad_A^{(k1)}(H),A]$$ are bounded for $`k=2\mathrm{}n`$. Let $`f`$ be a real $`C^{\mathrm{}}`$–function on $`R`$. Then, under a further condition given below, the commutator $`[H,f(A)]`$ has an expansion $$[H,f(A)]=\underset{k=1}{\overset{n1}{}}\frac{1}{k!}f^{(k)}(A)ad_A^{(k)}(H)+R_n$$ $`(3.1)`$ with a remainder estimate $$R_nc_nad_A^{(n)}(H)\underset{k=0}{\overset{n+2}{}}𝑑x(1+\|x\|)^{kn1}\|f^{(k)}(x)\|.$$ $`(3.2)`$ The further condition on $`f`$ is that the integrals (3.2) exist. The number $`c_n`$ is a numerical constant depending on $`n`$ but not on $`f`$, $`A`$ or $`H`$. In particular, the expansion (3.1) holds if $$f^{(k)}(x)=O(\|x\|^{n\epsilon k})(x\pm \mathrm{})$$ $`(3.3)`$ for $`k=1\mathrm{}n+2`$, i.e. if the function $`f(x)`$ grows not faster than $`\|x\|^{n\epsilon }`$, with corresponding slower growth of the successive derivatives. We will refer to (3.3) by saying that $`f`$ is of order $`n\epsilon `$. In that case (3.1) is defined in form sense on the domain of $`f^{(1)}(A)`$. Taking the adjoint of (3.1) and noting that $$ad_A^{(k)}(H)^{}=(1)^kad_A^{(k)}(H),$$ we also obtain $$[H,f(A)]=\underset{k=1}{\overset{n1}{}}\frac{1}{k!}.(1)^{k1}ad_A^{(k)}(H)f^{(k)}(A)R_n^{}.$$ $`(3.4)`$ This defines $`[H,f(A)]`$ as an operator on the domain of $`f^{(1)}(A)=f^{}(A)`$. In particular, if $`f`$ is of order $`1`$ and $`n2`$, then $`[H,f(A)]`$ is bounded. In the general case where $`H`$ is not bounded, we will work with operators $`g(H)`$, $`gC_0^{\mathrm{}}(R)`$, assuming that $$ad_A^{(k)}(g(H))\text{ is bounded for }k=1\mathrm{}n.$$ $`(3.5)`$ Then, if $`f`$ is of order $`<n`$, $$\begin{array}{cc}\hfill [g(H),f(A)]& =\underset{k=1}{\overset{n1}{}}\frac{1}{k!}f^{(k)}(A)ad_A^{(k)}((g(H))+R_n;\hfill \\ \hfill R_n& \mathrm{const}.ad_A^{(n)}((g(H)),\hfill \end{array}$$ $`(3.6)`$ with a constant depending on $`f`$ and $`n`$, and similarly for the adjoint expansion. All these formulas are particularly useful if the role of $`A`$ is played by a scaled operator, say $`s^1A`$, $`0<s<\mathrm{}`$. Then the commutor expansions are expansions in powers of $`s^1`$, e.g. $$[g(H),f(s^1A)]=\underset{k=1}{\overset{n1}{}}\frac{1}{k!}s^kf^{(k)}(A)ad_A^{(k)}(g(H))+s^nR_n.$$ $`(3.7)`$ A simple but useful observation is the following. Suppose that $`f(x)=0`$ for $`xR^+`$, and let $`\chi ^+`$ be the characteristic function of $`R^+`$. Then $$\chi ^+(A)g(H)f(s^1A)\mathrm{const}.s^n.$$ $`(3.8)`$ Proof. Since $`\chi ^+(A)f(s^1A)=0`$ we have $$\chi ^+(A)g(H)f(s^1A)=\chi ^+(A)[g(H),f(s^1A)].$$ Inserting the expansion (3.7) we notice that only the remainder $`s^nR_n`$ contributes, since $$\chi ^+(A)f^{(k)}(A)=0.$$ $`\text{ }\text{ }\text{ }\text{ }\text{ }`$ Commutator expansions of this type were introduced in \[SiSo1\] and have since become an important tool of operator analysis. 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warning/0002/astro-ph0002382.html	ar5iv	text	# Eclipse studies of the dwarf-nova Ex Draconis ## 1 Introduction Dwarf novae are mass-exchanging binaries in which a late type star (the secondary) overfills its Roche lobe and transfers matter to a companion white dwarf (the primary) via an accretion disc. These systems show recurrent outbursts of 2–5 magnitudes on timescales of a few weeks to months caused either by an instability in the mass transfer from the secondary star or by a thermal instability in the accretion disc which switches the disc from a low to a high-viscosity regime (Warner 1995 and references therein). During outburst most of the light arises from the bright, optically thick accretion disc, while in quiescence the dominant sources of light are the white dwarf and the bright spot formed by the impact of the infalling gas stream with the edge of the disc. Eclipsing dwarf novae are probably the best sites for the study of accretion physics as the occultation of the accretion disc and white dwarf by the secondary can be used to constrain the geometry and parameters of the binary, and tomographic techniques such as eclipse mapping (Horne 1985) and Doppler tomography (Marsh & Horne 1988) can be applied to probe the structure and dynamics of the accretion flow. EX Draconis (= HS1804+67) was detected in the Hamburger Quasar Survey (Bade et al. 1989) and shown to be an eclipsing dwarf nova with an orbital period of 5.04 hr by Barwig et al. (1993). From spectroscopic observations made in quiescence, Billington, Marsh & Dhillon (1996) found that the secondary star is of spectral type M1 to M2 and that it contributes almost all of the light at mid-eclipse. Their analysis showed that the inner face of the secondary is significantly irradiated by the white dwarf. They found a rotational broadening of $`v\mathrm{sin}i=140kms^1`$ and a radial velocity semi-amplitude of $`K_2=210kms^1`$ for the secondary star which leads to a spectroscopic mass ratio of $`q=0.8`$ when combined with the $`K_1=167kms^1`$ of Barwig et al. (1993). A relatively small value for the radius of the accretion disc ($`0.4R_{L1}`$) is derived but no explanation is given of how this estimate was made. In a follow up study using spectroscopy and photometry of EX Dra in quiescence and in outburst, Fiedler, Barwig & Mantel (1997) measured radial velocity semi-amplitudes of $`K_1=167kms^1`$ and $`K_2=223kms^1`$ and derived a spectroscopic model for the binary with $`q=0.75`$, $`i=84.2\text{}`$, $`M_1=0.75M_{}`$ and $`M_2=0.56M_{}`$. However, the radial velocity curve of the H$`\alpha `$ line shows a large phase shift ($`0.2`$ cycle) with respect to photometric conjunction which casts doubt on the derived value of $`K_1`$. They use the eclipse phases of the bright spot and white dwarf to derive a photometric mass ratio between 0.7 and 0.8, supporting the spectroscopic model. From the ratios of Ca I and TiO absorption features they infer a spectral type of M0 for the secondary star. Smith & Dhillon (1998) use the values of $`v\mathrm{sin}i`$ and $`K_2`$ of Billington et al. (1996) and the eclipse phase width $`\mathrm{\Delta }\varphi `$ of Fiedler et al. (1997) to infer a $`K_1=176kms^1`$. In this paper we present and discuss high-speed photometry of EX Dra in quiescence and in outburst. Section 2 describes the observations. In section 3 we present and discuss the eclipse lightcurves, provide an updated ephemeris, derive the binary parameters from the eclipse phases of the white dwarf and bright spot, and obtain estimates of the distance to the binary. The results are summarized in section 4. ## 2 Observations and data reduction Time-series of differential photometry of EX Dra in the $`V`$ and $`R`$ bands was obtained with a Wright Instruments CCD camera (1.76 arcsec/pixel, $`385\times 578`$ pixels) attached to the 0.9-m James Gregory Telescope of the University Observatory, St. Andrews, in 1995 and 1996. This pixel size is matched to the seeing at this sea-level site, where typical stellar images have FWHM values of 3.0 pixels, and are therefore well sampled. Exposure times ranged from 15 to 40 s with a dead-time between exposures of about 5 s to read the CCD chip. Details of the observations are listed in Table 1. The observations include five outbursts of EX Dra and sample various phases along the outburst cycle. The data was reduced using standard IRAF<sup>2</sup><sup>2</sup>2 IRAF is distributed by National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc., under contract with the National Science Foundation. procedures and included bias and flat-field corrections and cosmic rays removal. Photometry was obtained with the automated aperture photometry routine JGTPHOT (Bell, Hilditch & Edwin 1993). Fluxes were extracted for the variable and for five selected comparison stars in the field. The relative brightness of the comparison stars in all data sets is constant to better than 0.01 mag. We adopted a mean comparison star magnitude for each frame from the intensity-added values of these five stars. Time-series were constructed by computing the magnitude difference between the variable and the mean comparison star. From the dispersion in the magnitude difference of the comparison stars with similar brightness we estimate an uncertainty in the photometry of EX Dra of 0.025 mag in quiescence and better than 0.01 mag in outburst. Observations of spectrophotometric standard stars of Massey et al. (1998) and Massey & Gronwall (1990) were used to calibrate the photometry in the EX Dra field. These observations demonstrate that the transformation coefficients from the natural to the standard $`VR`$ system (Bessell 1983) are unity to a precision of one per cent. Hence differential instrumental magnitudes obtained from individual frames are differential $`V`$ and $`R`$ magnitudes. We found the mean of the combined $`V`$ and $`R`$ magnitudes of the five comparison stars to be $`V=13.21\pm 0.09`$ mag and $`R=12.84\pm 0.05`$ mag. We used the relations $`V=16.402.5\mathrm{log}_{10}f_\nu [\mathrm{mJy}]`$ and $`R=16.222.5\mathrm{log}_{10}f_\nu [\mathrm{mJy}]`$ (Lamla 1982) to transform the calibrated magnitudes to absolute flux units. ## 3 Results ### 3.1 Eclipse lightcurves We adopted the following convention regarding the phases: conjunction occurs at phase zero, the phases are negative before conjunction and positive afterwards. The lightcurves were phased according to the ephemeris of eq. 1 (section 3.2). Figure 1 shows the visual lightcurve of EX Dra for the period September 1995 to January 1996 from AAVSO and VSNET observations. Vertical dotted lines mark the epochs of our observations while filled circles show the corresponding $`R`$-band out-of-eclipse magnitudes. There were six recorded outbursts during this period (labeled from A to F in Fig. 1), with typical amplitudes of $`2.0`$ mag, duration of $`10`$ days, and average time between outbursts of $`20\pm 3`$ days. Outburst C was shorter ($`5`$ days) and weaker ($`\mathrm{\Delta }m1.5`$ mag) than the others and, unfortunately, was not covered by our observations. The visual magnitude is typically $`m_v12.7`$ at maximum and $`m_v15`$ in quiescence. At the end of outbursts A and E the star went through a faint state (hereafter named low state) before recovering its usual quiescent level. Our dataset frames eclipse lightcurves during most relevant phases through the outburst cycle: early rise to maximum (cycles 7812, 7813, 7817), late rise to maximum (7927), maximum light (7540, 7659, 7660, 7826 and 7827), end of maximum (7669, 7670, 7673, 7674 and 9583), early decline (7559, 7850), late decline (7564, 7970), low state (8002 to 8004), and quiescence (7569, 7798, 8007 to 8017, and 8074). Individual lightcurves of EX Dra in quiescence are shown in Fig. 2. Sharp changes in the slope reflect the occultation of a compact source at disc centre and of the bright spot at disc rim. The ingress of the central source and of the bright spot overlap in phase to form a unique sharp break in the slope. The egress of the central source is variable both in duration and in flux and sometimes is hardly visible (e.g., cycle 8074). The eclipses have a flat-bottomed, ‘U’ shape indicating that the eclipse is close from being total. There is no pronounced flickering (amplitude $`\stackrel{<}{}2.5`$ per cent). In some of the lightcurves the flux level is the same before and after eclipse while others show a perceptible orbital hump prior to eclipse (usually interpreted as being the result of anisotropic emission from the bright spot) and a slow recovery from eclipse where the bright spot egress is hardly discernible. Figure 3 shows lightcurves of EX Dra along the outburst cycle. The left panel shows the behaviour during rise and maximum light while the right panel shows the behaviour during decline and in the low state. The lightcurves were grouped by outburst phase. Only a subset of the outburst lightcurves are shown for clarity. The fact that, for a given outburst stage, the lightcurves of different outbursts have similar eclipse shapes and out-of-eclipse levels gives confidence that the observed sequence is representative of the general behaviour of EX Dra during outburst (at least for the epoch of our observations). Lightcurve 7812+7813 frames the early rise to maximum. The eclipse shape is asymmetric and mid-eclipse occurs after phase zero, indicating that the receding side of the disc is brighter. The mid-eclipse level and the total eclipse width are the same as in quiescence, showing that the brightening does not start in the outer disc regions. This is in agreement with the symmetric shape of the outbursts, with comparable rise and decline timescales (cf. Fig. 1), which is typical of inside-out outbursts (e.g., Cannizzo, Wheeler & Polidan 1986). The eclipse profile changes during the rise to maximum, from the asymmetric ‘U’ shape eclipse of a compact central source plus the bright spot at disc rim to a more symmetric ‘V’ shape indicating the partial occultation of a bright extended disc. The total width of the eclipse increases during the rise (from 0.196 cycle in quiescence to about 0.22 cycle at lightcurve 7927 and larger at maximum) indicating that the disc radius also increases as the system approaches maximum light. A precise measurement of the total width of the eclipse at maximum light is precluded due to the limited phase coverage of the corresponding lightcurve. The disc shrinks during decline (as indicated by the eclipse egress phase) until it reaches the quiescent radius close to the end of the outburst (lightcurve 7564+7970). The low state is characterized by an out-of-eclipse level $`15`$ per cent lower than the typical quiescent level and by a fairly symmetric eclipse shape, corresponding to the eclipse of a compact source at the disc centre with little evidence of a bright spot at the disc rim. The mid-eclipse level of lightcurves 7564+7970 and 8002+8003 is the same, showing that the decrease in brightness at this stage is due to the fading of the light from the inner disc regions. Flickering is much more pronounced in outburst than in quiescence. Flickers of an amplitude of $`1015`$ per cent can be seen in many lightcurves in outburst. None of our outburst lightcurves resembles the flat-bottomed outburst lightcurve of Fiedler et al. (1997, see their fig. 6). Their outburst lightcurve is a factor of only $`2`$ brighter than their typical quiescent lightcurve indicating that it corresponds to outbursts of lower amplitude than the ones sampled by our observations. ### 3.2 Revised ephemeris The ingress feature of the white dwarf and of the bright spot overlap in quiescent lightcurves of EX Dra (Fig. 2). Since the bright spot ingress depends on the variable disc radius, the mid-ingress time is not a stable feature of the eclipse. We therefore adopted the same procedure of Fiedler et al. (1997) and used the mid-egress times of the white dwarf plus the inferred duration of its eclipse (see section 3.3) to obtain a revised ephemeris for the mid-eclipse times. Mid-egress times were measured by computing the time of maximum derivative in a median-filtered version of the lightcurve (section 3.3). The uncertainty in determining mid-egress times depends on the time resolution and signal-to-noise of the lightcurve and is in the range $`(12)\times 10^4`$ d. The barycentric correction and the difference between universal times (UT) and dynamical ephemeris time scales are smaller than the uncertainties in the measured timings and were neglected. The new heliocentric (HJD) times of the egress of the white dwarf are listed in Table 2 with corresponding cycle number and uncertainties (quoted in parenthesis). We assumed equal errors of $`10^4`$ d for the timings of Fiedler et al. (1997) and combined them with the timings of Table 2 to obtain a least-squares linear fit with a reduced $`\chi ^2=13.8`$ for 54 d.o.f. and a standard deviation of $`\sigma =0.0019`$ cycle. The residuals with respect to the linear ephemeris show a clear cyclical behaviour and can be well described by a sinusoidal function. Assuming a duration of the eclipse of the white dwarf of $`\mathrm{\Delta }t=0.0228`$ d (section 3.3), the best-fit linear plus sinusoidal ephemeris for the mid-eclipse times is, $$T_{mid}=\mathrm{HJD}\mathrm{\hspace{0.33em}2\hspace{0.17em}448\hspace{0.17em}398.4530}(\pm 1)+\mathrm{0.209\hspace{0.17em}936\hspace{0.17em}98}(\pm 4)E+$$ $$+(8.2\pm 1.5)\times 10^4\mathrm{sin}\left[2\pi \frac{(E968)}{7045}\right]d,$$ (1) with $`\chi ^2=2.7`$ for 51 d.o.f. and $`\sigma =0.0010`$ cycle. Residuals with respect to the linear part of eq.(1) are listed in Table 2 and shown in Fig. 4. The sinusoidal term of eq.(1) is indicated in Fig. 4 as a dotted line. The amplitude (1.18 minutes) and timescale ($`4`$ years) of the period variation are similar to the quasi-periodic orbital period changes found in many other eclipsing CVs (Warner 1995 and references therein) and are possibly related to a solar-like magnetic activity cycle in the secondary star (Applegate 1992; Richman, Applegate & Patterson 1994). It may also be possible that the period changes are due to a third body in the system, as suggested by Fiedler et al. (1997), if the variation proves to be strictly periodic. Regular observations of eclipse timings during the next decade are required in order to check the stability of the period of the variation and test the above hypotheses. ### 3.3 Binary parameters #### 3.3.1 Measuring eclipse phases The ingress/egress phases of the occultation of the compact central source (hereafter CS) and of the bright spot (BS) by the secondary star provide information about the geometry of the binary system and the relative sizes of these components (e.g., Wood et al. 1986). We used the lightcurves of the low state – where the effect of the BS is minimal on the eclipse shape – to measure the ingress and egress phases of the CS and to derive the width of its eclipse as well as the duration of its ingress/egress feature. The contact phases can be identified as rapid changes in slope visible in the lightcurves of the low state (Fig. 3) and were measured with the aid of a cursor on a graphic display of a median filtered version of the lightcurve. We defined $`\varphi _{c1},\varphi _{c2}`$ as those phases during which the CS disappears behind the secondary star and $`\varphi _{c3},\varphi _{c4}`$ as the phases corresponding to its reappearance from eclipse. The mid-ingress (egress) phases ($`\varphi _{ci},\varphi _{ce}`$) were computed as the phases at which half of the central source light is eclipsed and also as the phases of minimum (maximum) derivative in the lightcurve (e.g., Wood, Irwin & Pringle 1985). Figure 5 illustrates the measurement procedure with the derivative for the combined lightcurve 8002+8003. The ingress/egress of the CS can be seen as those intervals for which the derivative curve is significantly different from zero. The width at half-peak intensity of these features yields a preliminary estimate of their duration. A spline function is fitted to the remaining regions in the derivative curve to remove the contribution from the extended and slowly varying eclipse of the disc. Estimates of the CS flux at ingress (egress) are obtained by integrating the flux in the spline-subtracted derivative curve between the first and second (third and fourth) contact phases. The lightcurve of CS is then reconstructed by assuming that the flux is zero between ingress and egress and that it is constant outside eclipse. The reconstructed CS lightcurve can be seen in Fig. 5(c) and the lightcurve after removal of the CS component is shown in Fig. 5(d). The measured contact phases, mid-ingress and mid-egress phases of the CS from the lightcurves of the low state are collected in Table 3. The quoted mid-ingress/egress values are the average of both procedures described above and have an estimated error of 0.0005 cycle. The duration of the CS eclipse (the eclipse of the disc centre) is defined as $$\mathrm{\Delta }\varphi =\varphi _{ce}\varphi _{ci},$$ (2) and the mid-eclipse phase (the inferior conjunction of the binary) is written as $$\varphi _0=1/2(\varphi _{ce}+\varphi _{ci}).$$ (3) These quantities are collected in Table 3. The mean of the measurements from all lightcurves yields $`\mathrm{\Delta }\varphi =0.1085\pm 0.0006`$ cycle (=0.0228 d), where the quoted error is the standard deviation of the mean. Similarly, we have $`\varphi _0=+0.0001\pm 0.0003`$ cycle, which indicates that the centre of the CS eclipse corresponds to phase zero. The difference between the first and second (third and fourth) CS contact phases yield the phase width of the CS ingress (egress), $`\mathrm{\Delta }_{\mathrm{ci}}`$ ($`\mathrm{\Delta }_{\mathrm{ce}}`$). These quantities are also listed in Table 3. A mean from all values of $`\mathrm{\Delta }_{\mathrm{ci}}`$ and $`\mathrm{\Delta }_{\mathrm{ce}}`$ yields $`\mathrm{\Delta }_{\mathrm{cs}}=0.0082\pm 0.0003`$ cycle. BS ingress/egress phases ($`\varphi _{bi},\varphi _{be}`$) were measured from the lightcurves in quiescence in which it was possible to simultaneously identify the eclipse of BS and the egress of CS. We measured the CS contact and mid-egress phases and used the derived value of $`\mathrm{\Delta }\varphi `$ to reconstruct the lightcurve of CS assuming that the flux and duration of its ingress feature are the same as in egress. Mid-ingress/egress phases of BS were measured in the lightcurves after removal of the CS component, which provide and unblended, clean view of the BS ingress feature (Fig. 6). The eclipse parameters measured from the lightcurves in quiescence are listed in Table 4. The BS eclipse in lightcurve 7569 starts earlier and ends later than in the other lightcurves, indicating a relatively larger disc radius at this epoch (see section 3.3.2). #### 3.3.2 Mass ratio, inclination and disc radius Making the usual assumption that the secondary star fills its Roche lobe and given the duration of the eclipse of the central parts of the disc, $`\mathrm{\Delta }\varphi `$, there is a unique relation between the mass ratio $`q=M_2/M_1`$ and the binary inclination $`i`$ (Bailey 1979; Horne 1985). From Table 3, the width of the eclipse in EX Dra is $`\mathrm{\Delta }\varphi =0.1085`$. This gives the constraint $`q>0.64`$, with $`q=0.64`$ if $`i=90\text{}`$. When combined with the measured eclipse phases of the CS and BS, this relation gives a unique solution for $`q`$, $`i`$, and $`R_{bs}/R_{L1}`$, where $`R_{bs}`$ is the distance from disc centre to the BS (usually taken to be the disc radius) and $`R_{L1}`$ is the distance from disc centre to the inner lagrangian point L1 (e.g., Smak 1971; Cook & Warner 1984). Fig. 7(a) shows a diagram of ingress versus egress phases for the measurements of the CS and BS in Tables 3 and 4. Measurements of the CS ingress and egress are shown as the cluster of small diamonds around phases ($`0.054,+0.054`$) in the lower portion of the diagram. Eclipse phases of BS are indicated by crosses. Theoretical gas stream trajectories corresponding to a set of pairs $`(i,q)`$ are also shown. The trajectories were computed by solving the equations of motion in a coordinate system synchronously rotating with the binary, using a 4th order Runge-Kutta algorithm (Press et al. 1986) and conserving the Jacobi integral constant to one part in $`10^6`$. The correct mass ratio, and hence inclination, are those for which the calculated stream trajectory passes through the observed position of the bright spot. This yields $`q=0.72\pm 0.06`$ and $`i=85_2^{+3}`$ degrees, where the uncertainties are taken from the standard deviation of the points about the trajectory of best fit. Fig. 7(b) shows the geometry of the binary system for $`q=0.72`$. For this mass ratio, the relative size of the primary Roche lobe is $`R_{L1}/a=0.534\pm 0.009`$, where $`a`$ is the orbital separation. BS eclipse phases are clustered at two distinct positions along the best-fit stream trajectory. The squashed circles in Fig. 7(a) represent the accretion discs whose edges pass through these positions for the adopted mass ratio. This corresponds to disc radii of $`R_{bs}/R_{L1}=(0.50\pm 0.01)`$ and $`(0.56\pm 0.01)`$ with the bright spot making angles of, respectively, $`\alpha _{bs}=20\text{}\pm 1\text{}`$ and $`16\text{}\pm 1\text{}`$ with respect to the line joining both stars. Circles with these radii are depicted in Fig. 7(b). The larger disc radius comes from measurements of BS phases just after the end of an outburst (lightcurve 7569, end of outburst A) while the remaining points correspond to lightcurves well into quiescence (lightcurves 8007, 8008, 8012, 8016+8017). This result suggests that the accretion disc of EX Dra shrinks (by at least $`12`$ per cent) during quiescence – a similar behaviour to that found in other dwarf novae (e.g., Smak 1984, 1991; Wood et al. 1989). The calculated radii are a factor of 2–3 larger than the radius expected for zero-viscosity discs, $`R_d/R_{L1}=0.19`$ (Flannery 1975), but are smaller than the radius expected for pressureless discs, $`R_d/R_{L1}=0.66`$ (Paczyński 1977). Lightcurve 8074 gives the best phase coverage of the orbital hump in our dataset. The hump can be well described by a sinusoid of amplitude 0.6 mJy and maximum at orbital phase $`0.17\pm 0.01`$ cycle. The direction of hump maximum (i.e., maximum visibility of the bright spot) is indicated in Fig. 7(b) by an arrow; it is clearly different from the radial direction of the bright spot. If the hump maximum is normal to the plane of the shock at the bright spot site then the shock lies in a direction between the stream trajectory and the edge of the disc, making an angle of $`41\text{}\pm 4\text{}`$ with the latter. #### 3.3.3 Masses and radii of the component stars An estimate of the binary parameters of EX Dra may be obtained by combining the inferred mass ratio $`q`$ with the empirical main sequence mass-radius relation of Smith & Dhillon (1998), $$R_2/R_{}=\alpha \left(M_2/M_{}\right)^\beta ,$$ (4) where $`\alpha =0.91\pm 0.09`$ and $`\beta =0.75\pm 0.04`$. The latter assumption seems reasonably well justified by the good quality of the fit to the measured masses of secondary stars in CVs and of field main sequence stars. The primary-secondary mass diagram for EX Dra can be seen in Fig. 8. The constraints from the mass ratio and the empirical mass-radius relation are shown as thick solid lines. We also plotted lines corresponding to the mass functions for a radial velocity of the white dwarf of $`K_1=167kms^1`$ (Fiedler et al. 1997) and for a radial velocity of the secondary star of $`K_2=210kms^1`$ (Billington et al. 1996). The four relations are consistent at the 1-$`\sigma `$ level. A Monte Carlo propagation code was used to estimate the errors in the calculated parameters. The values of the input parameters $`q`$ and ($`\alpha ,\beta )`$ are independently varied according to Gaussian distributions with standard deviation equal to the corresponding uncertainties. The results, together with their 1-$`\sigma `$ errors, are listed in Table 5. The quoted $`K_1,K_2`$ and $`v\mathrm{sin}i`$ are the predicted values of, respectively, the radial velocity of the primary and secondary stars and the secondary star rotational velocity; $`a`$ is the binary separation, and the remaining parameters are self-explanatory. The cloud of points in Fig. 8 was obtained from a set of $`10^4`$ trials using this code. The highest concentration of points indicates the region of most probable solutions. Table 5 also lists the estimated parameters of Billington et al. (1996), Fiedler et al. (1997) and Smith & Dhillon (1998). Our model of the EX Dra binary is in reasonably good agreement with the models independently derived by those authors from spectroscopic measurements. We now turn our attention to the compact central source. An estimate of its diameter can be obtained from the duration of its ingress/egress feature in the lightcurve, $`\mathrm{\Delta }_{cs}`$, using the approximate relations (Ritter 1980), $$R_{cs}/a=\pi z(q)\mathrm{\Delta }_{cs}\mathrm{sin}\theta ,\mathrm{cos}\theta =\frac{a}{R_2}\mathrm{cos}i,$$ (5) where $`R_{cs}`$ and $`R_2`$ are the radii of the central source and of the secondary star, respectively, and $`z(q)`$ is the distance (in units of $`a`$) from disc centre to the point tangent to the surface of the secondary that marks the beginning/end of the eclipse of CS (Baptista et al. 1989). $`z(q)`$ is a slow varying function and is usually close to unity. In our case, for $`q=0.72`$ we have $`z=0.924`$. For the following exercise we adopted the value of $`\mathrm{\Delta }_{cs}=0.0082`$ cycle inferred from the lightcurves of the low state (section 3.3.1). We note that the width of the CS ingress/egress feature is usually larger than this in the lightcurves in quiescence (see Table 4); from the mean $`B`$ lightcurve in quiescence of Fiedler et al. (1997) (see their fig. 7) we estimate a width of the egress feature of 0.010 cycle. Assuming that the compact central source is the white dwarf, the substitution of Kepler’s third law into equation (5) yields a relation between the mass and the radius of the white dwarf than can be combined with the Hamada-Salpeter (1961) mass-radius relation (c.f. Nauenberg 1972) to eliminate $`R_{cs}`$ and solve for $`M_1(q,\mathrm{\Delta }_{cs})`$ (e.g., Baptista et al. 1998). This relation (plotted in Fig. 8 as a dashed line) predicts unreasonably low white dwarf masses of $`M_10.2M_{}`$, in clear disagreement with the results in Table 5. The discrepancy is not alleviated by the use of a mass-radius relation for hot white dwarfs (Koester & Shönberner 1986; Vennes, Fontaine & Brassard 1995), the inclusion of possible spherical distortion effects due to fast rotation of the white dwarf or the consideration of strong limb-darkening effects (Wood & Horne 1990), or by adopting the larger value of $`\mathrm{\Delta }_{cs}`$ from the quiescent lightcurves. The measured $`\mathrm{\Delta }_{cs}`$ is simply too large for a $`0.7M_{}`$ white dwarf in a binary with the orbital period of EX Dra. Together with the variability in flux and duration of the ingress/egress feature this leads to the conclusion that the observed compact central source is not a bare white dwarf. The substitution of the parameters of Table 5 into equation (5) yields $`R_{cs}(\mathrm{\Delta }_{cs}=0.0082)=0.23a=0.037R_{}=3.36R_1`$. Therefore, the white dwarf in EX Dra seems surrounded by an extended, variable atmosphere or boundary layer of at least 3 times its radius. In this regard EX Dra is similar to the long period, eclipsing dwarf nova IP Peg – where the white dwarf seems to be wrapped in a thick boundary layer more than twice its radius (Wood & Crawford 1986) – and is clearly different from the short-period dwarf novae OY Car, Z Cha and HT Cas, where the central source seems to be a bare white dwarf (Wood & Horne 1990). This result suggests that different physical conditions may exist in the inner disc regions of the short period and of the long period dwarf novae, possibly related to the distinct mass accretion rates of these groups, although the statistics of eclipsing dwarf novae is still very low on both sides of the CV period gap. From the derived parameters, we predict a duration of the ingress/egress of the white dwarf of $`\mathrm{\Delta }_{wd}=0.0024\pm 0.0006`$ cycle. ### 3.4 Distance estimates The flux densities at mid-eclipse of the flat-bottomed lightcurves of the low state (Fig. 3) yield an upper limit to the contribution of the secondary star that can be used to set a lower limit on the distance to EX Dra (e.g., Baptista, Steiner & Cieslinski 1994). We find $`F_{\mathrm{mid}}(V)=0.79\pm 0.05`$ mJy and $`F_{\mathrm{mid}}(R)=1.56\pm 0.03`$ mJy, where the quoted values are the median of the fluxes in the phase range ($`0.02,+0.02`$) cycle and the uncertainties were derived from the median of the absolute deviations with respect to the median. This corresponds to an apparent magnitude of $`V_{\mathrm{mid}}=16.66\pm 0.07`$ mag and a color index of $`(VR)=+0.92\pm 0.07`$ mag. We estimated a reddening of $`E(BV)=0.15\mathrm{mag}kpc^1(A_v=4.8\times 10^4\mathrm{mag}pc^1)`$ for EX Dra from the galactic interstellar extinction contour maps of Lucke (1978), which gives a color excess of $`E(VR)=0.02`$ mag for a distance of $`D=290pc`$ (see below). This leads to a corrected, intrinsic color index of $`(VR)_0=+0.90\pm 0.07`$ mag. A main sequence star with this color index has a spectral type M$`0\pm 2`$, an absolute magnitude of $`M_v=9.2\pm 0.6`$ mag and an effective surface temperature of $`T_{\mathrm{eff}}=3850\pm 200K`$ (Schmidt-Kaler 1982; Pickles 1985). This spectral type is in good agreement with that inferred from the spectroscopy by Billington et al. (1996) and Fiedler et al. (1997). With the assumption that the observed properties of the secondary star in EX Dra are similar to those of a normal main sequence star of same mass, we replace these values into the equation, $$5\mathrm{log}D(pc)=V_{\mathrm{mid}}M_v+5A_vD(pc),$$ (6) to find a photometric parallax distance of $`D_{MS}=(290\pm 80)pc`$, where the quoted uncertainty is obtained from the propagation of errors in the input parameters and do not account for possible systematic errors. If we neglect the interstellar extinction, the inferred distance is reduced to $`D=280\pm 80pc`$. An alternative blackbody fit to the mid-eclipse fluxes including the interstellar extinction yields a solid angle of $`\theta _{BB}^2=[(R_2/R_{})/(D/kpc)]^2=3.8\pm 0.8`$, an $`T_{\mathrm{eff}}=3550\pm 250K`$, and a distance of $`D_{BB}=(290\pm 40)pc`$. The measured mid-eclipse fluxes and the best-fit main sequence and blackbody spectra are shown in the upper panel of Fig. 9. Another distance estimate can be obtained from the $`V`$ and $`R`$ flux densities of the compact central source in the lightcurves of the low state, since we also have an estimate of the physical dimension of this source. The fluxes at egress, which are free from contamination of light from the bright spot, were obtained by integrating the derivative between the third and fourth contact phases (see section 3.3.1). We find $`F_{cs}(V)=0.55\pm 0.05`$ mJy and $`F_{cs}(R)=0.42\pm 0.04`$ mJy. We fitted the observed flux densities from synthetic photometry with white dwarf atmosphere models (Gänsicke, Beuermann & de Martino 1995) allowing for the effect of interstellar extinction as estimated above. The best-fit model has $`T_{cs}=(28\pm 3)\times 10^3K,\mathrm{log}g=8`$ and $`\theta _{cs}^2=(3.3\pm 0.5)\times 10^3`$. The measured central source fluxes and best-fit white dwarf atmosphere model are shown in the lower panel of Fig. 9. The corresponding distance depends on the geometry and effective emitting area of the central source as seen by an observer on Earth. A spherical central source has an effective area of $`A_{sp}=\pi R_{cs}^2`$ as projected onto the plane of the sky. In this case, if the inner disc is optically thin, both hemispheres of the central source are seen and a distance of $`D=640\pm 50pc`$ is obtained. If the inner disc is opaque the lower hemisphere of the central source is occulted and the distance is reduced to $`D=450\pm 40pc`$. Both values are in agreement with the lower limit derived from the contribution of the secondary star. Alternatively, if we assume that the distance is $`D=290pc`$, the inferred $`\theta _{cs}^2`$ allow us to constrain the geometry of the central source. At this distance the effective area of the central source is reduced to 21 per cent of $`A_{sp}`$. Since the equatorial diameter of the central source is set by the width of its ingress/egress feature to be $`2R_{cs}`$, this implies that the polar diameter is significantly smaller than $`R_{cs}`$. Hence, the distance estimate from the flux of the secondary star and that from the flux of the central source can be reconciled if the central source has a toroidal shape with an equatorial diameter of $`2\times 0.037R_{}`$ and a vertical thickness of $`0.012R_{}`$ (if the inner disc is optically thin) or $`0.024R_{}`$ (if the inner disc is opaque). ## 4 Summary The results of the analysis of $`V`$ and $`R`$ high speed photometry of EX Dra in quiescence and through outburst can be summarized as follows: 1. During the period of the observations EX Dra showed outbursts with typical amplitudes of $`2.0`$ mag, duration of $`10`$ days, and average time between outbursts of $`20\pm 3`$ days. The observed amplitudes are larger than those found by Billington et al. (1996). 2. The lightcurves during outburst were grouped by outburst phase. The analysis of these lightcurves indicates that the outbursts do not start in the outer disc regions and, therefore, favours the disc instability model. The disc expands during the rise to maximum (as indicated by the increasing width of the eclipse) and shrinks during decline. The decrease in brightness at the later stages of the outburst is due to the fading of the light from the inner disc regions. 3. At the end of two outbursts the system was seen to go through a phase of lower brightness (named the low state), characterized by the fairly symmetric eclipse of a compact source at disc centre with little evidence of a bright spot at disc rim, and by an out-of-eclipse level $`15`$ per cent lower than the typical quiescent level. 4. New eclipse timings were measured from the lightcurves in quiescence and a revised ephemeris was derived. The residuals with respect to the linear ephemeris show a clear cyclical behaviour and can be well described by a sinusoid of amplitude 1.2 minutes and period $`4`$ years. This period variation is possibly related to a solar-like magnetic activity cycle in the secondary star. 5. Eclipse phases of the compact central source and of the bright spot were used to derive the geometry of the binary. By constraining the gas stream trajectory to pass through the observed position of the bright spot we find $`q=0.72\pm 0.06`$ and $`i=85_2^{+3}`$ degrees. 6. The binary parameters were estimated by combining the measured mass ratio with the assumption that the secondary star in EX Dra obeys the empirical main sequence mass-radius relation of Smith & Dhillon (1998). The set of derived parameters in listed in Table 5. 7. The observed changes in the position of the bright spot with time suggest that the accretion disc shrinks during quiescence by at least $`12`$ per cent. 8. The phase of hump maximum is distinct from the radial direction of the bright spot. If the hump maximum is normal to the plane of the shock at the bright spot site then the shock lies in a direction between the stream trajectory and the edge of the disc, making an angle of $`41\text{}\pm 4\text{}`$ with the latter. 9. The white dwarf seems surrounded by an extended, variable atmosphere or boundary layer of at least 3 times its radius. From the derived parameters, a duration of the ingress/egress of the white dwarf of $`\mathrm{\Delta }_{wd}=0.0024\pm 0.0006`$ cycle is predicted. 10. The fluxes at mid-eclipse of the lightcurves of the low state yield an upper limit to the contribution of the secondary star and lead to a lower limit photometric parallax distance of $`D_{MS}=290\pm 80pc`$. 11. The fluxes of the central source are well fitted by a white dwarf atmosphere model with $`T_{cs}=(28\pm 3)\times 10^3K,\mathrm{log}g=8`$ and solid angle $`\theta _{cs}^2=[(R_{cs}/R_{})/(D/kpc)]^2=(3.3\pm 0.5)\times 10^3`$. For a spherical central source, this leads to a distance of $`D=640\pm 50pc`$ if the inner disc is optically thin. The distance estimates from the mid-eclipse fluxes and from the fluxes of the central source can be reconciled if the central source has a toroidal shape with an equatorial diameter of $`2\times 0.037R_{}`$ and a vertical thickness of $`0.012R_{}`$ (if the inner disc is optically thin) or $`0.024R_{}`$ (if the inner disc is opaque). The analysis of the set of lightcurves through outburst with eclipse mapping techniques yields an uneven opportunity to investigate the changes in the structure of an outbursting accretion disc and is the subject of another paper (Baptista & Catalán 1999, 2000). ## Acknowledgments We are grateful to Yvonne Unruh for valuable help in obtaining the data at JGT and to Boris Gänsicke for kindly providing the white dwarf atmosphere models. In this research we have used, and acknowledge with thanks, data from the AAVSO International Database and the VSNET that are based on observations collected by variable star observers worldwide. RB acknowledges financial support from CNPq/Brazil through grant no. 300 354/96-7. MSC acknowledges financial support from a PPARC post-doctoral grant during part of this work. This research was partially supported by PRONEX grant FAURGS/FINEP 7697.1003.00.
warning/0002/cond-mat0002324.html	ar5iv	text	# Transient Dynamics of Pinning ## Abstract We study the evolution of an elastic string into the pinned state at driving forces slightly below the depinning threshold force $`F_c`$. We quantify the temporal evolution of the string by an activity function $`A(t)`$ representing the fraction of active nodes at time $`t`$ and find three distinct dynamic regimes. There is an initial stage of fast decay of the activity; in the second, intermediate, regime, an exponential decay of activity is observed; and, eventually, the fast collapse of the string towards its final pinned state results in an decay in the activity with $`A_r(t_pt)^\psi `$, where $`t_p`$ is the pinning time in the finite system involved. Driven dynamics of disordered elastic media, which has become a paradigm for a diversity of physical systems, is characterized by the existence of a depinning force $`F_c`$, which separates the pinned state at “sub-threshold” forces $`F<F_c`$ and the sliding regime at “super-threshold” forces $`F>F_c`$. At finite temperatures the system can move at sub-threshold forces as a result of thermally activated jumps between the disorder-induced meta-stable states. The corresponding activation barriers diverge at small forces, $`U(F)\mathrm{}`$ as $`F0`$, giving rise to glassy dynamics with a highly nonlinear response to infinitesimal drives (creep), aging, memory, and hysteretic effects . On the qualitative phenomenological level, the observed memory and hysteretic behaviors can be viewed as a result of slow relaxation of the relevant activation barriers upon altering the drive, yet the detailed quantitative description of transient glassy dynamics remains an appealing task. The static sub-threshold characteristics of the string were studied in and critical depinning and transient behavior in the sliding regime in . In this letter we investigate numerically and analytically the transient sub-threshold zero-temperature dynamics of an elastic string, taken as an exemplary system for a driven elastic disordered medium. We study the time evolution of the string as it slows and eventually stops upon applying a constant driving force in the vicinity of the depinning threshold, $`FF_c`$, where the transient period and the distance traveled by the string are large and diverge as $`FF_c`$. The graphic representation of a typical time evolution of the string velocity $`v(x,t)`$ is given in Fig. 1. It demonstrates the avalanche-like sub-threshold dynamics of the string. The avalanches advance via the propagation of kinks along the string. We quantify the temporal evolution by the activity function $`A(t)`$ defined as the fraction of the moving nodes at time $`t`$, which is proportional to the velocity of the center of mass of the string , and find three distinct regimes of the temporal evolution of the string: (i) an initial regime of a fast decay of activity, according to a power law; (ii) the exponential decay of activity, $`A\mathrm{exp}\left[(t/t_\mathrm{\Delta })^\gamma \right]`$ and (iii) the final avalanche-like immobilization stage, which appears only for strings of finite length for which the activity vanishes and the string gets pinned at a finite time $`t_p`$. Shortly before $`t_p`$ the activity function drops rapidly to zero following the power law, $`A_r(t_pt)^\psi `$, with $`\psi 0.63`$. Model – We consider an elastic string at zero temperature placed on a disordered plane and subject to a constant drive $`F`$ not exceeding the critical depinning force $`F_c`$. The overdamped equation of motion is $`\eta y/t=\delta /\delta y`$, where $`y(x,t)`$ is its lateral displacement at time $`t`$, $`\eta `$ is a friction coefficient. The string has a length $`L`$ and is aligned along the $`x`$ direction. Its energy is $`={\displaystyle _0^L}𝑑x\left[{\displaystyle \frac{C}{2}}\left({\displaystyle \frac{y}{x}}\right)^2+U(y(x),x)\right],`$ with the string elastic constant $`C`$. The string and the media are periodic in $`x`$ with periodicity $`L`$. In the simulations we use a discrete version of the continuum equations. The $`x`$ direction is discretized by a set of lines (rails) of equal spacing, with unit length, constraining the motion of discrete string elements (nodes). Energies are measured in the units of the elastic constant $`C`$. The pinning sites are randomly distributed on the rails according to a specified density $`\rho =0.1`$. The pinning sites are triangular wells with unit width and the depth/strength $`U_p=0.1`$. Forward Euler time marching is used to advance the system in time. The initial state for the string is taken as a straight line. The statistics were collected on $`30`$ samples of length $`L=8192`$. Results – To link to past research, we measured the geometric characteristics of the subthreshold behavior. The roughness of the string at time $`t`$ is $`w_t(x)=[y(x+x^{},t)y(x^{},t)]^2^{1/2}`$, where $`\mathrm{}`$ denotes a positional average over $`x^{}`$ and a disorder average. The roughness increases with $`x`$ (for $`xL`$) until it saturates at $`x\xi _{}(t)`$ with $`w_t(x)\xi _{}(t)`$ (see Fig. 2). Correlation lengths increase with time and reach the final values $`\xi _{}`$ and $`\xi _{}`$ in the pinned state ($`tt_p`$), where the roughness follows a power law $`w_{\mathrm{}}x^\zeta `$ for $`1x\xi _{}`$ . We find the roughening exponent $`\zeta `$ to grow linearly as the applied force approaches its threshold value $`F_c`$ (see inset of Fig. 3). This behavior allows for an alternative determination of $`F_c`$ resting on the analytical estimate $`\zeta =1`$ at the transition . This contrasts with the usual procedure where two parameters $`F_c`$ and $`\zeta `$ have to be simultaneously determined as fitting parameters. The criterion $`\zeta (F_c)=1`$ yields as an estimate for the critical depinning force $`F_c=0.3016`$, which is in excellent agreement with our value of the critical force determined independently from the double-logarithmic plot of roughness ($`F_c=0.3105`$) and also with the earlier finding ($`F_c=0.3058`$). Note a discrepancy with the results in Ref. , where $`\zeta =1.25`$ at the depinning transition has been found in conflict with the analytical result of . The reason for this discrepancy is not yet clear. Next, we studied the behavior of the pinning distance $`D_p`$, which is defined as the ensemble averaged distance that the center of mass of the string traveled before being completely pinned $`D_p=(1/L)_{i=1}^L\|y_i(t_p)y_i(0)\|`$. We verified the scaling behavior of $`D_p(FF_c)^\alpha `$ (cf. Fig. 3). For the given string length $`L=8192`$ the finite-size effects become negligible for the force dependence. We found $`\alpha =2.56\pm 0.08`$ for $`F_c=0.3016`$ in good agreement with the results of for the 1D CDW. The other relevant quantities $`\xi _{}`$ and $`\xi _{}`$, determined from the saturation of the roughness also scale with exponents grouping around $`\alpha 2.4`$. The results are summarized in the Table I, which also gives the values of the critical force determined independently from the divergences of the quantities involved. Having established the geometric characteristics of subthreshold dynamics, we now turn to a detailed study of the temporal evolution of transient behavior. We characterize the temporal evolution of a finite system toward the pinned state by the activity function $`A(t)`$ defined as the fraction of the moving string nodes moving at the given instant $`t`$. For a discrete system evolving in time by steps $`\delta t`$, we define a node as active at time $`t`$ if that node has moved at least one lattice unit in the previous 100 time steps. The activity function is then the ensemble average of the fraction of active nodes. In Fig. 4 we illustrate the temporal behavior of the activity function for the representative subthreshold force $`F=0.29`$. It is important to quantify not only the fraction of active nodes, but also the spatial structure of pinned and moving nodes, which can be quantified by the characteristic length $`L_{\mathrm{pin}}`$ of pinned segments (connected region of pinned nodes). Initially, $`L_{\mathrm{pin}}`$ will be of the order of unity, since the string position is not correlated with the location of the pinning centers and in the final pinned state $`L_{\mathrm{pin}}=L`$ (see inset of Fig. 2). The initial transient regime (i) is expected to last until $`t_\xi `$ given by $`L_{\mathrm{pin}}(t_\xi )=\xi _{}`$ . At $`t>t_\xi `$ relatively large segments of the string are pinned, and the remaining activity is due to the motion of mobile segments over rare regions, where the pinning centers are underrepresented. These mobile segments will advance until the driving force is balanced by other pinning centers or by the string tension. We now construct a lower estimate for the velocity of an infinitely long string as a function of time. We expect the finite string in regime (ii) to follow this dynamics, before finite-size effects set in in regime (iii). We divide the rails in elements of unit length. Then a finite fraction $`p`$ of elements is free of a pinning force. We overestimate the effect of pinning by assuming that the other elements contain an infinitely strong pinning force. Let us now consider a string segment that is pinned on two rails of distance $`l`$, say at $`y(0)=y(l)=0`$. If the rails in between are free of pinning centers, the segment relaxes into the configuration where the driving force is balanced by the string tension. A simple estimate shows that the front of the segment will then reach a distance $`y_lFl^2/C`$ after a time $`t_ly_l\eta /F\eta l^2/C`$. Before this time the ends of length $`x(t)\sqrt{Ct/\eta }`$ of the string segment come to rest, and the center of mass of the segment moves with a velocity $`v_l(t)(F/\eta )[1\sqrt{Ct/\eta l^2}]`$. During this relaxation the string segment moves over an area $`𝒜_lFl^3/C`$. The relaxation takes place in this undisturbed way only if all the rail segments in this area are free of pinning, which occurs with probability $`p_lp^{𝒜_l}`$. To find a lower bound to the string velocity we divide the whole string into segments of arbitrary length $`l`$, which serves as variational parameter. Although the string can be taken as unpinned initially, the velocity is underestimated by formally considering the links between these sections as pinned. A fraction $`p_l`$ of such segments relaxes without encountering pinning rail elements and contributes $`p_lv_l(t)`$ to the average string velocity. This contribution is a lower bound to the velocity for all possible $`l>\sqrt{Ct/\eta }`$, and the best estimate is obtained by $`v_{\mathrm{min}}(t)\mathrm{max}_lp_lv_l(t)p^{(F/C)(Ct/\eta )^{3/2}}`$, which is finite at all times. Consequently, the activity function is also expected to decay not faster than $`A(t)\mathrm{exp}(\mathrm{const}.t^{3/2})`$. In the finite system this exponential decay stops when the characteristic distance between the moving segments $`L_{\mathrm{pin}}l/p_l`$ becomes of the order of the system size and the string experiences a deficiency of the moving nodes. Our conjecture is that the origin of this “superfast” stopping regime is the avalanche-like disappearance of the moving nodes of the string. Once pinned, a node cannot depin any more; moreover, it causes the immediate pinning of the neighboring nodes, in a close analogy with the avalanche clustering at the final stage of the coagulation process . This analogy suggests the power-law termination of the motion according to $`A_r(t_pt)\|t_pt\|^\psi `$, where $`t_p`$ is the time of the total immobilization. Although $`t_p`$ is subject to strong sample-to-sample fluctuations because it is determined by fluctuations on the scale of the system size, the final regime can be recognized only by averaging the activity considered as a function of the reverse time $`t_pt`$; that is, the final regime can be seen only if the averaging over $`N`$ samples is performed according to $`A_r(t_pt)=(1/N)_{i=1}^NA_i(t_p^{(i)}t)`$ in contrast to the usual average $`A(t)=(1/N)_{i=1}^NA_i(t)`$. We give lower and upper limits on the exponent $`\psi `$. Consider the process in which the node that came to rest enhances the probability that the nearest node will stop at the next moment of time. In this case the boundary between localized and mobile nodes will move some average velocity, and the activity will decrease as $`A_rt_pt`$, which means that the upper limit on $`\psi `$ is $`\psi _{\mathrm{upper}}=1`$. This type of behavior can be expected when $`F`$ is far below $`F_c`$. Another limiting type of terminal dynamics occurs if immobile nodes do not affect dynamics in mobile regions. In this case the boundary between the immobile and mobile regions moves randomly, and the finite string has a nonzero chance of stopping when the boundaries enclosing the moving region collide with each other. Considering this process in reverse time, we conclude that $`A_r(t_pt)^{1/2}`$, and we get a lower limit $`\psi _{\mathrm{lower}}=1/2`$. Note, that the possibility of fragmentation of the moving segment of the string into several subsegments does not affect this estimate. In the more general case of a moving manifold of internal dimension $`d`$ ($`d=1`$ for the string), similar arguments result in $`d/2<\psi <d`$. Numerical simulations clearly show distinct relaxational regions. After the initial relaxation (regime (i)), the activity crosses over to the exponential decay regime, $`A(t)\mathrm{exp}\left[(t/t_\mathrm{\Delta })^\gamma \right]`$, $`\gamma 1`$. The precision of our data is not sufficient to determine a more precise estimate of $`\gamma `$; from the above consideration we could expect $`1<\gamma 3/2`$. The relaxation time $`t_\mathrm{\Delta }`$ diverges as the applied force approaches the threshold value (see Fig. 3). However, although one could expect that $`t_\mathrm{\Delta }=t_\xi `$ and therefore that $`t_\mathrm{\Delta }\|F_cF\|^{\nu z}`$ (since $`\xi \|FF_c\|^\nu `$), with $`z\nu =4/3`$ using the exponents of the depinning transition , we have found a different exponent: $`t_\mathrm{\Delta }\|F_cF\|^\mathrm{\Delta }`$ with $`\mathrm{\Delta }=2.74`$. A careful examination of a final stage of the relaxation process reveals a very sharp drop of the activity function just before the complete immobilization of the string. To analyze this stage, we introduce the reverse time representation: we present the activity function as a function of $`t_r=t_pt`$, counting time backwards from the moment of achieving final immobilization. Using this representation enables us to carry out an ensemble averaging despite the fact that different samples have different pinning times. The plot of $`\mathrm{ln}A_r`$ vs. $`\mathrm{ln}(t_pt)`$ reveals a power-law final behavior (see Fig. 5): $`A_r(t_pt)^\psi `$, with $`\psi =0.63\pm 0.06`$, which is within the above upper and lower estimates. We note that the reverse time plot of $`\mathrm{ln}A_r`$ reconfirms the existence of the intermediate regime with the essentially same characteristic time $`t_\mathrm{\Delta }`$ as was found from the direct time picture. In summary, we investigated subthreshold dynamics of a pinned elastic string and identified three distinct transient regimes: (i) a fast initial relaxation, (ii) an intermediate exponential decay of the activity resulting from residual motion in the exponentially rare regions free of defects and (iii) a novel avalanche-like terminal relaxation to the pinned state, $`A_r(t_pt)^\psi `$, resulting from the finite-sized effects. This final stage exhibits striking similarity to coagulation dynamics. This work was supported by Argonne National Laboratory through the U.S. Department of Energy, BES-Material Sciences and BES-MICS, under contract No. W-31-109-ENG-38, and by the NSF-Office of Science and Technology Centers under contract No. DMR91-20000 Science and Technology Center for Superconductivity.
warning/0002/cond-mat0002117.html	ar5iv	text	# Role of Secondary Motifs in Fast Folding Polymers: A Dynamical Variational Principle \[ ## Abstract A fascinating and open question challenging biochemistry, physics and even geometry is the presence of highly regular motifs such as $`\alpha `$-helices in the folded state of biopolymers and proteins. Stimulating explanations ranging from chemical propensity to simple geometrical reasoning have been invoked to rationalize the existence of such secondary structures. We formulate a dynamical variational principle for selection in conformation space based on the requirement that the backbone of the native state of biologically viable polymers be rapidly accessible from the denatured state. The variational principle is shown to result in the emergence of helical order in compact structures. \] A fundamental problem in every day life is that of packing with examples ranging from fruits in a grocery, clothes and personal belongings in a suitcase, atoms and colloidal particles in crystals and glasses, and amino acids in the folded state of proteins. The simplest problem in packing consists of determining the spatial arrangement that accomodates the highest packing density of its constituent entities with the result being a crystalline structure. Besides packing considerations, dynamical effects play a significant role when rapid packing/unpacking is entailed, as in the formation of amorphous glasses where crystallization is dynamically thwarted or in the more familiar suitcase problem. Fast packing has been recognized as a central issue for biopolymers, such as proteins, since the early work of Levinthal . Further, the native conformations display extremely regular motifs, such as $`\alpha `$-helices or $`\beta `$-sheets. In this Letter we postulate a direct connection between the dynamics of rapid folding and the emergence of secondary motifs in the native state conformations. In fact, an intuitive approach to rapid and reproducible folding might be to create neat patterns of lower dimensional manifolds than the physical space and bend and curl them into the final folded state. For proteins, secondary structures such as $`\alpha `$-helices and $`\beta `$-sheets are indeed patterns in low dimensions. There are two key aspects distinguishing a protein from a generic heteropolymer: the specially selected sequence of amino acids and the three-dimensional structure that it folds reversibly into. For a given target native structure, the selection mechanism in sequence space is the principle of minimal frustration . The chosen sequences are such that their target native states are reached through a funnel-like landscape which facilitates the harmonious fitting together of pieces to form the whole. The three-dimensional structure impacts on the functionality of the protein and a fascinating issue is the elucidation of the selection mechanism in conformation space that picks out certain viable structures from the innumerable ones with a given compactness. Earlier studies have shown that there is a direct link between viable native conformations and high designability . Recently , it was observed that the natural folds of proteins have a much larger density of nearby structures than generic (artificial) conformations of the same character and that the exceedingly large geometrical accessibility of natural proteins may be related to the presence of secondary motifs. The realization that proteins have secondary structures arose with early crystallographic studies and the brilliant deduction of Pauling et al. of the ability of an $`\alpha `$-helix of the correct pitch to accomodate hydrogen bonds, thus promoting its stability. Inspired by the findings of Pauling, helix-coil transition models have been used to study the thermodynamics of helix formation . The models encompass features that ensure the helical nature of the low-energy states by assuming first that that monomers can be in a helical state and by then introducing co-operative interactions that favor helical regions. It is interesting to note, however, that the number of hydrogen bonds is nearly the same when a sequence is in an unfolded structure in the presence of a polar solvent or in its native state rich in secondary structure content . It has also been suggested that the $`\alpha `$-helix is an energetically favorable conformation for main-chain atoms but the side-chain suffers from a loss of entropy . Nelson et al. have shown both numerically and experimentally that non-biological oligomers fold reversibly like proteins into a specific three-dimensional structure with high helical content driven only by solvophobic interactions. Recent studies have attempted to explain the emergence of secondary structure from geometrical principles rather than invoking detailed chemistry. Despite the concerted efforts of several groups, a simple general explanation remains elusive. In particular, the work of Yee et al. , Hunt et al. , and Socci et al. have shown that compactness alone can only account for a small secondary structure content. These facts are also corroborated by the recent study of the kinetics of homopolymer collapse, where no evidence was found for the formation of local regular structures . We propose a selection mechanism in structure space in the form of a variational principle postulating that, among all possible native conformations, a protein backbone will attain only those which are optimal under the action of evolutionary pressure favouring rapid folding. Our goal is to elucidate the role played by the bare native backbone independent of the selection in sequence space and hence of the (imperfectly-known) inter-amino-acid potentials. We therefore choose to employ a Go-like model with no other interaction that promotes or disfavours secondary structures. The model is a sequence-independent limiting case of minimal frustration which, for a given target native state conformation, favours the formation of native contacts – the energy of a sequence in a conformation is simply obtained as the negative of the number of contacts in common with the target conformation. We will consider two non-consecutive amino acids to be in contact if their separation is below a cutoff $`r_0=6.5`$ Å(the results are qualitatively similar when slightly different values of $`r_0`$ in the range $`68`$ Åare chosen). The energy of structure $`\mathrm{\Gamma }`$ in the Go model is given by $$H(\mathrm{\Gamma })=\frac{1}{2}\underset{i,j}{}\mathrm{\Delta }_{i,j}(\mathrm{\Gamma })\mathrm{\Delta }_{i,j}(\mathrm{\Gamma }_0)$$ (1) where the sum is taken over all pairs of amino acids, $`\mathrm{\Gamma }_0`$ is the target structure, $`\mathrm{\Delta }_{i,j}(\mathrm{\Gamma })`$ is the contact map of structure $`\mathrm{\Gamma }`$: $$\mathrm{\Delta }_{ij}(\mathrm{\Gamma })=\{\begin{array}{cc}1\hfill & R_{ij}<r_0\text{ and }\|ij\|>2;\hfill \\ 0\hfill & \text{ otherwise, }\hfill \end{array}$$ (2) where $`R_{ij}`$ is the distance of amino acids $`i`$ and $`j`$. The polypeptide chain is modelled as a chain of beads subject to steric constraints . We adopted a discrete representation similar to the one of Covell and Jernigan , in which each bead occupies a site of an FCC lattice with lattice spacing equal to 3.8 Å. Such a representation is able to describe the backbone of natural proteins to better that 1 Å rmsd per residue (equal to the best experimental resolution) and preserves typical torsional angles. All discretized structures were subject to a suitable constraint: any two non-consecutive residues cannot be closer than $`4.65`$ Å due to excluded volume effects and the distance between consecutive residues can fluctuate between 2.6 Å$`<d<`$ 4.7 Å. Such constraints were determined by an analysis of the coarse-grainings of several proteins of intermediate length ($`100`$ residues). In order to enforce a realistic global compactness for a backbone of length $`L`$, the number of contacts in all the target structures considered was chosen to be around $`N=1.9L`$ while, locally, no residue was allowed to make contact with four or more consecutive residues. In order to assess the validity of the variational principle, it is necessary to evaluate the typical time, $`t(\mathrm{\Gamma }_0)`$, taken to fold into a given target structure, $`\mathrm{\Gamma }_0`$, followed by a selection of the structures $`\mathrm{\Gamma }_0`$, that have the smallest folding times. To do this, an initial set of ten conformations was generated by collapsing a loose chain starting from random initial conditions. In each case, we modified the random initial conformation by using Monte Carlo dynamics: we move up to 3 consecutive beads to unoccupied discrete positions that do not violate any of the physical constraints and accept the moves according to the standard Metropolis rule. The energy is given by eq. (1), while the temperature for the MC dynamics was set to 0.35. This value was chosen in preliminary runs so that it was higher than the temperature below which the sequence is trapped in metastable states but comparable to the folding transition temperature so that conformations with significant overlap with the native state are sampled in thermal equilibrium. For each structure, as a measure of the folding time we took the median over various attempts (typically 41) of the total number of Monte Carlo moves necessary to form a pre-assigned fraction of native contacts, typically 66%, starting from a random conformation. Our results were unaltered on increasing this fraction to 75%; indeed, this fraction could be progressively increased towards 100% with successive generations without increase in the computational cost since better and better folders are obtained. A new generation of ten structures is created by “hybridizing” pairs of structures of the previous generation ensuring that structures with small folding times are hybridized more and more frequently as the number of generations, $`g`$, increases . To do this, each of the two distinct parent structures to be hybridized, $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ are chosen with probability proportional to $`\mathrm{exp}[(g1)ft)/1000]`$, where $`g`$ is the index of the current generation (initially equal to 1), $`ft`$ is the median folding time. Then, a hybrid map is created by taking the union of the two parent maps: $$\mathrm{\Delta }_{ij}^{Union}=\mathrm{max}(\mathrm{\Delta }_{ij}(\mathrm{\Gamma }_1),\mathrm{\Delta }_{ij}(\mathrm{\Gamma }_2)).$$ (3) Because it is not guaranteed that $`\mathrm{\Delta }^{Union}`$ corresponds to a three-dimensional structure obeying the same physical constraints as $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$, the corresponding hybrid $`\mathrm{\Gamma }`$ is constructed by taking one of the two parent structures (or alternatively a random one) as the starting conformation and carrying out MC dynamics favouring the formation of each of the contacts in the union map (i.e. using eq. (1) with $`\mathrm{\Delta }_{ij}(\mathrm{\Gamma }_0)`$ substituted by $`\mathrm{\Delta }_{ij}^{Union}`$). The dynamics is carried out starting from a temperature of $`0.7`$ and then decreasing it gradually over a sufficiently long time (typically thousands of MC steps) to achieve the maximum possible overlap with the union map, while simultaneously maintaining the realistic compactness. The resulting hybrid structure is typically midway between the two parent structures, in that it inherits native contacts from both of them. We adopted the following definition in order to obtain an objective and unbiased way to quantitatively estimate the presence of secondary content: a given residue, $`i`$ was defined to belong to a secondary motif if, for some $`j`$, one of these conditions held: $`a)\mathrm{\Delta }_{i1,j1}`$ $`=`$ $`\mathrm{\Delta }_{i,j}=\mathrm{\Delta }_{i+1,j+1}=\mathrm{\Delta }_{i,j+1}`$ (4) $`=`$ $`\mathrm{\Delta }_{i+1,j+2}=\mathrm{\Delta }_{i1,j}=1;`$ (5) $`b)\mathrm{\Delta }_{i+1,j1}`$ $`=`$ $`\mathrm{\Delta }_{i,j}=\mathrm{\Delta }_{i1,j+1}=\mathrm{\Delta }_{i,j+1}`$ (6) $`=`$ $`\mathrm{\Delta }_{i+1,j}=\mathrm{\Delta }_{i1,j+2}=1.`$ (7) The former \[latter\] identifies the presence of helices and parallel \[anti-parallel\] $`\beta `$ sheets in natural proteins, which can be identified by the visual inspection of contact matrices and appears as thick bands parallel or orthogonal to the diagonal. The upper plot of Fig. 1 shows the decrease of the typical folding time over the generations for chains of length 25, while the middle panel shows the accompanying increase in the number of residues in secondary motifs (secondary content). The bottom panel shows a milder decrease of the contact order (i.e. a larger number of short-range contacts) as the generations evolved, in agreement with the experimental findings of Plaxco et al. One of the optimal structures of length 25 is shown in Figure 2a. Due to the absence of any chirality bias in our structure space exploration, the helix does not have a constant handedness. The signature of the secondary motifs in the optimal structures is clearly visible in the contact maps of Figure 3, which are not sensitive to structure chirality. Strikingly, the variational principle selects conformations with significant secondary content as those facilitating the fastest folding. The correlation of the emergence of secondary structures with decrease of folding times is shown in the plot of Fig. 4. We verified that the hybridization procedure is not biased towards low contact order by iterating it for various generations and hybridizing the structures at random. Even after dozens of generations, the generated structures had secondary contents of about 1/3-1/4 of the true extremal structures. The very high secondary content in optimal conformations was found to be robust against changes in chain length or compactness of the target structure. On requiring that the structure be more compact, bundles of helices emerge \[see Fig. 2b\] along with an increase in contact order, signalling the presence of some longer range contacts, which are necessitated in order to accomodate the shorter radius of gyration. It is noteworthy that our calculations lead predominantly to $`\alpha `$-helices and not $`\beta `$ sheets, a fact accounted for by the demonstration that steric overlaps and the associated loss of entropy lead to the destabilization of helices in favor of sheets , the appearance of such sheets only in sufficiently long proteins and the much slower folding rate of $`\beta `$-sheets compared to $`\alpha `$-helices . It is remarkable that the same requirement of rapid folding is sufficient to lead to a selection in both sequence and structure space underscoring the harmony in the evolutionary design of proteins. The results and strategies presented here ought to be applicable in protein-engineering contexts, for example by ensuring optimal dynamical accessibility of the backbone of proteins. A systematic collection of the rapidly-accessible structures of various length should also lead to the creation of unbiased libraries of protein folds. Acknowledgements This work was supported by INFM, INFN sez. di Trieste, NASA, NATO and The Donors of the Petroleum Research Fund administered by the American Chemical Society. We thank F. Seno and A. Trovato for useful discussions.

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